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# References
IUHET-423
Supersymmetric left-right model and scalar potential
Biswajoy Brahmachari
Physics Department, Indiana University, Bloomington IN-47405, USA
Abstract
We construct a scalar potential of supersymmetric left-right model in the limit when supersymmetry is valid.
Left-right symmetric model is a natural extension of standard model. The symmetry breaking chain can be written as
$`G_{LR}[SU(2)_L\times SU(2)_R\times U(1)_{BL}]`$ $`{\displaystyle \genfrac{}{}{0pt}{}{\mathrm{\Delta }_R}{}}G[SU(2)_L\times U(1)_Y]`$ (1)
$`{\displaystyle \genfrac{}{}{0pt}{}{H_1H_2}{}}G_0[U(1)_{em}].`$
Symmetry breaking mechanism is normally understood in the following way. One uses symmetries of $`G_{LR}`$ to write down allowed terms of a classical scalar potential. A stable breaking of symmetry is obtained when field configuration (values of fields) is such that the potential is at its minimum at all space-time points. We know that these values of fields are the VEVs and we also know that the VEVs obey residual symmetries. Hence electric charge and color symmetries remain.
Mininum energy state of Higgs scalars need not have the symmetry properties of the gauge bosons. So the symmetry of the combined system of fermions, gauge bosons and Higgs scalars can be broken. In this paper we will first state a miminal Higgs scalar spectrum of left-right symmetric model. Then we will write them in the form of square matrices and apply Hamiltion-Cayley theorem to these matrices to see what happens. It will lead to polymonial equations (quadratic for this case) which are satisfied by these matrices. We will physically interpret these equations as minimization conditions.
Minimal Higgs choice for a model based on $`G_{LR}`$ is
$$\mathrm{\Delta }_R=(1,3,2)\mathrm{\Delta }_L=(3,1,2)\varphi =(2,2,0).$$
One can write these VEVs as matrices in a $`SU(2)_L\times SU(2)_R`$ basis where rows are $`SU(2)_R`$ multiplets and columns are $`SU(2)_L`$ multiplets. In writing so one can suppress abelian $`U(1)_{BL}`$ as well as non-abelian color degrees of freedoms. Then one gets
$$\varphi =\left(\begin{array}{cc}\kappa _1& 0\\ 0& \kappa _2\end{array}\right)\mathrm{\Delta }_L=\left(\begin{array}{cc}0& 0\\ v_L& 0\end{array}\right)\mathrm{\Delta }_R=\left(\begin{array}{cc}0& v_R\\ 0& 0\end{array}\right).$$
(2)
We can read-off that $`\varphi `$ breaks both $`SU(2)_L`$ and $`SU(2)_R`$ symmetries. However, $`\mathrm{\Delta }_R`$ breaks only the $`SU(2)_R`$ symmetry. From experiments one knows that $`SU(2)_R`$ symmetry is broken at a higher scale. This is because gauge bosons corresponding to broken $`SU(2)_R`$ symmetries are yet to be observed. Even though mass splitting between fermions within a $`SU(2)_L`$ multiplet parametrize $`SU(2)_L`$ breaking, values of parameters $`\kappa _1,\kappa _2`$ remain unknown up to magnitudes of Yukawa couplings of corresponding fermions. Similarly we can think of similar manifestations of $`SU(2)_R`$ and $`SU(2)_L\times SU(2)_R`$ symmetries in the masses of extra fermions. Thus we must find a way to study allowed values of $`\kappa _1,\kappa _2,v_L`$ and $`v_R`$ from theory. This is a motivation to further study the scalar potential.
A number of studies of minimizing scalar potential exist. Typically, one writes the most general potential using gauge symmetries and then it is minimized by taking derivatives of potential with respect to VEVs and equating them to zero. This set of equations constrain parameter space of VEVs. We consider a reverse situation. We ask whether given only matrix forms of VEVs is it possible to construct a minimal potential which leads to desired symmetry breaking for all possible values of $`\kappa _1,\kappa _2,v_L`$ and $`v_R`$ ? In other words, instead of writing the complete potential and study the ranges of $`\kappa _1,\kappa _2,v_L`$ and $`v_R`$, is it possible to obtain the unique subset of terms of the potential which allow all possible values of $`\kappa _1,\kappa _2,v_L`$ and $`v_R`$ ? This what we answer below.
The Hamilton-Cayley theorem states: Every square matrix must satisfy its own characteristic equation. That is, if
$$\mathrm{det}(𝐀\lambda 𝐈)=c_n\lambda ^n+c_{n1}\lambda ^{n1}+\mathrm{}+c_2\lambda ^2+c_1\lambda +c_0$$
(3)
then
$$c_n𝐀^n+c_{n1}𝐀^{n1}+\mathrm{}+c_2𝐀^2+c_1𝐀+c_0=0.$$
(4)
Using equations (2) (3) and (4) we get
$$\varphi ^2(\kappa _1+\kappa _2)\varphi +\kappa _1\kappa _2=0\mathrm{\Delta }_L^2=0\mathrm{\Delta }_R^2=0.$$
(5)
These equations are satisfied for any value of $`\kappa _1,\kappa _2,v_L`$ and $`v_R`$.
Suppersymmtry allows only trilinear terms in the superpotential. Thus minimization conditions are at most quadratic. It is a nice coincidence that $`2\times 2`$ matrices lead to quadratic characteristic equations. The term $`\kappa _1\kappa _2`$ is present in the minimization condition. We must have a linear term in superpotential. However there is no singlet scalar in the model. So we must have
$$\mathrm{Either}:\kappa _10\mathrm{or}:\kappa _20$$
(6)
Then the superpotential reads as,
$`W=W_1+W_2`$
$`W_1={\displaystyle \frac{1}{3}}\varphi ^3{\displaystyle \frac{1}{2}}\kappa _1\varphi ^2`$
$`W_2=M_L\mathrm{\Delta }_L^2+M_R\mathrm{\Delta }_R^2`$ (7)
$`W_2`$ vanishes identically, for all values of $`v_L,v_R,M_L,M_R`$. It does not contribute to energy if the VEVs are of the form of (2). $`W_1`$ however needs to be minimized and all possible values of $`\kappa _1`$ or $`\kappa _2`$ are not allowed. We had to chose either $`\kappa _1`$ or $`\kappa _2`$ to vanish. So we have got a negative answer to our question. This is our result.
Thus we have constructed a superpotential of supersymmetric left-right model using Hamilton-Cayley theorem. We had to chose either $`\kappa _1`$ or $`\kappa _2`$ to vanish. We have chosen $`\kappa _2=0`$. This means that down sector of fermions remains massless.
This research was supported by U.S Department of Energy under the grant number DE-FG02-91ER40661.
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# Protostellar Collapse with Various Metallicities
## 1 Introduction
In the standard theory of cosmic structure formation, galaxies are believed to have originated from primeval density fluctuations. Many authors have been studying the galaxy formation process in this context and have made remarkable progress in understanding such processes as the gravitational growth of the density fluctuation, its decoupling from the cosmic expansion and virialization, and the subsequent radiative cooling of the baryonic gas (e.g., Padmanabhan 1993). However, processes after the cooling of baryonic gas, namely, the transformation of a pregalactic cloud into a cluster of stars, are relatively poorly known. It is inevitably necessary to investigate them in studying the galaxy formation process since the very existence of stellar components is one distinct feature of galaxies. These processes can be viewed as successive fragmentations of a cloud into fragments, or in other words, protostellar clouds, and their contraction into stars (e.g., Hayashi 1984).
Whether a cloud fragments or collapses into a single object depends on the the central flatness, that is, the axial ratio of the isodensity contour in the central region (Tsuribe & Inutsuka 1999). In particular, warm clouds do not fragment before adiabatic core formation, whereas clumps sufficiently more massive than the Jeans mass first collapse disklike and next fragment most likely into filamentary clouds (Miyama, Narita, & Hayashi 1987a, 1987b). The filamentary clouds fragment again into protostellar cloud cores after further contraction. Uehara et al. (1997) investigated the gravitational collapse of metal-free filamentary clouds using one-zone approximation and found that the minimum mass of fragments, or protostellar cores, is essentially Chandrasekhar mass, i.e., $`1M_{\mathrm{}}`$. Nakamura & Umemura (1999) confirmed Uehara et al. (1997)’s result by performing one-dimensional hydrodynamical calculations. In this paper, we will discuss the collapse of protostellar cores into stars.
The distinction between the star formation in the galaxy formation epoch and that in present-day star-forming regions mainly resides in the difference of the composition of gas that stars are made from. Metallicity, and therefore the gas to dust ratio, is lower for earlier star formation. In present-day star forming regions, protostellar clouds remain roughly 10K in a wide range of densities owing to efficient dust emission (e.g., Hayashi & Nakano 1965) in the course of contraction.
On the other hand, many authors have studied the thermal and chemical evolution of collapsing primordial clouds with one-zone approximation (e.g., Matsuda, Sato, & Takeda 1969; Yoneyama 1972; Calberg 1981; Palla, Salpeter, & Stahler 1983; Izotov & Kolesnik 1984; Lahav 1986; Puy & Signore 1997) or by hydrodynamical calculations (Matsuda et al. 1969; Villere & Bodenheimer 1987; Haiman, Thoul, & Loeb 1996; Omukai & Nishi 1998; Oliveira et al. 1998). The collapse of metal-free protostellar clouds, which is relevant for the first star formation in the universe, is induced by radiative cooling owing to molecular hydrogen lines, and their temperatures are about 1000K.
Above are two extremes in composition of protostellar gas, namely, the former for the high-metallicity end and the latter for the low-metallicity end.
Then, how does the collapse of slightly metal-polluted clouds proceed? Even in the galaxy formation epoch, significant metal enrichment can occur since the lifetime of massive stars ($`10^6`$ yr) is shorter than the free-fall time of pregalactic clouds ($`10^8`$ yr). In fact, both Ly$`\alpha `$ forest clouds and Population II stars are slightly polluted by heavy elements. Therefore, it is important to study star formation from slightly metal-polluted gas. In this paper, we aim to fill the gap between the protostellar collapse of primordial and present-day interstellar gas.
This topic has been investigated in the pioneering works of Low & Lynden-Bell (1976) and Yoshii & Sabano (1980). However, the following points remain to be fulfilled; in their works, (i) chemical reactions had not been solved, (ii) molecular coolants, which become important at high densities, had not been considered, and (iii) the evolution of clouds as dense as the stellar density ($`>10^{22}\mathrm{cm}^3`$) had not been studied. With these points in mind, we investigate in this paper the thermal and chemical evolution of collapsing spherical clouds as a function of metallicity.
The outline of this paper is as follows. In §2, we describe the method of our calculations. In §3, results of our calculations are presented. We summarize our work in §4.
## 2 Model
### 2.1 Basic Equations
We consider a spherical cloud with mean metallicity $`Z`$. The helium concentration is assumed to be $`y_{\mathrm{He}}=0.0972`$ for all clouds.<sup>1</sup><sup>1</sup>1The concentration of element X is defined by
$$y_X=n_X/n,$$
(1) where $`n`$ and $`n_X`$ are the number densities of hydrogen nuclei and nuclei of element X. Similarly, we write for each atomic, molecular, or ionic species
$$y(x)=n(x)/n,$$
(2) where $`n(x)`$ is the number density of species $`x`$. Note $`y(\mathrm{H}_2)=1/2`$ for fully molecular gas. We assume that the dust-to-gas ratio is proportional to the mean metallicity, or in other words, that a fixed fraction of the heavy elements in the interstellar medium condenses into dust grains. We adopt the Pollack et al. (1994) model of grains in molecular clouds (see §2.2.2). For the case of local interstellar clouds ($`Z1Z_{\mathrm{}}`$), the mass fraction of grain we used is $`0.934\times 10^2`$ in the lowest temperature regime, and the gas-phase elemental abundances are $`y_\mathrm{C}=0.927\times 10^4`$, and $`y_\mathrm{O}=3.568\times 10^4`$, respectively, which correspond to 46% of oxygen and 72% of carbon depleted onto grains. These values are reduced proportionally for lower metallicity cases. We normalize the mean metallicity (and therefore also the dust-to-gas ratio and gas-phase metallicity in our model) relative to the local interstellar values $`Z_{\mathrm{local}}`$ and denote relative metallicity by $`zZ/Z_{\mathrm{local}}`$.
We calculate the time evolution of the central density, temperature, and chemical composition of the collapsing core.
First, we assume that the dynamics is described by the free-fall relation,
$$\frac{d\rho }{dt}=\frac{\rho }{t_{\mathrm{ff}}},$$
(3)
where $`\rho `$ is the density in the central region and the free-fall time is
$$t_{\mathrm{ff}}\sqrt{\frac{3\pi }{32G\rho }}.$$
(4)
The thermal evolution is followed by solving the energy equation
$$\frac{de}{dt}=p\frac{d}{dt}(\frac{1}{\rho })^{(\mathrm{net})},$$
(5)
where
$$e=\frac{1}{\gamma _{\mathrm{ad}}1}\frac{kT}{\mu m_\mathrm{H}}$$
(6)
is the specific internal energy,
$$p=\frac{\rho kT}{\mu m_\mathrm{H}}$$
(7)
is the pressure for an ideal gas, $`\gamma _{\mathrm{ad}}`$ is the adiabatic exponent, $`T`$ is the temperature, $`\mu `$ is the mean molecular weight, $`m_\mathrm{H}`$ is the mass of hydrogen nucleus, and $`^{(\mathrm{net})}`$ is the net energy loss rate per unit mass (see §2.2).
We neglect any effect owing to rotation or magnetic fields for simplicity. In this case, the actual collapse is expected to proceed like the Penston-Larson similarity solution (Penston 1969; Larson 1969). According to this solution, the cloud consists of two parts, that is, the central core region, which has flat density distribution, and the envelope, where the density decreases outward as $`r^2`$. The size of the central flat region is roughly given by the local Jeans length $`\lambda _\mathrm{J}=\pi c_\mathrm{s}/\sqrt{G\rho }`$ in the core. Since we focus on the evolution of the central region, optical depth is estimated by that across one local Jeans length:
$$\tau _\nu =\kappa _\nu \rho \lambda _\mathrm{J}.$$
(8)
This formulation of Jeans length shielding is the same procedure as that of Low & Lynden-Bell (1976) and Silk (1977), although they introduced it as a result of successive fragmentation into Jeans mass clouds, not because of the Penston-Larson collapse.
### 2.2 Cooling/Heating processes
In addition to compressional heating, we include cooling owing to (i) atomic and molecular line radiation, (ii) energy transfer between the gas and the dust grains, which will be emitted as infrared radiation from the grains, (iii) continuous radiation from the gas, and (iv) cooling and heating associated with chemical reactions. Then
$$^{(\mathrm{net})}=_{\mathrm{line}}+_{\mathrm{gr}}+_{\mathrm{cont}}+_{\mathrm{chem}},$$
(9)
where the terms in the right-hand side of the equation correspond to the cooling rates owing to the processes (i)-(iv) above in the same order. Although we treat clouds with metallicity up to $`1Z_{\mathrm{}}`$, we are concerned mainly with the lower metallicity environments that are relevant to the epoch of galaxy formation. In this paper, we do not try to reproduce the present Galactic environment and neglect any external radiation (e.g., UV radiation, cosmic rays, cosmic background radiation, etc.) for simplicity. Then, the considered heating sources are the compressional work and H<sub>2</sub> formation. As will be mentioned below, lower metallicity clouds tend to be warmer. Therefore, the compressional heating, which is proportional to the temperature of the cloud, becomes more important if the external heating rate is constant.
#### 2.2.1 Line Cooling
As line cooling agents, we include atomic fine-structure transitions of CI, CII, and OI and molecular rovibrational transitions of H<sub>2</sub>, CO, OH, and H<sub>2</sub>O:
$$_{\mathrm{line}}=_{\mathrm{CI}}+_{\mathrm{CII}}+_{\mathrm{OI}}+_{\mathrm{H}_2}+_{\mathrm{CO}}+_{\mathrm{OH}}+_{\mathrm{H}_2\mathrm{O}}.$$
(10)
The line cooling rate of species $`x`$ is given by
$$_x=\frac{1}{\rho }\underset{(ij)}{}n(x,i)A_{ij}ϵ_{ij}h\nu _{ij},$$
(11)
where $`n(x,i)`$ is the population density of species $`x`$ in level $`i`$, $`A_{ij}`$ is the spontaneous transition probability, $`ϵ_{ij}`$ is the escape probability, and $`h\nu _{ij}`$ is the energy difference between levels $`i`$ and $`j`$. The population density $`n(x,i)`$ is obtained from a solution of the equation of statistical equilibrium
$$n(x,i)\underset{ji}{\overset{n}{}}R_{ij}=\underset{ji}{\overset{n}{}}n(x,j)R_{ji},$$
(12)
where $`n`$ is the total number of lines included and
$$R_{ij}=\{\begin{array}{cc}A_{ij}ϵ_{ij}+C_{ij}\hfill & \text{for }i>j\hfill \\ C_{ij}\hfill & \text{for }i<j,\hfill \end{array}$$
(13)
ignoring the external radiation. Here $`C_{ij}`$ is the collisional rate from level $`i`$ to level $`j`$. Following Takahashi, Hollenbach, & Silk (1983), we use the escape probability for the case that the velocity is proportional to the radius in the central region:
$$ϵ_{ij}=\left(\frac{1\mathrm{e}^{\tau _{ij}}}{\tau _{ij}}\right)\mathrm{e}^{\tau _{\mathrm{cont}}}.$$
(14)
The optical depth averaged over the line and the continuum optical depth are given by
$$\tau _{ij}=\frac{A_{ij}c^3}{8\pi \nu _{ij}^3\eta _\mathrm{T}}[n(x,j)g_i/g_jn(x,i)]l_{\mathrm{sh}}/(2\mathrm{\Delta }v_\mathrm{D}),$$
(15)
and
$$\tau _{\mathrm{cont}}=(\kappa _{\mathrm{gr}}+\kappa _{\mathrm{gas}})\rho \lambda _\mathrm{J},$$
(16)
respectively, where $`\eta _\mathrm{T}`$ is the multiplicity factor, $`g_i`$ is the statistical weight of level $`i`$, $`\mathrm{\Delta }v_\mathrm{D}`$ is the velocity dispersion, and $`l_{\mathrm{sh}}`$ is the shielding length (Takahashi et al. 1983):
$$l_{\mathrm{sh}}=\mathrm{min}(\mathrm{\Delta }s_{\mathrm{th}},\lambda _\mathrm{J});\mathrm{\Delta }s_{\mathrm{th}}=2\mathrm{\Delta }v_\mathrm{D}/(\frac{dv}{dr})=6\mathrm{\Delta }v_\mathrm{D}t_{\mathrm{ff}},$$
(17)
where we have used the relation for the homogeneous collapse $`v=r/3t_{\mathrm{ff}}`$ and, in the case of the small velocity gradient, have assumed Jeans length shielding as introduced in §2. As a source of the velocity dispersion, we consider only the thermal motion of atoms and neglect microturbulent motions. Then
$$\mathrm{\Delta }v_\mathrm{D}=\sqrt{\frac{2kT}{\mu _xm_\mathrm{H}}},$$
(18)
where $`\mu _x`$ is the molecular weight of species $`x`$. The continuum Planck mean opacity of dust $`\kappa _{\mathrm{gr}}`$ and that of gas $`\kappa _{\mathrm{gas}}`$ are adopted from Pollack et al.(1994) and Lenzuni, Chernoff, & Salpeter (1991), respectively (see §2.2.2 and 2.2.3).
In general, the level population and the escape probabilities depend on each other. We then need an iterative procedure to find their consistent solution except for the case of H<sub>2</sub>. For H<sub>2</sub> lines, we need not iterate to find the consistent set of the level populations and the escape probabilities, since in our calculations, clouds become opaque to H<sub>2</sub> lines only at the density exceeding the critical density, and in such a case, molecules populate following a Boltzmann distribution regardless of the escape probabilities.
The related parameters of transitions are given in Hollenbach & McKee (1989) for CI, CII, and OI, in Hollenbach & McKee (1979;1989) for OH and H<sub>2</sub>O <sup>2</sup><sup>2</sup>2The OH and H<sub>2</sub>O collision cross sections and H<sub>2</sub>O multiplicity factor are given in Hollenbach & McKee (1989). We adopted other parameters from Hollenbach & McKee (1979)., and in McKee et al. (1982) for CO. We computed the population of H<sub>2</sub> following the procedure of Hollenbach & McKee (1979) using the collision coefficient given in Hollenbach & McKee (1989). For H<sub>2</sub>, we considered the first three vibrational states ($`v=02`$) with rotational levels up to $`J=20`$ in each vibrational state. We take into account only the ground vibrational level for other molecular lines since cooling rates owing to vibrational transitions are nearly always dominated by grain cooling or H<sub>2</sub> cooling (Hollenbach & McKee 1979).
#### 2.2.2 Gas-Grain Heat Transfer
We adopt the Pollack et al. (1994) model of grains in molecular clouds. In their model, silicates, organics, troilite, metallic iron, and water ice constitute the most abundant grain species and the dominant sources of opacity in molecular cloud cores. The Planck mean opacity due to the grains in the case of mean metallicity $`Z=Z_{\mathrm{local}}`$ is
$$\kappa _{\mathrm{gr}}=4.0\times 10^4T_{\mathrm{gr}}^2\mathrm{cm}^2\mathrm{g}^1,$$
(19)
in $`T_{\mathrm{gr}}50\mathrm{K}`$, where $`T_{\mathrm{gr}}`$ is the effective grain temperature (see below). The mass fraction of grain is $`0.934\times 10^2`$ below the water vaporization temperature (100-200 K, depending on the density), and that of each component; water ice, volatile organics, refractory organics, troilite, orthopyroxene, olivine and metallic iron are $`1.19\times 10^3,6.02\times 10^4,3.53\times 10^3,5.69\times 10^4,7.33\times 10^4,2.51\times 10^3`$, and $`2.53\times 10^4`$, respectively (from Table 2 of Pollack et al. 1994). For the grain size distribution
$$n(a)\{\begin{array}{cc}a^{3.5}\hfill & \text{(}0.005\mu \mathrm{m}<a<1\mu \mathrm{m}\text{)}\hfill \\ a^{5.5}\hfill & \text{(}1\mu \mathrm{m}<a<5\mu \mathrm{m}\text{)}\hfill \end{array}$$
(20)
(Pollack et al. 1994) that is based on that derived by Mathis, Rumpl, & Nordsieck (1977), the energy transfer rate from the gas to the dust grains per unit mass is given by
$$_{\mathrm{gr}}=1.1\times 10^5n(\frac{f_{\mathrm{gr}}}{\rho _{\mathrm{gr}}})(\frac{T}{1000\mathrm{K}})^{1/2}[10.8\mathrm{exp}(75\mathrm{K}/T)](TT_{\mathrm{gr}})$$
(21)
(Hollenbach & McKee 1979), where $`T_{\mathrm{gr}}`$ is an effective grain temperature, which is determined by the energy balance for dust grains;
$$_{\mathrm{gr}}=4\sigma T_{\mathrm{gr}}^4\kappa _{\mathrm{gr}}\beta _{\mathrm{cont}},$$
(22)
where $`\kappa _{\mathrm{gr}}`$ is the Planck mean opacity owing to the dust grains and $`f_{\mathrm{gr}}/\rho _{\mathrm{gr}}`$ is the total volume of dust grains per unit mass of gas and $`f_{\mathrm{gr}}/\rho _{\mathrm{gr}}=5.3\times 10^3`$ at the lowest temperatures (i.e., below the water ice vaporization temperature). These quantities are taken from tables in Pollack et al.(1994).
The continuum energy transport rate $`\beta _{\mathrm{cont}}`$ decreases as $`\tau _{\mathrm{cont}}^2`$ owing to radiative diffusion in the optically thick case (e.g., Masunaga et al. 1998). Then,
$$\beta _{\mathrm{cont}}=\mathrm{min}(1,\tau _{\mathrm{cont}}^2).$$
(23)
#### 2.2.3 Continuum of Gas
Using the Planck mean opacity of gas $`\kappa _{\mathrm{gas}}`$, the cooling rate owing to the gas continuum is given by
$$_{\mathrm{cont}}=4\sigma T^4\kappa _{\mathrm{gas}}\beta _{\mathrm{cont}}.$$
(24)
We take the continuum Planck mean opacity for metal-free gas from Lenzuni et al. (1991), which includes all the important continuum processes, specifically, bound-free absorption by $`\mathrm{H}^0`$ and $`\mathrm{H}^{}`$; free-free absorption by $`\mathrm{H}^0`$, $`\mathrm{H}^{}`$, $`\mathrm{H}_2`$, $`\mathrm{H}_2^{}`$, $`\mathrm{H}_2^+`$, $`\mathrm{H}_3^+`$, $`\mathrm{He}^0`$, $`\mathrm{He}^{}`$; photodissociation of $`\mathrm{H}_2`$, and $`\mathrm{H}_2^+`$ by thermal radiation; Rayleigh scattering by $`\mathrm{H}^0`$, $`\mathrm{H}_2`$, $`\mathrm{He}^0`$; Thomson scattering by $`e^{}`$; and collision-induced absorption by $`\mathrm{H}_2`$ due to collisions with $`\mathrm{H}_2`$, $`\mathrm{He}^0`$, and $`\mathrm{H}^0`$. Among the above processes, the most important cooling mechanism is H<sub>2</sub> collision-induced continuum, which dominates the cooling at about $`10^{16}\mathrm{cm}^3`$ for low metallicity (i.e., $`<10^6Z_{\mathrm{}}`$) clouds.
#### 2.2.4 Chemical Cooling/Heating
Following Hollenbach & McKee (1979), we assume the heat deposited per a formed molecular hydrogen as $`0.2+4.2(1+n_{\mathrm{cr}}/n)^1\mathrm{eV}`$ for H<sub>2</sub> formation on grain surfaces (reaction 23 in Appendix), $`3.53(1+n_{\mathrm{cr}}/n)^1\mathrm{eV}`$ for H<sub>2</sub> formation by H<sup>-</sup> process (reaction 8), $`1.83(1+n_{\mathrm{cr}}/n)^1\mathrm{eV}`$ for H<sub>2</sub> formation by H$`{}_{}{}^{+}{}_{2}{}^{}`$ process (reaction 10), and $`4.48(1+n_{\mathrm{cr}}/n)^1\mathrm{eV}`$ for H<sub>2</sub> formation by the three-body reactions (reactions 19 and 20), where
$$n_{\mathrm{cr}}=\frac{10^6T^{1/2}}{1.6y(\mathrm{H})\mathrm{exp}[(400/T)^2]+1.4y(\mathrm{H}_2)\mathrm{exp}[12000/(T+1200)]}\mathrm{cm}^3.$$
(25)
Collisional dissociation and ionization absorb the same energy as the binding energy, that is 4.48 eV per H<sub>2</sub> dissociation, 13.6 eV per H ionization, 24.6 eV per He ionization, and 54.4 eV per He<sup>+</sup> ionization.
### 2.3 Chemical reactions
We solve nonequilibrium chemistry involving the four elements H, He, C and O, that contains the following 45 species: $`\mathrm{H}`$, $`\mathrm{H}_2`$, $`\mathrm{e}^{}`$, $`\mathrm{H}^+`$, $`\mathrm{H}_2^+`$, $`\mathrm{H}_3^+`$, $`\mathrm{H}^{}`$, $`\mathrm{He}`$, $`\mathrm{He}^+`$, $`\mathrm{He}^{++}`$, $`\mathrm{HeH}^+`$, $`\mathrm{C}`$, $`\mathrm{C}_2`$, $`\mathrm{CH}`$, $`\mathrm{CH}_2`$, $`\mathrm{CH}_3`$, $`\mathrm{CH}_4`$, $`\mathrm{C}^+`$, $`\mathrm{C}_2^+`$, $`\mathrm{CH}^+`$, $`\mathrm{CH}_2^+`$, $`\mathrm{CH}_3^+`$, $`\mathrm{CH}_4^+`$, $`\mathrm{CH}_5^+`$, $`\mathrm{O}`$, $`\mathrm{O}_2`$, $`\mathrm{OH}`$, $`\mathrm{CO}`$, $`\mathrm{H}_2\mathrm{O}`$, $`\mathrm{HCO}`$, $`\mathrm{O}_2\mathrm{H}`$, $`\mathrm{CO}_2`$, $`\mathrm{H}_2\mathrm{CO}`$, $`\mathrm{H}_2\mathrm{O}_2`$, $`\mathrm{O}^+`$, $`\mathrm{O}_2^+`$, $`\mathrm{OH}^+`$, $`\mathrm{CO}^+`$, $`\mathrm{H}_2\mathrm{O}^+`$, $`\mathrm{HCO}^+`$, $`\mathrm{O}_2\mathrm{H}^+`$, $`\mathrm{H}_3\mathrm{O}^+`$, $`\mathrm{H}_2\mathrm{CO}^+`$, $`\mathrm{HCO}_2^+`$ and $`\mathrm{H}_3\mathrm{CO}^+`$. Chemical reactions included are listed in Appendix. The reactions of H and He chemistry are mainly taken from Abel et al. (1997) and Galli & Palla (1998). Other important H<sub>2</sub> forming processes included are the three-body reactions (Palla et al. 1983) and the reaction on the surfaces of the dust grains (Tielens & Hollenbach 1985). In addition, we supplement H, He, C and O chemical reactions involving the above species from Millar, Farquhar, & Willacy (1997).
## 3 Results
In this section, we present the results obtained by the method described in §2 and discuss our analysis.
Figure 1 displays the evolutionary trajectories of protostellar clouds whose metallicities are $`z=0,10^6,10^4,10^2,1`$. Here we set the initial condition to be $`T=100\mathrm{K},n=1\mathrm{c}\mathrm{m}^3,y(e)=1\times 10^4,y(\mathrm{H}_2)=1\times 10^6`$. All the carbon is assumed to be in the form of CII, while oxygen is OI at the beginning.
The evolutionary trajectories for two other initial conditions are shown in Figure 2, as well as the same curves in Figure 1. We can see from Figure 2 that trajectories of clouds with a fixed composition converge rapidly toward the dense region for any initial conditions (e.g., Hayashi & Nakano 1965; Low & Lynden-Bell 1976). We have also tested the sensitivity to initial chemical compostions for two cases – (i) where the ionization degree $`y(e)=1\times 10^3`$ and (ii) all the carbon is assumed to be CI,– and have found the results similar to those shown in Figure 2. Hereafter, we discuss only on the case presented in Figure 1.
In general, as seen in Figure 1, the temperatures of lower metallicity clouds are higher because of their lower radiative cooling ability (i.e., lower radiative cooling rate for the same temperature and density) as long as the clouds are transparent to continuum and continue to collapse dynamically, owing to the efficient radiative cooling, in other words, during what is known as the “first collapse” stage.
On account of the total lack of metals and grains, the only cooling agent in the temperature range below $`10^4`$K for primordial ($`z=0`$) clouds is rovibrational transitions of molecular hydrogen (e.g., Matsuda et al. 1969). In the primordial gas, H<sub>2</sub> is formed mainly by the H<sup>-</sup> process,
$`\mathrm{H}+e^{}`$ $``$ $`\mathrm{H}^{}+\gamma ;`$ (26)
$`\mathrm{H}+\mathrm{H}^{}`$ $``$ $`\mathrm{H}_2+e^{},`$ (27)
(Saslaw & Zipoy 1967) until the density reaches $`10^8\mathrm{cm}^3`$, where the three-body reactions
$$3\mathrm{H}\mathrm{H}_2+\mathrm{H}$$
(28)
$$2\mathrm{H}+\mathrm{H}_22\mathrm{H}_2$$
(29)
become efficient (Palla et al. 1983). H<sub>2</sub> line emission via electric quadrupole transitions is the only efficient cooling agent until the number density exceeds $`10^{14}\mathrm{cm}^3`$, where H<sub>2</sub> continuum emission via collision-induced dipole transitions begins to overwhelm. The drop in temperature at $`n10^{16}\mathrm{cm}^3`$ is due to the quadrature dependence of the collision-induced emission coefficient on the density. The cloud becomes optically thick to the continuum at $`10^{16}\mathrm{cm}^3`$. (See Fig. 3 a)
The gravitational contraction of metal-free protostellar clouds has been investigated with Omukai & Nishi (1998) by detailed hydrodynamical calculations assuming spherical symmetry. Our evolutionary trajectory for the primordial cloud agrees well with those of Omukai & Nishi (1998) at the number density $`n10^{10}\mathrm{cm}^3`$. At lower densities, the temperature of ours is higher owing to the influence of Omukai & Nishi’s (1998) initial condition. They started the calculation from a cloud in hydrostatic equilibrium at $`n10^6\mathrm{cm}^3`$. As long as the density is not so much higher than the initial value, the collapse is slower than the free-fall rate, which we have assumed in this paper. As a result, the temperature is lower because of the lower compressional heating.
The evolution of the cloud with $`z=10^6`$ is essentially the same as that of the primordial cloud, although there is minor deviation caused by extra cooling owing to water molecules and dust grains. (See Fig. 3 b) The dust thermal emission dominates the cooling just before the H<sub>2</sub> continuum becomes effective. However, it is only temporary and the grains rapidly sublimate at $`n10^{14}\mathrm{cm}^3`$. The most refractory grain composition that contributes to the Planck mean opacity is iron, which sublimate at about 1200K at this density (Table 3 of Pollack et al. 1994). The slight amount of dust grains in the cloud with $`z=10^6`$ does not contribute to H<sub>2</sub> formation significantly, as seen in Figure 4.
For the cloud with $`z=10^4`$, H<sub>2</sub>, which is formed mainly on the grain surfaces instead of the H<sup>-</sup> process, is the major coolant at low densities ($`<10^4\mathrm{cm}^3`$). Water molecules, which are mainly formed at $`10^6\mathrm{cm}^3`$ by the reaction
$$\mathrm{H}_2+\mathrm{OH}\mathrm{H}_2\mathrm{O}+\mathrm{H},$$
(30)
dominate the cooling in the range $`10^5\mathrm{cm}^3<n<10^{10}\mathrm{cm}^3`$. In present-day molecular cloud cores, usually neither H<sub>2</sub>O nor OH plays a key role in the thermal balance since the temperature there is lower than the excitation energy of these molecular species (Neufeld, Lepp, & Melnick 1995). However, in lower metal environments, the temperature is higher because of the lower radiative cooling. In that case, H<sub>2</sub>O can be a principal cooling agent. In the cloud with $`z=10^4`$, grain surface reactions exceed the H<sup>-</sup> process as a mode of H<sub>2</sub> formation, although hydrogen does not become fully molecular until the three-body reaction begins to work. At $`10^{10}\mathrm{cm}^3`$, the grain thermal emission dominates the cooling and causes the temperature to drop to about 100K. The cloud becomes optically thick to the grain emission at about $`10^{13}\mathrm{cm}^3`$. (See Fig. 3 c)
For the cloud with $`z>10^2`$, cooling is dominated by atomic lines and CO lines in lower densities (that is, $`n<10^6\mathrm{cm}^3`$ for $`z=10^2`$, and $`n<10^3\mathrm{cm}^3`$ for $`z=1`$, respectively), and by the dust thermal emission in higher densities. (See Fig. 3 d, e) The temperature of our $`z=1`$ cloud drops far less than 10K, which is appropriate in present-day protostellar clouds. This is merely because we have ignored external radiation, and the results should not be taken seriously as a model of present-day molecular cloud cores.
Clouds with metallicity $`z>10^4`$ become opaque to the thermal radiation of dust grains before the grain vaporization occurs. Once the cloud becomes opaque, it begins to contract adiabatically since the radiative cooling rate drops rapidly. According to some hydrodynamical simulations (e.g, Larson 1969; Masunaga et al. 1998), a transient core in hydrostatic equilibrium (called a “first core”) forms after the cloud has more contractions. The hydrostatic core continues the adiabatic contraction by accreting the envelope matter until the temperature reaches $`2000`$ K, where molecular hydrogen begins to dissociate.
Clouds with lower metallicity (i.e., $`z<10^6`$) become opaque to the H<sub>2</sub> continuum instead of the dust thermal emission, as described above. Since the temperature at that time is already near the dissociation value, molecular hydrogen begins to dissociate just after the cloud becomes opaque. As a result, no transient hydrostatic core forms in these cases as has been demonstrated by the hydrodynamical calculations of Omukai & Nishi (1998).
When the ratio of specific heats $`\mathrm{\Gamma }=\mathrm{dlog}p/\mathrm{dlog}\rho `$ exceeds the critical value 4/3, corresponding to a gradient of 1/3 in the density-temperature ($`nT`$) plane (i.e., Figure 1), the bulk motion of the cloud is decelerated and the free-fall assumption is no longer valid. However, the trajectories in the $`nT`$ plane are not altered since the clouds contract almost adiabatically anyway in the opaque stage because of the rapidly declining radiative cooling rate (e.g., Larson 1969; Narita, Nakano, & Hayashi 1970).
In the course of the adiabatic contraction of the transient cores, all the evolutionary trajectories converge to a certain line (Fig. 1, dashed line) in spite of their different composition and histories. We discuss the reason for the convergence in the following. The trajectories in the $`nT`$ plane just prior to when the clouds become opaque are determined by thermal balance between the compressional heating and the radiative cooling. The compressional heating rate (per unit mass) is given by
$$𝒢_{\mathrm{comp}}=p\frac{d}{dt}(\frac{1}{\rho })=\frac{c_\mathrm{s}^2}{t_{\mathrm{ff}}},$$
(31)
where $`p`$ is the pressure and $`c_\mathrm{s}`$ is the sound speed in the central region of the cloud. The radiative cooling at that time is dominated by continuum for all clouds, namely, the dust thermal emission for clouds with relative metallicity $`z>10^4`$, or the H<sub>2</sub> continuum emission for those with lower metallicity. The radiative cooling rate is given by
$$_{\mathrm{rad}}=4\kappa _\mathrm{B}\sigma T^4,$$
(32)
where $`\kappa _\mathrm{B}`$ is the Planck mean continuum opacity. In the equation above, we have used the relation that the grain temperature is the same as the gas temperature, $`T_{\mathrm{gr}}=T`$, which is valid in a cloud so dense that it becomes opaque to the continuum. Note that expression (32) holds for any clouds that cool mainly by continuum by choosing an appropriate value of opacity, where the opacity is an increasing function of relative metallicity $`z`$. Equating (31) and (32), for each value of $`\kappa _\mathrm{B}(z)`$, we obtain a trajectory just before the cloud becomes opaque (we call it Tr($`\kappa _\mathrm{B}`$));
$$\mathrm{Tr}(\kappa _\mathrm{B}):\frac{c_\mathrm{s}^2}{t_{\mathrm{ff}}}=4\kappa _\mathrm{B}\sigma T^4.$$
(33)
As the cloud of opacity $`\kappa _\mathrm{B}`$ collapses along the trajectory represented by (33), the central part of the cloud becomes opaque to the continuum at the point
$$\mathrm{P}(\kappa _\mathrm{B}):\tau _\mathrm{J}(\kappa _\mathrm{B},n,T)=\kappa _\mathrm{B}\rho \pi \frac{c_\mathrm{s}}{\sqrt{G\rho }}=1,$$
(34)
where we have used the assumption that the size of the central region is approximately $`\lambda _\mathrm{J}`$. Eliminating $`\kappa _\mathrm{B}`$ in equations (33) and (34), we obtain the locus L of points P for variable $`\kappa _\mathrm{B}`$;
$$\mathrm{L}:T=(\frac{k^3}{12\sigma ^2m_\mathrm{p}})^{1/5}n^{2/5}.$$
(35)
On the other hand, the trajectory of the cloud of opacity $`\kappa _\mathrm{B}`$ in the opaque stage, O($`\kappa _\mathrm{B}`$), is the line that starts from the point P($`\kappa _\mathrm{B}`$) and whose gradient is approximately equal to the value for adiabatic contraction $`\gamma _{\mathrm{ad}}1`$. The adiabatic exponent $`\gamma _{\mathrm{ad}}`$ for molecular hydrogen is 7/5 in the temperature range $`T200`$K, while it is 2/3 in the lowest temperature. The gradient of the trajectory O($`\kappa _\mathrm{B}`$) in the range $`T200`$K is 2/5. This value happens to be the same as the gradient of the line L. Thus, the trajectories in the opaque stage O all coincide with the line L, which can be written only by physical constants. When the clouds become opaque (i.e., at the point P), their specific entropy has the same value for all clouds regardless of their metallicity.
As a cloud climbs up along line L, the temperature reaches the dissociation value ($`2000`$K) at $`n3\times 10^{16}\mathrm{cm}^3`$. The ratio of specific heats are reduced below the critical value 4/3. Then the cloud begins to collapse dynamically again, or in other words, it begins the second collapse (e.g., Larson 1969). The trajectories in the second collapse stage are also unique because the specific entropy has the same value for all clouds regardless of their metallicity. When most hydrogen molecules have been dissociated, the ratio of specific heats rises above the critical value 4/3 again. After the central part of the cloud contracts almost adiabatically to some extent, a hydrostatic core forms. This core is referred to as the stellar or second core in literature. We cannot accurately estimate the size of the stellar core from our simple one-zone treatment. However, hydrodynamical calculations tell us that the number density and mass of the stellar core at its formation time are on the order of $`10^{22}\mathrm{cm}^3`$ and $`10^3M_{\mathrm{}}`$ respectively, both for present-day (Larson 1969) and primordial protostars (Omukai & Nishi 1998). This value is expected to be universal, namely, independent of metallicity since the equations of state, or trajectories in the $`nT`$ plane, are the same in the second collapse stage.
The stellar core, or protostar, grows in mass by accretion of the envelope material. The temperature rises accordingly and eventually the core will become an ordinary star. Consequently, the final mass of stars is determined by the subsequent accretion onto the core, although the size of the protostar at the time of formation is the same. Stahler, Shu, & Taam (1980) argued that a rough estimate of the protostellar mass accretion rate $`\dot{M}`$ can be obtained from the relation
$$\dot{M}c_\mathrm{s}^3/G,$$
(36)
where $`c_\mathrm{s}`$ is the isothermal sound speed in the initial protostellar cloud. Therefore, the mass accretion rate is higher for a lower metallicity cloud because of the higher temperature of the protostellar cloud (Stahler, Palla, & Salpeter 1986). The higher accretion rate and lower opacity (then smaller radiation force) may result in the higher mass of formed stars for a lower metallicity cloud (e.g., Wolfire & Cassinelli 1987). However, to address these issues fully, further investigations are needed.
## 4 Summary
We have investigated the thermal and chemical evolution of collapsing protostellar clouds with different metallicities. The varied range of metallicity spans the local interstellar value ($`1Z_{\mathrm{}}`$) to zero.
The evolution of the clouds is summarized as follows. While the clouds are transparent to continuous radiation, the temperature of clouds with lower metallicity is higher since their radiative cooling rates are lower for the same density and temperature. However, after the clouds become opaque and begin adiabatic contraction, their evolutionary trajectories converge to a line that is determined only by physical constants. Thereafter, the trajectories coincide with each other regardless of their metallicity. Consequently, the physical dimension of the stellar core at the time of formation is the same for clouds with any composition.
We have also discussed analytically the reasons for the convergence.
The author wishes to express his cordial thanks to Humitaka Sato and Naoshi Sugiyama for their continual interest and advice, to Hiroshi Koyama, Ryoichi Nishi, and Toru Tsuribe for useful discussions, and to the referee for improving this manuscript.
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# References
A Novel Approach to Spherical Harmonics Expansion for Electron Transport in Semiconductors
S. F. Liotta and A. Majorana
Dipartimento di Matematica - University of Catania -
Viale A. Doria 6, 95125 Italy
Abstract. A set of equations is derived from the Boltzmann kinetic equation describing charge transport in semiconductors. The unknowns of these equations depend on the space-time coordinates and the electron energy. The non-parabolic and parabolic band approximation are treated in detail. In these cases, the set of equations is equivalent to those obtained in the spherical-harmonics expansion. Stationary and homogeneous solutions are explicitly treated. In order to solve numerically the equations, a cut in the energy range is introduced. The modified model maintains the physical characteristics of the original equations. The solution of an asymptotic equation is found and compared to the numerical solutions.
Keywords: Boltzmann equation; Semiconductors; Spherical harmonics expansion
1. Introduction
The motion of electrons in a crystal is governed by complex physical laws so that accurate models are required to achieve correct results. The drift-diffusion model has been widely used in the past, and it has provided a good description of relevant physical mechanisms. In modern devices, whose sizes are in the submicron range, non-equilibrium effects play an important role and are not adequately modeled by the drift-diffusion approach. Since in many processes a more accurate description than the hydrodynamical setting is required, Boltzmann equation or Monte Carlo simulations are employed. A fully kinetic treatment of carrier dynamics guarantees accurate results but requires very expensive numerical procedures in order to solve realistic problems.
To reduce the complexity of the use of the full Boltzmann equation, many authors (see ref. and references therein) have introduced simpler models, assuming particular forms for the probability density function. The aim is to yield results, which are more accurate than those obtained by hydrodynamical models but less expensive than direct particle simulations or numerical treatment of the Boltzmann equation. In this framework, a well-known model is derived using a spherical-harmonics expansion , , , from the Boltzmann equation.
The paper is organized as follows. In Sec. 2, we introduce the Boltzmann Equation. The collisional operator takes into account the interactions between electrons and phonons, i.e. vibrations of the lattice. This is assumed to be in thermal equilibrium. A set of kinetic equations is derived in Sec. 3. These are obtained for a general expression of the microscopic kinetic energy, so that, for example, they can be used also for non-parabolic bands. The equations are of partial differential and difference type. Equations are specialized in Sec. 4 in the parabolic band approximation, and is Sec. 5, where the Kane band model is assumed. In both cases we show that the equations are equivalent to those obtained in the spherical harmonics framework. In Sec. 6 we write the equations in the stationary and homogeneous case. The rest of the paper is devoted to the study of these equations. In Sec. 7 an asymptotic equation, valid for large value of the electric field, is found and analyzed; the explicit solution is found. In Sec. 8 we discuss some difficulties related to the infinite energy range and we propose a suitable solution by introducing a cut in the scattering rate of the collisional operator. The numerical scheme to solve the equations is presented in Sec. 9, and the results are shown in the last section.
Throughout the paper, boldface and lightface symbols denote vectors and scalar quantities respectively.
2. Basic equations
We consider an electron gas, which moves in a lattice, subjects to an external electric field $`𝐄`$. This can be applied or related, by Poisson equation, to the density of the gas. For the electron gas in a semiconductor device the Boltzmann equation , writes
(1)
$$\frac{f}{t}+\frac{1}{\mathrm{}}_𝐤\epsilon _𝐱f\frac{\text{e}}{\mathrm{}}𝐄_𝐤f=Q(f),$$
where the unknown function $`f(t,𝐱,𝐤)`$ represents the probability of finding an electron at the position $`𝐱\text{}^3`$, with the wave-vector $`𝐤\text{}^3`$, at time $`t`$. The parameters $`\mathrm{}`$ and e are the Planck constant divided by $`2\pi `$ and the positive electric charge, respectively. The symbol $`_𝐤`$ stands for the gradient with respect to the variables $`𝐤`$ and $`_𝐱`$ that with respect to the space coordinates $`𝐱`$. The particle energy $`\epsilon `$ is an assigned nonnegative continuous function, defined on $`\text{}^3`$, of the wave-vector $`𝐤`$. Here we assume that the low-density approximation holds, so that $`Q`$ is linear in $`f`$. If $`K(𝐤,\stackrel{~}{𝐤})`$ is the sum of the scattering kernels, which describe the nature of the inelastic collisions (for example, electron–polar or non-polar optical phonon scattering), and $`K_0(𝐤,\stackrel{~}{𝐤})`$ is the kernel for the elastic collisions (for example, electron–impurity scattering), then the collisional operator writes , ,
(2) $`Q(f)=(𝗇_q+1){\displaystyle _\text{}^3}K(𝐤,\stackrel{~}{𝐤})\delta (\epsilon (\stackrel{~}{𝐤})\epsilon (𝐤)\mathrm{}\omega )f(t,𝐱,\stackrel{~}{𝐤})𝑑\stackrel{~}{𝐤}`$
$`+𝗇_q{\displaystyle _\text{}^3}K(\stackrel{~}{𝐤},𝐤)\delta (\epsilon (\stackrel{~}{𝐤})\epsilon (𝐤)+\mathrm{}\omega )f(t,𝐱,\stackrel{~}{𝐤})𝑑\stackrel{~}{𝐤}`$
$`+{\displaystyle _\text{}^3}K_0(𝐤,\stackrel{~}{𝐤})\delta (\epsilon (\stackrel{~}{𝐤})\epsilon (𝐤))f(t,𝐱,\stackrel{~}{𝐤})𝑑\stackrel{~}{𝐤}\overline{\nu }(𝐤)f(t,𝐱,𝐤),`$
where
(3) $`\overline{\nu }(𝐤)=𝗇_q{\displaystyle _\text{}^3}K(𝐤,\stackrel{~}{𝐤})\delta (\epsilon (\stackrel{~}{𝐤})\epsilon (𝐤)\mathrm{}\omega )𝑑\stackrel{~}{𝐤}`$
$`+(𝗇_q+1){\displaystyle _\text{}^3}K(\stackrel{~}{𝐤},𝐤)\delta (\epsilon (\stackrel{~}{𝐤})\epsilon (𝐤)+\mathrm{}\omega )𝑑\stackrel{~}{𝐤}+{\displaystyle _\text{}^3}K_0(\stackrel{~}{𝐤},𝐤)\delta (\epsilon (\stackrel{~}{𝐤})\epsilon (𝐤))𝑑\stackrel{~}{𝐤}.`$
Often, in the following we omit to write explicitly that $`f`$ depends also on $`t`$ and $`𝐱`$, whereas we do not operate with these variables.
The constant quantity $`𝗇_q`$ represents the thermal-equilibrium number of phonons and is given by
$$𝗇_q=\frac{1}{\mathrm{exp}\left(\frac{\mathrm{}\omega }{k_BT_L}\right)1}$$
where $`k_B`$ is the Boltzmann constant and $`T_L`$ is the constant lattice temperature.
The symbol $`\delta (\epsilon (\stackrel{~}{𝐤})\epsilon (𝐤)\pm \mathrm{}\omega )`$ means the composition of the real-valued function $`\epsilon (\stackrel{~}{𝐤})\epsilon (𝐤)\pm \mathrm{}\omega `$ and the Dirac distribution $`\delta `$. Since the function $`\epsilon `$ may have many different expressions, it is not possible in general, by using a unique explicit technique, to transform each of the integral operators (S0.Ex2) in equivalent operators without Dirac distributions. Therefore, we shall maintain the original form of (S0.Ex2) in the following.
As usual, due to the isotropy of phase-space, the kernels $`K(𝐤,\stackrel{~}{𝐤})`$ and $`K_0(𝐤,\stackrel{~}{𝐤})`$ are symmetric with respect to the wave-vectors $`𝐤`$ and $`\stackrel{~}{𝐤}`$. Moreover, since the kernels are scalar, they depend on $`𝐤`$ and $`\stackrel{~}{𝐤}`$ only through the scalar quantities $`\epsilon (𝐤)`$, $`\epsilon (\stackrel{~}{𝐤})`$ and $`𝐤\stackrel{~}{𝐤}`$. This implies in particular that the collision frequency $`\nu `$ depends on $`𝐤`$ only through the variable $`\epsilon `$.
3. Energy-kinetic equations
In the framework of the Boltzmann equation for a perfect gas, a classical procedure, due to Grad , to reduce the dimension of the space of the independent coordinates, is the moment method. It consists in expanding the ratio between the probability density $`f`$ and a local maxwellian function. Inserting this expansion in the Boltzmann equation, an infinite sequence of differential equations are obtained. In general no finite set of equations is uncoupled from the rest of the system. To obtain a determined system up to certain order $`m`$ (usually 13) only $`m`$ equation are considered and the coefficients of the expansion of higher order are assumed negligible.
In order to obtain a new set of equations, where the unknown functions depend on the space-time coordinates and the microscopic kinetic energy, we follow the scheme analyzed in . This scheme partially recalls the moment method.
If $`\text{P}(𝐤)`$ is a polynomial and $`u`$ a real variable, we multiply the Boltzmann equation (1) by $`\text{P}(𝐤)\delta (\epsilon (𝐤)u)`$ and then formally integrate with respect to the variable $`𝐤`$ over the whole space $`\text{}^3`$. It is possible to show that the following equation can be derived:
(4) $`{\displaystyle \frac{\text{ }}{t}}{\displaystyle _\text{}^3}f(t,𝐱,𝐤)\text{P}(𝐤)\delta (\epsilon (𝐤)u)𝑑𝐤`$
$`+_𝐱{\displaystyle _\text{}^3}𝐯(𝐤)f(t,𝐱,𝐤)\text{P}(𝐤)\delta (\epsilon (𝐤)u)𝑑𝐤`$
$`\text{e}𝐄{\displaystyle \frac{\text{ }}{u}}{\displaystyle _\text{}^3}f(t,𝐱,𝐤)\text{P}(𝐤)𝐯(𝐤)\delta (\epsilon (𝐤)u)𝑑𝐤`$
$`+{\displaystyle \frac{\text{e}}{\mathrm{}}}𝐄{\displaystyle _\text{}^3}f(t,𝐱,𝐤)\left[_𝐤\text{P}(𝐤)\right]\delta (\epsilon (𝐤)u)𝑑𝐤={\displaystyle _\text{}^3}Q(f)\text{P}(𝐤)\delta (\epsilon (𝐤)u)𝑑𝐤.`$
Here the molecular velocity is defined by means of the formula
(5)
$$𝐯(𝐤):=\frac{1}{\mathrm{}}_𝐤\epsilon (𝐤).$$
For any choice of the function $`\text{P}(𝐤)`$ a new equation is obtained. As it happens for the moment method, each equation contains, except for particular cases at least two unknown functions. For example, if we choose $`\text{P}(𝐤)1`$, then the left hand side of eq. (S0.Ex7) contains in general the following unknowns
(6) $`N(t,𝐱,u)`$ $`:=`$ $`{\displaystyle _\text{}^3}f(t,𝐱,𝐤)\delta (\epsilon (𝐤)u)𝑑𝐤,`$
(7) $`𝐕(t,𝐱,u)`$ $`:=`$ $`{\displaystyle _\text{}^3}f(t,𝐱,𝐤)𝐯(𝐤)\delta (\epsilon (𝐤)u)𝑑𝐤.`$
The scalar quantity $`N(t,𝐱,u)`$ is the probability density function to find an electron with energy $`u`$ at time $`t`$ and position $`𝐱`$. The function $`N`$ is non-negative, because the probability density $`f`$ is always non-negative. We note that the presence of the delta-function implies that $`N`$ and $`𝐕`$ vanish for all $`(t,𝐱,u)`$ such that $`uu_0=\mathrm{min}\left\{\epsilon (𝐤):𝐤\text{}^3\right\}`$. If the function $`N`$ is known, then we can obtain all the hydrodynamical scalars. For example, to find the hydrodynamical density $`\rho (t,𝐱)`$, it is sufficient to integrate $`N`$ with respect to all values of energy $`u`$; i.e.
$$\rho (t,𝐱):=_{u_0}^+\mathrm{}N(t,𝐱,u)𝑑u.$$
The physical meaning of the other quantity $`𝐕(t,𝐱,u)`$ can be understood with a similar argument.
In the rest of this article we assume the function $`\epsilon (𝐤)`$ to be even (i.e. $`\epsilon (𝐤)=\epsilon (𝐤)`$) and the kernels $`K`$ and $`K_0`$ to depend only on $`\epsilon (𝐤)`$ and $`\epsilon (\stackrel{~}{𝐤})`$. The first hypothesis is verified for the most common models used in the numerical simulations. The second are introduced in order to simplify the treatment of the equations. Let
$$\text{K}(\epsilon (𝐤),\epsilon (\stackrel{~}{𝐤})):=K(𝐤,\stackrel{~}{𝐤})\text{ and }\text{K}_0(\epsilon (𝐤),\epsilon (\stackrel{~}{𝐤})):=K_0(𝐤,\stackrel{~}{𝐤}).$$
Then eqs. (S0.Ex2)-(S0.Ex3) become
$`Q(f)`$ $`=`$ $`(𝗇_q+1)\text{K}(\epsilon (𝐤),\epsilon (𝐤)+\mathrm{}\omega )N(t,𝐱,\epsilon (𝐤)+\mathrm{}\omega )`$
$`+𝗇_q\text{K}(\epsilon (𝐤)\mathrm{}\omega ,\epsilon (𝐤))N(t,𝐱,\epsilon (𝐤)\mathrm{}\omega )`$
$`+\text{K}_0(\epsilon (𝐤),\epsilon (𝐤))N(t,𝐱,\epsilon (𝐤))\nu (\epsilon (𝐤))f(t,𝐱,𝐤)`$
with
$`\nu (\epsilon (𝐤))`$ $`=`$ $`(𝗇_q+1)\text{K}(\epsilon (𝐤)\mathrm{}\omega ,\epsilon (𝐤))\sigma (\epsilon (𝐤)\mathrm{}\omega )`$
$`+𝗇_q\text{K}(\epsilon (𝐤),\epsilon (𝐤)+\mathrm{}\omega )\sigma (\epsilon (𝐤)+\mathrm{}\omega )+\text{K}_0(\epsilon (𝐤),\epsilon (𝐤))\sigma (\epsilon (𝐤)),`$
where
(10)
$$\sigma (u):=_\text{}^3\delta (\epsilon (𝐤)u)𝑑𝐤$$
is the density of states.
We are interested in deducing two equations in the only variables $`N`$ and $`𝐕`$. By choosing $`\text{P}(𝐤)1`$ and $`\text{P}(𝐤)=𝐯(𝐤)`$ eqs. (S0.Ex7), (S0.Ex11) and (S0.Ex12) give the following equations
(11) $`{\displaystyle \frac{N}{t}}+_𝐱𝐕\text{e}𝐄{\displaystyle \frac{𝐕}{u}}=G(N)`$
(12) $`{\displaystyle \frac{V_i}{t}}+{\displaystyle \underset{j=1}{\overset{3}{}}}\left[{\displaystyle \frac{\mathrm{\Pi }_{ij}}{x_j}}\text{e}E_j{\displaystyle \frac{\mathrm{\Pi }_{ij}}{u}}+{\displaystyle \frac{\text{e}}{\mathrm{}}}E_j\mathrm{\Delta }_{ij}\right]=\nu (u)V_i(i=1,2,3),`$
where
(13) $`G(N):=(𝗇_q+1)\text{K}(u,u+\mathrm{}\omega )\sigma (u)N(t,𝐱,u+\mathrm{}\omega )`$
$`+𝗇_q\text{K}(u,u\mathrm{}\omega )\sigma (u)N(t,𝐱,u\mathrm{}\omega )`$
$`\left[𝗇_q\text{K}(u,u+\mathrm{}\omega )\sigma (u+\mathrm{}\omega )+(𝗇_q+1)\text{K}(u,u\mathrm{}\omega )\sigma (u\mathrm{}\omega )\right]N(t,𝐱,u),`$
(14) $`\nu (u)=𝗇_q\text{K}(u,u+\mathrm{}\omega )\sigma (u+\mathrm{}\omega )`$
$`+(𝗇_q+1)\text{K}(u,u\mathrm{}\omega )\sigma (u\mathrm{}\omega )+\text{K}_0(u,u)\sigma (u),`$
(15) $`\mathrm{\Pi }_{ij}:={\displaystyle _\text{}^3}f(t,𝐱,𝐤)v_i(𝐤)v_j(𝐤)\delta (\epsilon (𝐤)u)𝑑𝐤,`$
(16) $`\mathrm{\Delta }_{ij}:={\displaystyle _\text{}^3}f(t,𝐱,𝐤){\displaystyle \frac{v_i(𝐤)}{k_j}}\delta (\epsilon (𝐤)u)𝑑𝐤.`$
Also the tensor $`\mathrm{\Delta }_{ij}`$ is symmetric, due to eq. (5). We note that the effects of the elastic collisions, through the kernel $`\text{K}_0`$, disappear in the operator $`G(N)`$.
The set of eqs. (11)-(12) contains $`N`$ and $`𝐕`$ but also two tensors $`\mathrm{\Delta }_{ij}`$ and $`\mathrm{\Pi }_{ij}`$. Therefore, except in the case of spatially homogeneous solutions with null electric field, we need to assume some relations (usually called constitutive equations in the thermodynamic framework), which link $`\mathrm{\Delta }_{ij}`$ and $`\mathrm{\Pi }_{ij}`$ with $`N`$ and $`𝐕`$. These relations should also depend on the form of the microscopic kinetic energy $`\epsilon `$.
We feel that at this stage it would be worthwhile to make a simple observation. We note that the linearity of the Boltzmann equation (1) is maintained in eqs. (11)-(12). To maintain this feature, we suggest to assume constitutive relations of the kind
(17) $`\mathrm{\Pi }_{ij}(t,𝐱,u)`$ $`=`$ $`p_{ij}(u)N(t,𝐱,u)+p_i(u)V_j(t,𝐱,u)+p_j(u)V_i(t,𝐱,u)`$
(18) $`\mathrm{\Delta }_{ij}(t,𝐱,u)`$ $`=`$ $`d_{ij}(u)N(t,𝐱,u)+d_i(u)V_j(t,𝐱,u)+d_j(u)V_i(t,𝐱,u),`$
where the symmetric tensors $`p_{ij}`$ and $`d_{ij}`$ and the vectors $`p_i`$ and $`d_i`$ must be determined. In the next sections we propose a simple choice in the cases of the parabolic and non-parabolic band approximation.
4. The parabolic case
The simplest and widely used expression for the microscopic kinetic energy is
(19)
$$\epsilon (𝐤)=\frac{\mathrm{}^2}{2m^{}}k^2,$$
where $`m^{}`$ is the value of the effective electron mass in the parabolic mass approximation. In this case the tensor $`\mathrm{\Delta }_{ij}`$ reduces to
(20)
$$\mathrm{\Delta }_{ij}=\frac{\mathrm{}}{m^{}}N(t,𝐱,u)\delta _{ij}(i,j=1,2,3),$$
where $`\delta _{ij}`$ is the Kroneker symbol; so that no further assumption on eq. (18) is required. The tensor $`\mathrm{\Pi }_{ij}`$ is determined by assuming that (17) holds with $`p_i(u)0`$ and moreover that the constitutive relation is exact if $`f`$ depends on $`𝐤`$ only through the variable $`\epsilon `$. For $`f(t,𝐱,𝐤)=\stackrel{~}{f}(t,𝐱,\epsilon (𝐤))`$, we have
$`N(t,𝐱,u)`$ $`=`$ $`\stackrel{~}{f}(t,𝐱,u)\sigma (u)`$
$`\mathrm{\Pi }_{ij}(t,𝐱,u)`$ $`=`$ $`\stackrel{~}{f}(t,𝐱,u){\displaystyle _\text{}^3}v_i(𝐤)v_j(𝐤)\delta (\epsilon (𝐤)u)d𝐤(i,j=1,2,3),`$
so that
$$\sigma (u)p_{ij}(u)=_\text{}^3v_i(𝐤)v_j(𝐤)\delta (\epsilon (𝐤)u)𝑑𝐤.$$
Now, it is easy to verify that
$$\sigma (u)=4\sqrt{2}\pi \left(\frac{\sqrt{m^{}}}{\mathrm{}}\right)^3\theta (u)\sqrt{u},$$
where $`\theta `$ is the Heaviside step function. Then, for every $`u0`$,
$$_\text{}^3v_i(𝐤)v_j(𝐤)\delta (\epsilon (𝐤)u)d𝐤=\frac{2}{3}\frac{u}{m^{}}\delta _{ij}\sigma (u)(i,j=1,2,3).$$
Hence, eq. (12) becomes
(21)
$$\frac{𝐕}{t}+\frac{1}{3m^{}}\left[2u_𝐱N2\text{e}𝐄\frac{\text{ }}{u}(uN)+3\text{e}𝐄N\right]=\nu (u)𝐕.$$
In the following the functions $`N`$ and $`𝐕`$ will be the unknowns, which must satisfy eqs. (11) and (21). It is evident that the knowledge of these quantities cannot provide the probability density $`f`$ uniquely. Also, it is very hard to prove that, if $`N`$ and $`𝐕`$, at the initial time, derive from a probability density $`f`$, there exists a non-negative probability density, which gives $`N`$ and $`𝐕`$ for all times.
Actually the situation is similar to that of the moments in the hydrodynamical equations derived by the Grad method. Here usually, not even for the initial data, the existence of a probability density, which gives all the moments occurring in the equations, is considered.
Therefore, we limit to assuming only the simple conditions that $`N`$ is nonnegative and both the variables $`N`$ and $`𝐕`$ vanish as $`u0`$.
We note that these unknowns are simply related to the first two terms of the spherical harmonic expansion of the distribution function
(22)
$$f(t,𝐱,𝐤)=f_0(t,𝐱,\epsilon (𝐤))+𝐤𝐟_1(t,𝐱,\epsilon (𝐤))+\mathrm{}.$$
It is a simple matter to show that
$`f_0(t,𝐱,\epsilon (𝐤))`$ $`=`$ $`{\displaystyle \frac{N(t,𝐱,\epsilon (𝐤))}{\sigma (\epsilon (𝐤))}},`$
$`𝐟_1(t,𝐱,\epsilon (𝐤))`$ $`=`$ $`{\displaystyle \frac{3\mathrm{}}{2\epsilon (𝐤)\sigma (\epsilon (𝐤))}}𝐕(t,𝐱,\epsilon (𝐤)).`$
The corresponding equations for $`f_0`$ and $`𝐟_1`$ (see ref. ) coincide with our eqs. (11) and (21). This is a consequence of the choice on the approximation of the tensor $`\mathrm{\Pi }_{ij}`$. It is evident that different choices give different sets of equations. In this paper we do not investigate this possibility.
Similar equations were proposed by Hänsch in ref. in the framework of quantum kinetic transport. Also in that case a problem of equating the number of equations and unknowns arose. It was solved by simple physical arguments, but the equations are different from eqs. (11)-(21).
5. The non-parabolic band approximation
For large values of electron energy eq. (19) is not adequate. The following suitable expression was given by Kane :
(23)
$$\epsilon (1+\alpha _k\epsilon )=\frac{\mathrm{}^2}{2m^{}}k^2,$$
where $`\alpha _k`$ is a constant. Now, it is simple to verify that
$$\sigma (u)=4\sqrt{2}\pi \left(\frac{\sqrt{m^{}}}{\mathrm{}}\right)^3\theta (u)\sqrt{u(1+\alpha _ku)}(1+2\alpha _ku).$$
So eq. (11) can be written explicitly. The second equation requires only to determine $`\mathrm{\Pi }_{ij}`$ and $`\mathrm{\Delta }_{ij}`$ according to eqs. (17)-(18). Due to the spherical symmetry of the band we assume $`p_j(u)=d_j(u)=0`$. Analogously to the parabolic case, we require that eqs. (17)-(18) hold if $`f`$ depends on $`𝐤`$ only through $`\epsilon `$. In the same manner we can see that that
$$p_{ij}(u)=g(u)\delta _{ij},d_{ij}(u)=h(u)\delta _{ij}$$
where
$$g(u)=\frac{2}{3}\frac{u(1+\alpha _ku)}{m^{}(1+2\alpha _ku)^2}$$
and
$$h(u)=\frac{1}{m^{}(1+2\alpha _ku)}\frac{4\alpha _k}{3m^{}}\frac{u(1+\alpha _ku)}{(1+2\alpha _ku)^3}.$$
Therefore eq. (12) becomes
(24)
$$\frac{𝐕}{t}+g(u)_𝐱N\text{e}𝐄\frac{\text{ }}{u}\left[g(u)N\right]+\text{e}𝐄h(u)N=\nu (u)𝐕.$$
Also in this case the equations for $`N`$ and $`V`$ are equivalent to those obtained by means of the spherical harmonics expansion . For $`\alpha _k=0`$ we recover eq. (21).
6. Homogeneous solutions with constant electric field
We consider equations (11) and (24), and we look for a solution depending only on $`u`$. Since in this case the electric field must be a constant vector, we choose the reference frame so that only the first component of $`𝐄`$ is different from zero. As a consequence the only significant component of $`𝐕`$ is the first. We limit to assume constant kernels K and $`\text{K}_0`$. This fact holds when a homogeneous, intrinsic silicon at room temperature is considered.
It is useful to introduce dimensionless variables. Let be
$$t_{}:=\left[4\sqrt{2}\pi \left(\frac{\sqrt{m^{}}}{\mathrm{}}\right)^3\sqrt{k_BT_L}\text{K}𝗇_q\right]^1,\mathrm{}_{}:=\sqrt{\frac{k_BT_L}{m^{}}}t_{},u_{}:=k_BT_L,$$
$$w:=\frac{u}{u_{}},n(w):=u_{}\mathrm{}_{}^3N(u),v(w):=u_{}\mathrm{}_{}^2t_{}V_1(u),$$
$$\alpha :=\frac{\mathrm{}\omega }{k_BT_L},a:=\frac{𝗇_q+1}{𝗇_q}=e^\alpha ,\beta =\frac{\text{K}_0}{𝗇_q\text{K}},\stackrel{~}{\alpha }:=\alpha _ku_{},\zeta :=\text{e}E\frac{\mathrm{}_{}}{u_{}}.$$
It is a simple matter that the following equations can be obtained
(25) $`\zeta {\displaystyle \frac{dv}{dw}}=s(w)\left[an(w+\alpha )+n(w\alpha )\right]\left[s(w+\alpha )+as\left(w\alpha \right)\right]n(w)`$
(26) $`\zeta \left[{\displaystyle \frac{d\text{ }}{dw}}(r(w)n)+q(w)n\right]=\left[s(w+\alpha )+as\left(w\alpha \right)+\beta s(w)\right]v(w),`$
where
(27) $`s(w)`$ $`=`$ $`\theta (w)\sqrt{w(1+\stackrel{~}{\alpha }w)}(1+2\stackrel{~}{\alpha }w)`$
(28) $`r(w)`$ $`=`$ $`{\displaystyle \frac{2}{3}}{\displaystyle \frac{w(1+\stackrel{~}{\alpha }w)}{(1+2\stackrel{~}{\alpha }w)^2}}`$
(29) $`q(w)`$ $`=`$ $`{\displaystyle \frac{1}{1+2\stackrel{~}{\alpha }w}}{\displaystyle \frac{4\stackrel{~}{\alpha }}{3}}{\displaystyle \frac{w(1+\stackrel{~}{\alpha }w)}{(1+2\stackrel{~}{\alpha }w)^3}}.`$
As previously mentioned, we recall that we assume the conditions
$$n(0)=v(0)=0.$$
We note that in this case eq. (26) gives $`v(0)=0`$ as a consequence of the boundary condition $`n(0)=0`$.
We look for solutions of eqs. (25)-(26) satisfying the following conditions
(30) $`n(0)=0,\underset{w+\mathrm{}}{lim}v(w)=0`$
(31) $`n(w)0\text{for every }w0\text{ and }{\displaystyle _{\text{}}}n(w)𝑑w>0.`$
Since eqs. (25)-(26) are linear and homogeneous, if a solution $`(n(w),v(w))`$, satisfying the above conditions exists, then, for every $`c>0`$, also $`(cn(w),cv(w))`$ is solution.
It is important to note that the conservation of particle number holds for the system (25)-(26). In fact, from eq. (25) it is easy to verify that
(32)
$$_0^+\mathrm{}s(w)\left[an(w+\alpha )+n(w\alpha )\right]\left(s(w+\alpha )+as(w\alpha )\right)n(w)dw=0,$$
for any nonnegative function $`n`$ such that $`n(w)s(w)`$ is integrable and $`n(0)=0`$. This property forces to choose the boundary on $`v`$ in (30), being $`v(0)=0`$.
7. Asymptotic equations
A simple approximate solution of the previous problem can be obtained in the parabolic band case. We introduce a new variable $`\psi `$ defined by
$$n(w)=\sqrt{w}\psi (w).$$
Then, taking into account that now $`\stackrel{~}{\alpha }=0`$, eqs. (25)-(26) become
(33) $`\zeta {\displaystyle \frac{dv}{dw}}=\sqrt{w}\left[a\sqrt{w+\alpha }\psi (w+\alpha )+\sqrt{(w\alpha )_+}\psi ((w\alpha )_+)\right]`$
$`\text{ }\sqrt{w}\left(\sqrt{w+\alpha }+a\sqrt{(w\alpha )_+}\right)\psi (w)`$
(34) $`{\displaystyle \frac{2}{3}}\zeta \sqrt{w^3}{\displaystyle \frac{d\psi }{dw}}=\left(\sqrt{w+\alpha }+a\sqrt{(w\alpha )_+}+\beta \sqrt{w}\right)v(w),`$
where $`(z)_+=\mathrm{max}\{z,0\}`$.
We look for asymptotic equations of system (S0.Ex32)-(34) for large values of the energy $`w`$. To this scope we expand the coefficients of the equations up to the zeroth order (for example $`\sqrt{w+\alpha }\sqrt{w}`$). We obtain a new set of equations
(35) $`\zeta {\displaystyle \frac{dv_A}{dw}}`$ $`=`$ $`w\left\{\left[a\psi _A(w+\alpha )+\psi _A(w\alpha )\right](a+1)\psi _A(w)\right\}`$
(36) $`v_A(w)`$ $`=`$ $`{\displaystyle \frac{2\zeta }{3(1+a+\beta )}}w{\displaystyle \frac{d\psi _A}{dw}},`$
where the subscript $`A`$ indicates the new unknowns. Inserting $`v_A(w)`$ in the first equation and neglecting a coefficient proportional to $`w^1`$, a single equation in $`\psi _A`$ is obtained
(37)
$$\frac{2\zeta ^2}{3(a+1+\beta )}\frac{d^2\psi _A}{dw^2}+a\psi _A(w+\alpha )+\psi _A(w\alpha )(a+1)\psi _A(w)=0.$$
This is a linear difference-differential equation. It is easy to see that $`\mathrm{exp}(\lambda w)`$ is a solution of eq. (37) if and only if $`\lambda `$ satisfies the transcendent equation
(38)
$$\frac{2\zeta ^2}{3(a+1+\beta )}\lambda ^2+ae^{\lambda \alpha }+e^{\lambda \alpha }(a+1)=0.$$
In Appendix A, we prove that this equation admits only two solutions: $`\lambda =0`$ and $`\overline{\lambda }(1,0)`$. We give also a simple approximation for $`\overline{\lambda }`$.
By means of this solution of eq. (37), we get the original quantities
(39)
$$n_A(w)=c\sqrt{w}e^{\overline{\lambda }w},v_A(w)=\frac{2\zeta w}{3(1+a+\beta )}\overline{\lambda }ce^{\overline{\lambda }w},$$
where $`c`$ is a positive constant, in order to satisfy the condition (31). Despite this solution is obtained for large values of $`w`$ it often agrees with the numerical results (which we present in next section) in the whole range of $`w`$.
In the limit case $`\zeta =0`$ (no electric field) $`\overline{\lambda }=1`$ and the solution of (37) gives the correct function $`N`$ obtained from the equilibrium maxwellian distribution function .
For high electric fields approximate solutions are known (see , for example). These are obtained by means of an asymptotic expansion of the original differential-difference equations, which are transformed in simple ordinary differential equations and are explicitly solved. Our solution, obtained by the asymptotic equation, is different with respect to the previous ones, because we do not approximate the difference terms $`f_0(ϵ\pm \mathrm{}\omega )`$ ($`\omega `$ is the constant optical phonon frequency), appearing in the equations, using Taylor formula. In fact the asymptotic equation remains of differential-difference type.
We are interested in calculating the hydrodynamical velocity of electrons
$$v_h:=\frac{{\displaystyle _{u_0}^+\mathrm{}}V_1(u)𝑑u}{{\displaystyle _{u_0}^+\mathrm{}}N(u)𝑑u}=\sqrt{\frac{k_BT_L}{m^{}}}\frac{{\displaystyle _0^+\mathrm{}}v(w)𝑑w}{{\displaystyle _0^+\mathrm{}}n(w)𝑑w}.$$
By a simple calculation, it is possible to verify that, using (39), the following value is obtained
$$\left|v_{A_h}\right|=\sqrt{\frac{k_BT_L}{m^{}}}\frac{4}{3(1+a+\beta )\sqrt{\pi }}\zeta \sqrt{|\overline{\lambda }|}.$$
If we use the approximate value (see Appendix A) of $`\overline{\lambda }`$, we obtain the saturation velocity
$$\underset{\zeta +\mathrm{}}{lim}\left|v_{A_h}\right|=\sqrt{\frac{k_BT_L}{m^{}}}\sqrt{\frac{8\alpha (a1)}{3\pi (1+a+\beta )}}.$$
8. Finite energy model
The difficulties to solve numerically eqs. (25)-(26) are given by the deviating arguments in the right side of equations and by the infinite interval for the variable $`w`$. The first suggests to use a difference scheme to discretize the equations with the constraint that the step $`\mathrm{\Delta }w`$ be such that $`{\displaystyle \frac{\alpha }{\mathrm{\Delta }w}}`$ is an integer. The second difficulty is solved by considering a finite interval for the variable $`w`$. This truncation must be made carefully. To explain this aspect of the problem, we recall that we are looking for a solution which describes the no-runaway phenomena. It is a particular configuration of the gas, where the applied electric field is balanced by the effect of the collisions with the background. In other words, the energy that an electron receives from the external electric field is transferred to the phonons so that a spatial homogeneous and static configuration is maintained.
In order to find this solution also in the case of a finite range for the energy $`w`$, we introduce a cut in the kernels K and $`\text{K}_0`$. Let us choose a positive integer $`\overline{N}`$. Put $`M=\overline{N}\mathrm{}\omega `$, we consider the interval $`[0,M]`$ and the function $`H(u)`$ which coincides with the Heaviside step function $`\theta (u)`$ except for $`u=0`$ where $`H(u)=0`$. Then the new kernels are defined by
$`\text{K}^M(\epsilon (𝐤),\epsilon (\stackrel{~}{𝐤}))`$ $`=`$ $`\text{K}(\epsilon (𝐤),\epsilon (\stackrel{~}{𝐤}))H(M\mathrm{max}\{\epsilon (𝐤),\epsilon (\stackrel{~}{𝐤})\}),`$
$`\text{K}_0^M(\epsilon (𝐤),\epsilon (\stackrel{~}{𝐤}))`$ $`=`$ $`\text{K}_0(\epsilon (𝐤),\epsilon (\stackrel{~}{𝐤}))H(M\mathrm{max}\{\epsilon (𝐤),\epsilon (\stackrel{~}{𝐤})\}).`$
These kernels imply that a particle has zero probability of collision in the following two cases:
* the particle before the collision has an energy less than $`M`$, but after it should have an energy equal or greater than $`M`$;
* the particle before the collision has an energy equal or greater than $`M`$.
Using the new kernels eqs. (25)-(26) are substituted by the following
(40) $`\zeta {\displaystyle \frac{dv}{dw}}=A(n)(w)`$
(41) $`\zeta \left[{\displaystyle \frac{d\text{ }}{dw}}(r(w)n)+q(w)n\right]=B(v)(w),`$
where
$`A(n)(w):=s(w)\left[aH(\alpha \overline{N}w\alpha )n(w+\alpha )+H(\alpha \overline{N}w)n(w\alpha )\right]`$
$`\left[H(\alpha \overline{N}w\alpha )s(w+\alpha )+aH(\alpha \overline{N}w)s(w\alpha )\right]n(w)`$
$`B(v)(w):=\left\{H(\alpha \overline{N}w\alpha )s(w+\alpha )+H(\alpha \overline{N}w)\left[as(w\alpha )+\beta s(w)\right]\right\}v(w).`$
The corresponding boundary conditions becomes
(42) $`n(0)=0,v(\alpha \overline{N})=0`$
(43) $`n(w)0\text{for every }w0\text{ and }{\displaystyle _0^{\alpha \overline{N}}}n(w)𝑑w>0.`$
The boundary condition on $`v`$ is the natural consequence of the condition (30). The choice of the above boundary conditions and the cut are suggested by a simple physical picture. No-runaway phenomena may be considered as the limit for $`t+\mathrm{}`$ of a suitable spatial homogeneous solution of eqs. (11)(21). For time-depending solutions, the new kernels and the condition $`v(\alpha \overline{N})=0`$ imply that a particle with initial energy less than $`M`$ will have, for all time, energy less than $`M`$. For the same reason particles having energy greater than $`M`$ will have, for all time, energy greater than $`M`$. Then, the total number of the particles having energy less than $`M`$ will be constant in time. Different choices can violate the conservation of the particle number for electrons having energy less than $`M`$. In these cases the number of particles having energy less than $`M`$ could tend to zero or to infinity.
9. Numerical solutions
In order to discretize eqs. (40)-(41), we fix the step $`\mathrm{\Delta }w`$ and consider in the interval $`[0,\alpha \overline{N}]`$ the points
$$w_i=i\times \mathrm{\Delta }w(i=0,1,2,\mathrm{}.,\text{N})$$
where $`\text{N}={\displaystyle \frac{\alpha \overline{N}}{\mathrm{\Delta }w}}`$. If we integrate eqs. (40)-(41) in the interval $`[w_i,w_{i+1}]`$ exactly where it is possible and using trapezoidal rule otherwise, we obtain the linear algebraic system
(44) $`{\displaystyle \frac{2\zeta }{\mathrm{\Delta }w}}(v_{i+1}v_i)+A_{i+1}+A_i=0(i=0,1,\mathrm{},\text{N}1),`$
(45) $`v_0=0`$
(46) $`{\displaystyle \frac{2\zeta }{\mathrm{\Delta }w}}(n_{i+1}r_{i+1}n_ir_i)\zeta (q_{i+1}n_{i+1}+q_in_i)B_{i+1}B_i=0`$
$`(i=1,2,\mathrm{},\text{N}1),`$
where, for any function $`\phi (w)`$, $`\phi _i=\phi (w_i)`$. The corresponding boundary conditions are
(47) $`n_0=0,v_\text{N}=0`$
(48) $`n_i0i=0,1,\mathrm{}\text{N}\text{and }{\displaystyle \frac{\mathrm{\Delta }w}{2}}{\displaystyle \underset{i=0}{\overset{\text{N}1}{}}}\left(n_{i+1}+n_i\right)>0.`$
Equations (S0.Ex39) are obtained considering only the interval $`[w_i,w_{i+1}]`$ for $`i1`$, because in the first point $`w_0`$ we have written the exact equation $`v_0=0`$ (i.e. eq. (45)), which derives from eq. (41) with the boundary condition $`n(0)=0`$. It is simple to deduce from (44) the following equations
(49)
$$v_j=v_\text{N}+\frac{\mathrm{\Delta }w}{2\zeta }\underset{i=j}{\overset{\text{N}1}{}}(A_i+A_{i+1})(j=0,1,\mathrm{},\text{N}1).$$
Moreover, as a consequence of the appropriate choice of the cut, the conservation of the particle number holds (see ref. ), being
(50)
$$\underset{i=0}{\overset{\text{N}1}{}}(A_i+A_{i+1})=0.$$
Eq. (50) is the corresponding of eq. (32) introducing the cut in the kernels and performing the integral using the trapezoidal rule.
Now, eq. (49) gives $`v_0=v_\text{N}`$ for $`j=0`$. Due to the boundary condition $`v_\text{N}=0`$, there follows that (45) is a consequence of eqs. (49), or equivalently eq. (49) for $`j=0`$ coincides with eq. (45). For numerical convenience, we prefer to eliminate the case $`j=0`$ in (49). Now, we must add the condition (48). To select an unique solution, we choose the equation $`n_1=1`$. This, with the conditions $`n_i0`$ for $`i=2,3,\mathrm{}\text{N}`$ guarantees (48). This makes the number of equations equal to the number of the unknowns. For the sake of clarity, we rewrite the complete linear system
(51) $`v_\text{N}=0,v_0=0,n_0=0,n_1=1,`$
(52) $`v_j=v_\text{N}+{\displaystyle \frac{\mathrm{\Delta }w}{2\zeta }}{\displaystyle \underset{i=j}{\overset{\text{N}1}{}}}(A_i+A_{i+1})(j=1,2,\mathrm{},\text{N}1)`$
(53) $`{\displaystyle \frac{2\zeta }{\mathrm{\Delta }w}}\left[n_{i+1}r_{i+1}n_ir_i\right]\zeta (q_{i+1}n_{i+1}+q_in_i)B_{i+1}B_i=0`$
$`(i=1,2,\mathrm{},\text{N}1).`$
10. Numerical results and conclusion
We are interested into solve eqs. (40)-(41) in the case of a silicon bulk device. The appropriate values for the parameters are given in the following table.
| $`m^{}=0.32m_e`$ | $`T_L=300`$ K | $`\mathrm{}\omega =0.063`$ eV |
| --- | --- | --- |
| $`\text{K}={\displaystyle \frac{\left(D_tK\right)^2}{8\pi ^2\rho \omega }}`$ | $`D_tK=11.4`$ eV $`\stackrel{}{\text{A}}`$<sup>-1</sup> | $`\rho =2330`$ Kg m<sup>-3</sup> |
| $`\text{K}_0={\displaystyle \frac{k_BT_L}{4\pi ^2\mathrm{}v_0^2\rho }}\mathrm{\Xi }_d^2`$ | $`\mathrm{\Xi }_d=9`$ eV | $`v_0=9040`$ m sec<sup>-1</sup>. |
| $`\alpha _k=0.5eV^1`$ | | |
Here, $`m_e`$ denotes the electron rest mass. Using these parameters, we get $`\alpha 2.437`$ and $`\beta 5.986`$. The values for the electric field are $`|𝐄|=10^3`$, $`10^4`$ and $`10^5`$ Vcm<sup>-1</sup>. Except for low electric field there is a great difference between P (parabolic) and NP (non-parabolic) band. In fig. 1 we plot the velocity versus the electric field, obtained using the asymptotic solutions (39) and values yielded by numerical integration of eqs. (40)-(41) both in P and NP cases. In the first case the value for the saturation velocity obtained is $`1.27\times 10^5m/sec`$. The second model does not show a limit value as in the Monte Carlo simulations . For the largest value of $`|𝐄|`$ the velocity is $`0.95\times 10^5m/sec`$ The experimental value is $`1\times 10^5m/sec`$. In order to compare numerical and asymptotic solutions, we have chosen the condition
(54)
$$_0^{\alpha \overline{N}}n_A(w)𝑑w=_0^{\alpha \overline{N}}n(w)𝑑w$$
where $`n_A(w)`$ indicates the asymptotic solution. The integral in (54) of the function $`n_A`$ is performed analytically and the second numerically.
In such way the constant $`c`$ in (39) is chosen. The meaning of eq. (54) is evident and coherent with the criterion of the choice of the cut. We have made many numerical experiments, varying the values of the electric field, the step $`\mathrm{\Delta }w`$ and the cut $`\overline{N}`$. In fig. 2 the quantity $`N(u)`$, both for P and NP approximation, is shown for a low value of the electric field. Nevertheless a small range of the energy $`04\mathrm{}\omega `$ is sufficient to contain almost all the electrons, the numerical results for P case agree with the asymptotic solution. This is due to the fact that for a small values of the electric field $`\overline{\lambda }1`$. The quantity $`V_1(u)`$ is shown in fig. 3. Here the asymptotic solution is given by the formula
$$v_A(w)=\frac{2\zeta \sqrt{w^3}c\overline{\lambda }e^{\overline{\lambda }w}}{3\left(\sqrt{w+\alpha }+a\sqrt{w\alpha }+\beta \sqrt{w}\right)}$$
This choice is not coherent with the expansion, but it allows a better agreement with the numerical solutions of P case and it puts in evidence a discontinuity in the derivative at the point $`w=\alpha u=\mathrm{}\omega `$, due to the term $`\sqrt{w\alpha }`$. This square root is also present in the complete set of equations (25)-(26). Then, such discontinuity is also expected in the solutions. A deeper analysis could be necessary to prove mathematically this fact. Figs. 4 and 5 shown an intermediate case where the agreement is poor.
When the electric field is high (figs. 6 and 7) very small variations are observable. This agreement was reasonable for large values of $`u`$. In the other cases this is due to the number of particles having low energy which is enough small with respect to that of the whole gas.
The last figures show a quite surprising result in the parabolic case. We have used the same value of the electric field $`E`$ as in fig. 6 and 7. Here, a small value of $`\overline{N}`$ is used. The density $`N`$ agrees with the asymptotic solutions. The quantity $`V_1`$ subject to the boundary condition $`V_1(\overline{N}\alpha )=0`$, follows the asymptotic solutions almost to the end of the interval.
Acknowledgments
We acknowledge partial support from Italian Consiglio Nazionale delle Ricerche (Prog. N. 96.03855.CT01) and TMR project n. ERBFMRCT970157 “Asymptotic Methods in Kinetic Theory”.
Appendix A
The function
$$\eta (\lambda )=\frac{2\zeta ^2}{3(a+1+\beta )}\lambda ^2+ae^{\lambda \alpha }+e^{\lambda \alpha }(a+1)$$
for $`\zeta 0`$ is convex, because $`\eta ^{\prime \prime }(\lambda )0`$ for every real $`\lambda `$.
Since $`\eta (0)=0`$ and $`\eta ^{}(0)=\alpha (a1)>0`$, then there follows that the minimum is negative and is reached at a point $`\lambda _m<0`$.
From $`\eta (1)=2\zeta ^2\left(3(a+1+\beta )\right)^1`$ (recall that $`a=e^\alpha `$) we get $`1<\lambda _m<0`$. Then, there exists only $`\overline{\lambda }0`$ such that $`\eta (\overline{\lambda })=0`$. Moreover $`\overline{\lambda }[1,0[`$. An approximation of this value is obtained expanding $`e^{\pm \lambda \alpha }`$ around $`\lambda =0`$ up the term of second order. A second order algebraic equations follows
$$\left[\frac{2\zeta ^2}{3(a+1+\beta )}+\frac{\alpha ^2}{2}(a+1)\right]\lambda ^2+\alpha (a1)\lambda =0.$$
Then
$$\overline{\lambda }\frac{6(a+1+\beta )\alpha (a1)}{4\zeta ^2+3(a+1+\beta )\alpha ^2(a+1)}.$$
For $`\zeta 1`$ we get
$$\overline{\lambda }\frac{3\alpha (a1)(a+1+\beta )}{2\zeta ^2}.$$
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# Error Avoiding Quantum Codes and Dynamical Stabilization of Grover’s Algorithm
## I Introduction
According to a suggestion of Feynman quantum systems are not only of interest for their own sake but they might also serve for practical purposes. Thus they may be used for simulating other quantum systems which are less convenient to handle or they may be used for solving computational problems more efficiently than by any other classical means. Two well known examples demonstrating this latter point are Shor’s factorization algorithm and Grover’s search algorithm .
Quantum systems which are capable of performing quantum algorithms are called quantum computers. So far several physical systems have been considered as potential candidates for quantum computers, such as trapped ions , nuclear spins of molecules or in the context of cavity quantum electrodynamics atoms interacting with a single mode of the radiation field . To describe the operation of a quantum computer theoretically it is advantageous to refrain from a detailed physical description of the particular quantum system involved. Thus, analogous to the spirit of computer science it is more useful to concentrate on those particular aspects which are essential for the performance of quantum computation. On this abstract level a generic quantum computer consists of $`m`$ distinguishable smaller quantum systems which are frequently chosen as two-level systems with basis states $`|1`$ and $`|0`$, for example. The quantum information which can be stored in one of these two level systems is called a qubit. Thus the state space of a generic quantum computer is spanned by the so called computational basis which consists of the corresponding $`2^m`$ product states $`|b_0=|0\mathrm{}00`$, $`|b_1=|0\mathrm{}01`$, … $`|b_{2^m}=|1\mathrm{}11`$.
A typical quantum computation proceeds in several steps. Firstly, the quantum computer is prepared in an initial state. Secondly, a certain sequence of unitary transformations is performed which are called quantum gates and which usually entangle the $`m`$ qubits. Thirdly, the final result is measured. Typically the solution of a particular computational problem is obtained with a certain probability only. A general quantum algorithm takes advantage of an essential feature of quantum theory, namely the interference between probability amplitudes and the fact that the dimensionality $`D`$ of the state space of $`m`$ distinguishable qubits increases exponentially with the number of qubits, i.e. $`D=2^m`$. One of the best known quantum algorithms are the already mentioned Shor algorithm and Grover’s search algorithm . In this latter algorithm a particular sequence of quantum gates enables one to find a specific item out of an unsorted database much faster than with any other known classical mean. This quantum algorithm was already realized experimentally for a small number of qubits .
One of the main practical problems one has to overcome in the implementation of quantum algorithms are non-ideal performances of the quantum gates involved or random environmental influences which both tend to affect the relevant quantum coherence. To protect quantum computation against such errors two major strategies have been proposed recently, namely quantum error correction and error avoiding quantum codes . Quantum error correction rests on the assumption that nothing is known about the physical origin of the errors affecting the quantum computation. Its methods may be considered as an extension of classical error correction techniques to the quantum domain. The approach of the error avoiding quantum codes is different. They rest on the assumption that the physical origin of the errors which affect the quantum computation is known. The main idea of this latter approach is to encode the logical information in one of those subspaces of the relevant Hilbert space which is not affected by the physical interactions responsible for the occurance of errors . Both theoretical approaches to error correction rest on the concept of redundancy which is also fundamental for classical error correcting codes. Provided the physical origin of the errors affecting a quantum computation is known it is expected that error avoiding codes offer more effective means for stabilizing quantum algorithms. This expectation is based on two facts. Firstly, there is no need for control measurements which are an essential ingredient for any error correcting code. Secondly, usually a smaller number of physical qubits is needed for the representation of a given number of logical qubits.
In the subsequent discussion it is demonstrated that this is indeed the case. By considering Grover’s quantum search algorithm it is shown that non-ideal perturbations may be corrected dynamically in an efficient way with the help of an appropriate error avoiding quantum code. As a particular example, we discuss coherent errors which may arise from systematic detunings of the physical qubits of the quantum computer from the frequency of the light pulses which realize the required quantum gates. It is shown that the corresponding error avoiding quantum code with the lowest degree of redundancy is more efficient in encoding quantum information than the corresponding optimal error correcting code which saturates the quantum Hamming bound. The proposed error avoiding quantum code consists of states only which are factorizable in the computational basis. In this respect it differs significantly from the recently proposed error avoiding code of Ref., for example, which also involves entangled states. Such factorizable codes may offer practical advantages as far as the implementation of quantum gates in error avoiding subspaces is concerned.
The article is organized as follows: In Sec. II basic facts about Grover’s quantum search algorithm are summarized. It is demonstrated that for large databases the dynamics of this quantum algorithm can be described by a two-level Hamiltonian which implies Rabi oscillations between the initial state and the search state. In Sec. III general ideas underlying the construction of error avoiding quantum codes are discussed. An efficient error avoiding quantum code is presented which is capable of stabilizing Grover’s algorithm against a particular class of coherent errors. The redundancy of this code is discussed and compared with the one resulting from error correcting codes which saturate the quantum Hamming bound. Numerical examples demonstrating the stabilizing capabilities of this error avoiding quantum code are presented in Sec. IV.
## II Grover’s quantum search algorithm
Consider an unsorted database with $`N`$ items and a certain item $`x_0`$ you are searching for. As a particular example you can imagine a telephone directory with $`N`$ entries and a particular telephone number $`x_0`$ you are looking for. Furthermore, assume that you are given a black box, i.e. a so called oracle, which can decide whether an item is $`x_0`$ or not. Thus, in mathematical terms you are given a Boolean function
$$f(x)=\delta _{x,x_0}=\{\begin{array}{cc}1& x=x_0\\ 0& xx_0\end{array}$$
(1)
with $`\delta _{a,b}`$ denoting the Kronecker delta function. Usually the elements $`x`$ of the database are assumed to be described by the $`N`$ integers between zero and $`N1`$. Assuming that each application of the oracle requires one elementary step a classical random search process will require $`N1`$ steps in the worst case and one step in the best possible case. Thus, on the average a classical algorithm will need $`N/2`$ steps to find the searched item $`x_0`$. It has been shown by Grover that with the help of his quantum search algorithm this task can be performed in $`O(\sqrt{N})`$ steps with a probability arbitrarily close to unity. The basic idea of this quantum algorithm is to rotate the initial state of the quantum computation in the direction of the searched state $`|x_0`$ by a sequence of unitary quantum versions of the oracle. It will become apparent from the subsequent discussion that apart from Hadamard transformations the dynamics of this rotation are analogous to a Rabi oscillation between the initially prepared state and the searched state $`|x_0`$. It has been shown by Zalka that Grover’s quantum search algorithm is optimal.
### A Characteristic gate sequence of Grover’s search algorithm
In Grover’s quantum search algorithm every element of the database is represented by a state of the computational basis of the quantum computer. Thus a database which is represented by $`m`$ qubits has $`N=2^m`$ distinguishable elements. The state $`|\mathrm{0..0110..0}`$ of the computational basis, for example, corresponds to the element $`\mathrm{0..0110..0}`$ of the database in binary notation. The quantum oracle $`𝒰_f`$ is determined completely by the Boolean function of Eq.(1) and is represented by a quantum gate, i.e. by the unitary and hermitian transformation
$$𝒰_f:|x,a|x,f(x)a.$$
(2)
Thereby $`|x`$ is an arbitrary element of the computational basis and $`|a`$ is the state of an additional ancilla qubit which is discarded later. The symbol $``$ denotes addition modulo 2. This unitary form of the oracle depends on the Boolean function $`f(x)`$. As far as complexity estimates are concerned it is assumed that this unitary transformation requires one elementary step. This assumption is analogous to the complexity estimate of the corresponding classical version of this search problem.
For the subsequent discussion it is important to note that the elementary rotations in the direction of the searched quantum state $`|x_0`$ which are the key ingredient in Grover’s algorithm can be performed with the help of this unitary oracle. Thus such a rotation can be performed without explicit knowledge of the state $`|x_0`$. Its implicit knowledge through the values of the Boolean function $`f(x)`$ is already sufficient. For large values of $`N`$ it turns out that the number of elementary rotations needed to prepare state $`|x_0`$ is $`O(\sqrt{N})`$. To implement such an elementary rotation from the initial state $`|s=|0\mathrm{}0`$, for example, towards the final state $`|x_0`$ two different types of quantum gates are needed, namely Hadamard gates and controlled phase inversions.
A Hadamard gate is a unitary one-qubit operation. It produces an equal weighted superposition of the two basis states according to the rule
$`|0`$ $``$ $`{\displaystyle \frac{1}{\sqrt{2}}}(|0+|1),`$ (3)
$`|1`$ $``$ $`{\displaystyle \frac{1}{\sqrt{2}}}(|0|1)`$ (4)
or in matrix notation
$`H^{(2)}={\displaystyle \frac{1}{\sqrt{2}}}\left(\begin{array}{cc}1& 1\\ 1& 1\end{array}\right).`$
A $`m`$-qubit Hadamard gate $`H^{(2^m)}`$ is defined by the $`m`$-fold tensor product, i.e. $`H^{(2^m)}=H^{(2)}\mathrm{}H^{(2)}`$. Thus, for two qubits, for example, $`H^{(2^2)}`$ is represented by the matrix
$$H^{(2^2)}=\frac{1}{2}\left(\begin{array}{cccccccc}1& 1& 1& 1& & & & \\ 1& 1& 1& 1& & & & \\ 1& 1& 1& 1& & & & \\ 1& 1& 1& 1& & & & \end{array}\right).$$
(5)
The Hadamard transformation is hermitian and unitary. An arbitrary matrix element $`H_{i,j}^{(2^m)}`$ of a Hadamard transformation may be written in the general form
$$H_{i,j}^{(2^m)}=\frac{1}{\sqrt{2^m}}(1)^{ij}.$$
(6)
Thereby $`i`$ and $`j`$ denote binary numbers and the multiplication $``$ is bitwise modulo 2, i.e. for $`i=1`$, $`j=3`$ and $`m=2`$, one obtains $`H_{1,3}^{(4)}=(1/2)(1)^{(0111)}=(1/2)(1)^{(01+11)}=1/2`$. It has been shown by Grover that this Hadamard transformation can be replaced by any other unitary one-qubit operation.
The remaining quantum gates needed for the implementation of the necessary rotation are controlled phase inversions with respect to the initial and searched states $`|s=|0\mathrm{}0`$ and $`|x_0`$. A controlled phase inversion with respect to a state $`|x`$ changes the phase of this particular state by an amount of $`\pi `$ and leaves all other states unchanged. Thus the phase inversion $`I_s`$ with respect to the initial state $`|s`$ is defined by
$`I_s|s`$ $`=`$ $`|s,`$ (7)
$`I_s|x`$ $`=`$ $`|x(xs).`$ (8)
For two qubits, for example, its matrix representation is given by
$$I_s=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& 1\end{array}\right).$$
(9)
The controlled phase inversion $`I_{x_0}`$ with respect to the searched state $`|x_0`$ is defined in an analogous way. As state $`|x_0`$ is not known explicitly but only implicitly through the property $`f(x_0)=1`$ this transformation has to be performed with the help of the quantum oracle. This task can be achieved by preparing the ancilla of the oracle of Eq.(2) in state $`|a_0=1/\sqrt{2}(|0|1)`$. As a consequence one obtains the required properties for the phase inversion $`I_{x_0}`$, namely
$`|x,f(x)a_0|x,0a_0`$ (11)
$`=1/\sqrt{2}(|x,0|x,1)=|x,a_0\mathrm{for}xx_0,`$
$`|x,f(x)a_0|x,1a_0`$ (13)
$`=1/\sqrt{2}(|x,1|x,0)=|x,a_0\mathrm{for}x=x_0.`$
One should bear in mind that this controlled phase inversion can be performed with the help of the quantum oracle of Eq.(2) only without explicit knowledge of state $`|x_0`$.
Grover’s algorithm starts by preparing all $`m`$ qubits of the quantum computer in state $`|s=|0\mathrm{}0`$. An elementary rotation in the direction of the searched state $`|x_0`$ with the property $`f(x_0)=1`$ is achieved by the gate sequence
$$Q=I_sH^{(2^m)}I_{x_0}H^{(2^m)}.$$
(14)
In order to rotate the initial state $`|s`$ into state $`|x_0`$ one has to perform a sequence of $`n`$ such rotations and a final Hadamard transformation at the end, i.e.
$$|f=HQ^n|s.$$
(15)
The effect of the elementary rotation $`Q`$ is demonstrated in Fig.3 for the case of three qubits, i.e. $`m=3`$. The first Hadamard transformation $`H^{(2^3)}`$ prepares an equally weighted state. The subsequent quantum gate $`I_{x_0}`$ inverts the amplitude of the searched state $`|x_0=|111`$. Together with the subsequent Hadamard transformation and the phase inversion $`I_s`$ this gate sequence $`Q`$ amplifies the probability amplitude of the searched state $`|111`$. In this particular case an additional Hadamard transformation finally prepares the quantum computer in the searched state $`|111`$ with a probability of 0.88.
In order to determine the dependence of the ideal number of repetitions $`n`$ on the number of qubits $`m`$ it is convenient to analyze the repeated application of the gate sequence $`Q`$ according to Eq.(15) in terms of the two states $`|s`$ and $`|v=H^{(2^m)}|x_0`$ whose overlap is given by $`ϵ=s|v=s|H^{(2^m)}|x_0=2^{m/2}`$ for $`m`$ qubits. It is straightforward to show that the unitary gate sequence $`Q`$ preserves the subspace spanned by these two states, i.e.
$$Q\left(\begin{array}{c}|s\\ |v\end{array}\right)=\left(\begin{array}{cc}14ϵ^2& 2ϵ\\ 2ϵ& 1\end{array}\right)\left(\begin{array}{c}|s\\ |v\end{array}\right).$$
(16)
Thus $`Q`$ acts like a rotation in the plane spanned by states $`|s`$ and $`|v`$. The angle of rotation is given by $`\phi =\mathrm{arcsin}(2ϵ\sqrt{1ϵ^2})`$.
After j iterations the amplitude of state $`|v`$ is given by
$$\mathrm{sin}\left[(2j+1)ϵ\right].$$
(17)
Therefore, the optimal number $`n`$ of repetitions of the gate sequence $`Q`$ is approximately given by
$$n=\frac{\pi }{4\mathrm{arcsin}\left(2^{m/2}\right)}\frac{1}{2}\frac{\pi }{4}\sqrt{2^m}(2^m1).$$
(18)
### B Hamiltonian representation of Grover’s algorithm
If the database contains many elements, i.e. $`Nϵ^21`$, the repeated application of the elementary rotation which is essential for Grover’s search algorithm can be described by a Hamiltonian quantum dynamics. The elementary rotation $`Q`$ can be approximated by the relation
$$Q=\mathrm{𝟏}\tau i/\mathrm{}𝐇_G(ϵ)+O(ϵ^2)$$
(19)
which involves the Hamiltonian
$$𝐇_G=2iϵ\frac{\mathrm{}}{\tau }\left(|vs||sv|\right).$$
(20)
The elementary time $`\tau `$ might be interpreted as the physical time required for performing the elementary rotation $`Q`$. The Hamiltonian of Eq.(20) describes the dynamics of a quantum mechanical two level system whose degenerate energy levels $`|s`$ and $`|v`$ are coupled by a time-independent perturbation. In lowest order of $`ϵ`$ these degenerate energy levels are orthogonal. The resulting oscillations between these coupled energy levels are characterized by the Rabi frequency $`\mathrm{\Omega }=2s|v/\tau `$. Correspondingly, the repeated application of the elementary rotation $`Q`$ can be determined with the help of Trotter’s product formula , namely
$`Q^n`$ $`=`$ $`(I_sH^{(2^m)}I_{x_0}H^{(2^m)})^n=`$ (22)
$`\mathrm{exp}\left({\displaystyle \frac{i}{\mathrm{}}}𝐇_G\tau n\right)+O(ϵ^2n).`$
Thus, in the framework of this Hamiltonian description applying the elementary rotation $`Q`$ $`n`$ times is equivalent to a time evolution of the effective two-level quantum system over a time interval of magnitude $`n\tau `$. This Hamiltonian description demonstrates that the physics behind Grover’s quantum search algorithm is the same as the physics governing the Rabi oscillations between degenerate or resonantly coupled energy eigenstates. As the errors entering Eq.(22) are of order $`O(ϵ^2n)`$ this Hamiltonian description is applicable only as long as $`ϵ^2nn/2^m1`$. Thus for a given size of the database it is valid only as long as the number of iterations is sufficiently small, i.e. $`n2^m`$. However, as Grover’s search algorithm needs approximately $`(\pi \sqrt{2^m}/4)`$ steps to find the searched item the main condition which restricts the validity of this Hamiltonian description is a large size of the database, i.e. $`ϵ^21/N1`$.
### C An example of coherent errors
So far we have been concentrating on the ideal dynamics of Grover’s quantum search algorithm. However, in practical applications it is very difficult to realize this search algorithm in an ideal way. Usually the ideal dynamics are affected by numerous perturbations. Physically one may distinguish two different kinds of errors, namely incoherent and coherent ones. Typically incoherent perturbations originate from a coupling of the physical qubits of a quantum computer to an uncontrollable environment. As a consequence the resulting errors are of a stochastic nature. Coherent errors may arise from non-ideal quantum gates which lead to a unitary but non-ideal time evolution of the quantum algorithm. A simple example of this latter type of errors are systematic detunings from resonance of the light pulses with which the required quantum gates are realized on the physical qubits. In the Hamiltonian formulation of Grover’s algorithm such systematic detunings may be described by a perturbing Hamiltonian of the form
$$𝐇_d=\underset{i=1}{\overset{m}{}}\mathrm{}\omega _i\sigma _z^{(i)}.$$
(23)
In Eq.(23) it has been assumed that Grover’s quantum algorithm is realized by $`m`$ qubits and that the $`i`$-th qubit is detuned with respect to the ideal transition frequency by an amount $`\omega _i`$. The Pauli spin-operator of the $`i`$-th qubit is denoted $`\sigma _z^{(i)}`$. In the presence of these systematic detunings and for a large number of qubits the dynamics of Grover’s algorithm are described by the Hamiltonians of Eqs.(20) and (23).
In order to obtain insight into the influence of this type of coherent errors the performance of Grover’s algorithm under repeated applications of the elementary rotation $`Q`$ is depicted in Fig. 5. The dynamics of the ideal Grover algorithm are depicted by the dashed line for the case of three qubits, i.e. $`m=3`$. The Rabi oscillations with frequency $`\mathrm{\Omega }=2v|s/\tau `$ are clearly visible. The solid line shows the probability of observing the quantum computer in state $`|x_0`$ in a case in which all the qubits are detuned from their ideal resonance frequency. One notices the deviations from the ideal behaviour. Due to the coherent nature of the errors the time evolution of the non-ideal algorithm exhibits revival phenomena .
## III Error avoiding quantum codes
In general there are two different strategies for correcting errors in quantum information processing. If nothing is known about the physical origin of the errors affecting a qubit one can use general quantum error correcting schemes. They may be viewed as generalizations of classical error correction techniques to the quantum domain . Typically they involve a suitably chosen quantum code and a sequence of quantum measurements. This code has to map all possible states which may result from arbitrary environmental influences onto orthogonal states. According to basic postulates of quantum theory these orthogonal quantum states can be distinguished and based on the result of a control measurement one may restore the original quantum state. So far these general techniques have been applied mainly to the stabilization of static quantum memories . It is still an open question whether these general methods are also useful and efficient for the dynamical stabilization of quantum algorithms.
The second possible error correction strategy which seems to be well adopted also for stabilizing quantum algorithms is based on error avoiding quantum codes . However, these latter methods are applicable only, if the physical origin of the errors is known. The main idea is to encode the quantum information in those subspaces of the Hilbert space which are not affected by the errors. This aim is achieved by restricting oneself to degenerate eigenspaces of the relevant error operators. Thus, in the special case of a single error operator, say $`𝐄`$, the basis states $`\{|\psi _i\}`$ of such an error free subspace have to fulfill the relation
$$𝐄|\psi _i=c|\psi _i.$$
(24)
In the above mentioned example of coherent errors which may affect Grover’s algorithm this error operator is given by the Hamiltonian of Eq.(23), i.e. $`𝐄=𝐇_d`$. It is crucial for the success of an error avoiding code that the eigenvalue $`c`$ of Eq.(24) does not depend on the states belonging to the error free subspace. This implies that all possible elements of the error free subspace of the general form $`_i\alpha _i|\psi _i`$ are affected by the error operator in the same way, i.e.
$$𝐄(\underset{i}{}\alpha _i|\psi _i)=c(\underset{i}{}\alpha _i|\psi _i).$$
(25)
It is apparent that a non-trivial error avoiding code is possible only, if the eigenspace of the error operator $`𝐄`$ is degenerate.
### A An error avoiding quantum code stabilizing coherent errors
As an example for an error avoiding quantum code let us consider the case of coherent errors which may affect Grover’s quantum algorithm and which can be characterized by the Hamiltonian $`𝐇_d`$ of Eq.(23). In the simple case of equal detunings, i.e. $`\omega _1=\mathrm{}=\omega _m\omega `$, the error operator $`𝐄`$ reduces to the form
$$𝐇_e=\mathrm{}\omega \underset{i=1}{\overset{m}{}}\sigma _z^{(i)}.$$
(26)
It is easy to find highly degenerate error free subspaces of this error operator. All states with a fixed number of ones and zeroes constitute a degenerate eigenspace of $`𝐇_e`$. For an even number of qubits it is possible to find an error avoiding subspace with eigenvalue $`c=0`$ so that
$$(𝐇_G+𝐇_e)|\psi =𝐇_G|\psi $$
(27)
for all elements $`|\psi `$ of this subspace. For this purpose one is looking for quantum states with zero total spin. For four qubits, for example, this subspace is defined by the basis vectors $`|0011,|0101,|0110,|1001,|1010,|1100`$ and involves all states with the same number of zeros and ones. Four of these states may be used as a basis for the state space of two logical qubits. For these eigenstates the error Hamiltonian $`𝐇_e`$ maps onto zero, e.g.
$`𝐇_e|0011=\mathrm{}\omega {\displaystyle \underset{i=1}{\overset{m=4}{}}}\sigma _z^{(i)}|0011=\mathrm{}\omega (1+111)|0011=0.`$
This particular error avoiding code works ideal for equal detunings of all qubits from resonance. It is formed by quantum states which factorize in the computational basis. So it is expected that in this error free subspace the encoding of quantum information and the implementation of quantum gates is considerably easier than in cases in which the error avoiding codes involve entangled quantum states.
### B Implementation of quantum gates in an error free subspace
To realize a quantum algorithm in an error free subspace one has to implement the necessary quantum gates in such a way that they do not mix the error free subspace with its orthogonal complement . Consider two logical qubits, for example, which are encoded by four physical qubits. For this purpose one may choose the states $`|0011,|0101,|0110,|1001`$ which have been mentioned in the previous subsection. This error avoiding code works ideal for stabilizing Grover’s algorithm with respect to the error operator $`𝐇_e`$ of Eq.(26) provided it is possible to realize the required unitary transformations, namely Hadamard transformations and the controlled phase inversions.
Consider as an example a Hadamard transformation which acts in a two dimensional error avoiding subspace of this kind. Thus it is assumed that the two basis states of this error avoiding code are given by $`|01`$ and $`|10`$ and that they involve two physical qubits. Thus, we are looking for a transformation which performs the mappings
$`|01`$ $``$ $`1/\sqrt{2}(|01+|10),`$ (28)
$`|10`$ $``$ $`1/\sqrt{2}(|01|10)`$ (29)
and which does not mix the subspace spanned by $`|01`$ and $`|10`$ with the orthogonal space spanned by the basis states $`|00`$ and $`|11`$. In matrix notation we are looking for a unitary matrix of the form
$$\left(\begin{array}{cccc}\hfill & \hfill 0& \hfill 0& \hfill \\ \hfill 0& \hfill 1& \hfill 1& \hfill 0\\ \hfill 0& \hfill 1& \hfill 1& \hfill 0\\ \hfill & \hfill 0& \hfill 0& \hfill \end{array}\right)$$
(30)
with $``$ denoting arbitrary entries which ensure unitarity. Such a transformation can be achieved by the gate sequence $`CNOT_{21}(\mathrm{𝟏}\stackrel{~}{H}^{(2)})CNOT_{21}`$ with $`\stackrel{~}{H}^{(2)}=i\sigma _yH^{(2)}`$. Thereby $`CNOT_{21}`$ is a controlled-not operation with the first qubit as the target and the second qubit as the control qubit and $`\sigma _y`$ is the Pauli Matrix. Thus in matrix notation this relation yields
$`\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 0& 0& 1\\ 0& 0& 1& 0\\ 0& 1& 0& 0\end{array}\right)\left(\begin{array}{cccc}\hfill 1& \hfill 1& \hfill 0& \hfill 0\\ \hfill 1& \hfill 1& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 1& \hfill 1\\ \hfill 0& \hfill 0& \hfill 1& \hfill 1\end{array}\right)\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 0& 0& 1\\ 0& 0& 1& 0\\ 0& 1& 0& 0\end{array}\right)=`$ (43)
$`\left(\begin{array}{cccc}\hfill 1& \hfill 0& \hfill 0& \hfill 1\\ \hfill 0& \hfill 1& \hfill 1& \hfill 0\\ \hfill 0& \hfill 1& \hfill 1& \hfill 0\\ \hfill 1& \hfill 0& \hfill 0& \hfill 1\end{array}\right).`$ (48)
Obviously the final result does not mix the error avoiding subspace with its orthogonal complement. But in the intermediate steps such a mixing takes place. However, for practical purposes it is enough to ensure that the time spent by the quantum computer in the orthogonal complement of the error avoiding subspace is sufficiently small so that the resulting errors can be neglected for all practical purposes. Under these circumstances it is expected that the implementation of quantum algorithms in error avoiding subspaces is a powerful tool for stabilizing quantum codes.
### C Code size of error avoiding quantum codes
In order to estimate the redundancy which has to be introduced for stabilizing a quantum algorithm by an error avoiding quantum code let us consider the particular example of Sec. IIIA in more detail. It has been argued that in the case of coherent errors which can be characterized by the Hamiltonian of Eq.(26) an error avoiding quantum code can be constructed from basis states with equal numbers of ones and zeroes. In order to minimize the redundancy it is desirable to maximize the dimension of the resulting error avoiding subspace. If one starts with $`m`$ physical qubits the dimension $`D(m,q)`$ of the corresponding error avoiding subspace with $`q`$ qubits in state $`|1`$ and $`(mq)`$ qubits in state $`|0`$, for example, is given by
$$D(m,q)=\left(\begin{array}{c}m\\ q\end{array}\right)\frac{m!}{q!(mq)!}.$$
(49)
From elementary properties of binomial coefficients it is clear that $`D(m,q)`$ is maximum for $`q=m/2`$. Thus for an even number of qubits $`m`$ the largest possible dimension of the resulting error avoiding subspace is given by
$$D(m,m/2)=\frac{m!}{[(m/2)!]^2}2^m\sqrt{\frac{2}{m\pi }}(m1).$$
(50)
Thus, in this case it is possible to encode
$`l`$ $`=`$ $`\mathrm{log}_2D(m,m/2)`$ (52)
$`m{\displaystyle \frac{\mathrm{log}_2m}{2}}+\mathrm{log}_2\sqrt{2/\pi }(m1)`$
logical qubits with $`m`$ physical ones. It is instructive to compare the redundancy of this error avoiding code as described by Eq.(52) with the ones resulting from general error correcting quantum codes which saturate the quantum Hamming bound . If one wants to correct arbitrary errors of maximum length $`t`$ with a general error correcting quantum code the number of physical and logical qubits $`m`$ and $`l`$ have to fulfill the so called quantum Hamming bound, i.e.
$$2^l\underset{r=0}{\overset{t}{}}3^r\left(\begin{array}{c}m\\ r\end{array}\right)2^m.$$
(53)
This inequality reflects the fact that in a general error correcting quantum code the action of different error operators onto any of the logical qubits must lead to orthogonal quantum states. The dimension of the resulting Hilbert space as described by the left hand side of the inequality (53) has to be smaller than the dimension of the Hilbert space of all physical qubits. Thus the number of logical qubits obtainable by a general error correcting code which is capable of correcting all possible errors of maximum length one, i.e. $`t=1`$, cannot be larger than
$$l_>=m\mathrm{log}_2(3m+1).$$
(54)
Comparing Eq.(52) with Eq.(54) one realizes that the redundancy of this particular error avoiding quantum code is smaller than the one resulting from saturating the Hamming bound with a general error correcting code capable of correcting errors of maximum length one, i.e.
$$ll_>\frac{1}{2}\mathrm{log}_2m+\mathrm{log}_2\frac{3\sqrt{2}}{\sqrt{\pi }}>0(m1).$$
(55)
However, this reduction of redundancy is based on the fact that the error avoiding code obeying Eq.(52) can stabilize errors only which are described by the Hamiltonian of Eq.(26). Usually more general errors cannot be corrected with this code.
## IV Numerical examples
In the previous section we have developed an error avoiding quantum code which is capable of correcting coherent errors. These errors were assumed to be caused by systematic detunings of the physical qubits of the quantum computer from the frequency of the laser pulses implementing the action of the quantum gates. This error avoiding quantum code works perfect provided all physical qubits are detuned from the frequency of these laser pulses by the same amount. However, in realistic situations this case is scarcely realized. For the realistic assumption of unequal detunings in general the eigenstates of $`𝐇_d`$ are non-degenerate so that it is not possible to construct a perfect error avoiding quantum code. Therefore the practical question arises whether the presented error avoiding quantum code of Sec. III is still useful for stabilizing quantum algorithms against arbitrary systematic detunings.
The dynamics of Grover’s algorithm in the presence of arbitrary detunings are depicted in Fig. 7. The dashed line represents the ideal dynamics in the absence of detunings for the case of 6 qubits as evaluated from the Hamiltonian of Eq.(20). The characteristic Rabi oscillations are clearly apparent. The corresponding dynamics for 8 qubits in the presence of arbitrarily chosen detunings are depicted by the dotted line in Fig.7. It is apparent that in this case a quantum search for state $`|x_0`$ is not successful at all. However, as apparent from the solid line of Fig. 7 encoding the quantum information by the error avoiding code of Sec.III improves the performance considerably. Despite the fact that this error avoiding code has not been designed for these detunings it almost succeeds in finding the searched quantum state $`|x_0`$ after a number of iterations which is close to the ideal case (compare with Eq.(18)).
In order to obtain more insight into the stabilizing properties of this error avoiding code let us investigate the probability of success in the presence of arbitrary detunings in more detail. For this purpose we consider 8 physical qubits whose detunings $`\omega _i`$ are distributed randomly according to a normal distribution. According to Fig.6 these 8 physical qubits are capable of encoding 6 logical qubits. In Fig. 8 the average value of the maximum probability of finding the quantum computer in the searched state $`|x_0`$ is depicted for various values of the variance of the randomly chosen detunings. The lower sequence of dots (stars) refers to Grover’s algorithms without error avoiding encoding and the upper sequence of points (diamonds) refers to error avoiding encoding according to Sec. III. It is apparent that error avoiding encoding is very successful as long as the differences between the detunings of the qubits is sufficiently small. Only in extreme cases in which these differences become comparable to the typical magnitudes of the detunings this type this error avoiding code is no longer capable of stabilizing Grover’s algorithm in a satisfactory way.
## V Summary and Conclusions
It has been demonstrated that error avoiding quantum codes may offer efficient methods for stabilizing quantum codes dynamically against those types of errors whose physical origin is known. As a particular example we discussed the stabilization of Grover’s quantum search algorithm against coherent errors which may arise from systematic detunings of the physical qubits from the frequency of the light pulses implementing the quantum gates. Though originally the error avoiding quantum code has been constructed for the special case of equal detunings of all the qubits it has been shown that it is also capable of stabilizing this quantum algorithm in other cases to a satisfactory degree. The proposed error avoiding quantum code consists of quantum states only which are factorizable in the computational basis. This may offer advantages as far as the implementation of the necessary quantum gates in this error free subspace is concerned. Furthermore, this quantum code has also other noteworthy properties, such as a redundancy which is lower than the one of an optimal error correcting quantum code saturating the quantum Hamming bound. Though the stabilizing ability of error avoiding quantum codes has been demonstrated for one particular quantum code and one particular class of coherent errors only it is expected that similar capabilities are also found in more general cases which may also involve incoherent errors.
This work is supported by the DFG within the SPP ‘Quanteninformationsverarbeitung’. Stimulating discussions with Thomas Beth, Markus Grassl and Dominik Janzing are acknowledged.
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# GRAVITY-MODES IN ZZ CETI STARS IV. AMPLITUDE SATURATION BY PARAMETRIC INSTABILITY
## 1. INTRODUCTION
Within an instability strip of width $`\mathrm{\Delta }T_{\mathrm{eff}}10^3\mathrm{K}`$ centered at $`T_{\mathrm{eff}}1.2\times 10^3\mathrm{K}`$, hydrogen white dwarfs exhibit multiple excited gravity modes with $`10^2P10^3\mathrm{s}`$. Convective driving, originally proposed by Brickhill (nonad-brick90 (1990), nonad-brick91 (1991)), is the overstability mechanism (Goldreich & Wu nl-paperI (1998), hereafter Paper I). Individual modes maintain small amplitudes; typical fractional flux variations range from a few mma to a few tens of mma.<sup>1</sup><sup>1</sup>1$`1`$ mma of light variation is approximately $`0.1\%`$ fractional change in flux.
The nonlinear mechanism responsible for saturating mode amplitudes has not previously been identified. We demonstrate that parametric resonance between an overstable parent g-mode and a pair of lower frequency damped daughter g-modes sets an upper envelope to the parent modes’ amplitudes.<sup>2</sup><sup>2</sup>2The g-mode dispersion relation allows plenty of good resonances. Moreover, the envelope we calculate reproduces the broad trends found from observational determinations of mode amplitudes in ZZ Ceti stars. Our investigation follows pioneering work by Dziembowski & Krolikowska (nl-dziem85 (1985)) on overstable acoustic modes in $`\delta `$-Scuti stars. They showed that parametric resonance with damped daughter g-modes saturates the growth of the overstable p-modes at approximately their observed amplitudes.
This paper is comprised of the following parts. In §2 we introduce parametric instability for a pair of damped daughter modes resonantly coupled to an overstable parent mode. We evaluate the parent mode’s threshold amplitude and describe the evolution of the instability to finite amplitude. §3 is devoted to the choice of optimal daughter pairs. We discuss relevant properties of 3-mode coupling coefficients, and the constraints imposed by frequency resonance relations and angular selection rules. Evaluation of the upper envelope for parent mode amplitudes set by parametric resonance is the subject of §4. Numerical results are interpreted in terms of analytic scaling relations and compared to observations. §5 contains a discussion of a variety of issues leftover from this investigation. Detailed derivations are relegated to a series of Appendices.
The stellar models used in this investigation were provided by Bradley (nl-bradley96 (1996)). Their essential characteristics are $`M_{}=0.6M_{}`$, $`\mathrm{log}(g/\mathrm{cm}\mathrm{s}^2)=8.0`$, hydrogen layer mass $`1.5\times 10^4M_{}`$, and helium layer mass $`1.5\times 10^2M_{}`$.
## 2. PARAMETRIC INSTABILITY
In this section we introduce parametric instability (Landau & Lifshitz nl-landau76 (1976)). We present the threshold criterion for the instability, and discuss relevant aspects of the subsequent evolution. Depending upon the parameters, the modes either approach a stable steady-state or develop limit cycles. We describe the energies attained by the parent and daughter modes in either case. Many of the results in this section were obtained earlier by Dziembowski (nl-dziem82 (1982)).
### 2.1. Instability Threshold
Parametric instability in the context of our investigation is a special form of resonant 3-mode coupling. It refers to the destabilization of a pair of damped daughter modes by an overstable parent mode. The frequencies of the three modes satisfy the approximate resonance condition $`\omega _p\omega _{d_1}+\omega _{d_2}`$, where the subscripts $`p`$ and $`d_1,d_2`$ denote parent and two daughter modes, respectively.
Equations governing the temporal evolution of mode amplitudes are most conveniently derived from an action principle (Newcomb nl-newcomb62 (1962), Kumar & Goldreich nl-kumar89 (1989)). The amplitude equations take the form
$`{\displaystyle \frac{dA_p}{dt}}`$ $`=`$ $`+{\displaystyle \frac{\gamma _p}{2}}A_pi\omega _pA_p+i{\displaystyle \frac{3}{\sqrt{2}}}\omega _p\kappa A_{d_1}A_{d_2},`$ (1)
$`{\displaystyle \frac{dA_{d_1}}{dt}}`$ $`=`$ $`{\displaystyle \frac{\gamma _{d_1}}{2}}A_2i\omega _{d_1}A_{d_1}+i{\displaystyle \frac{3}{\sqrt{2}}}\omega _{d_1}\kappa A_pA_{d_2}^{},`$ (2)
$`{\displaystyle \frac{dA_{d_2}}{dt}}`$ $`=`$ $`{\displaystyle \frac{\gamma _{d_2}}{2}}A_{d_2}i\omega _{d_2}A_{d_2}+i{\displaystyle \frac{3}{\sqrt{2}}}\omega _{d_2}\kappa A_pA_{d_1}^{}.`$ (3)
Here, $`A_j`$ is the complex amplitude of mode $`j`$; it is related to the mode energy $`E_j`$ by $`|A_j|^2=E_j`$. The $`\gamma _j`$ ($`>0`$) denote linear energy growth and damping rates, and $`\kappa `$ is the nonlinear coupling constant (cf. §3.1). Our amplitude equations differ only in notation from those given by Dziembowski (nl-dziem82 (1982)).
The instability threshold follows from a straightforward linear stability analysis applied to equations (2) and (3). The effects of nonlinear interactions on the parent mode are ignored as is appropriate for infinitesimal daughter mode amplitudes. The critical parent mode amplitude satisfies (Vandakurov nl-vandakurov79 (1979), Dziembowski nl-dziem82 (1982)).
$$|A_p|^2=\frac{\gamma _{d_1}\gamma _{d_2}}{18\kappa ^2\omega _{d_1}\omega _{d_2}}\left[1+\left(\frac{2\delta \omega }{\gamma _{d_1}+\gamma _{d_2}}\right)^2\right],$$
(4)
where $`\delta \omega \omega _{d_1}+\omega _{d_2}\omega _p`$.
### 2.2. Dynamics
The amplitude equations (1)-(3) have a unique equilibrium solution given by
$`|A_p|^2`$ $`=`$ $`{\displaystyle \frac{\gamma _{d_1}\gamma _{d_2}}{18\kappa ^2\omega _{d_1}\omega _{d_2}}}\left[1+\left({\displaystyle \frac{2\delta \omega }{\gamma _{d_1}+\gamma _{d_2}\gamma _p}}\right)^2\right],`$ (5)
$`|A_{d_1}|^2`$ $`=`$ $`{\displaystyle \frac{\gamma _{d_2}\gamma _p}{18\kappa ^2\omega _{d_2}\omega _p}}\left[1+\left({\displaystyle \frac{2\delta \omega }{\gamma _{d_1}+\gamma _{d_2}\gamma _p}}\right)^2\right],`$ (6)
$`|A_{d_2}|^2`$ $`=`$ $`{\displaystyle \frac{\gamma _{d_1}\gamma _p}{18\kappa ^2\omega _{d_1}\omega _p}}\left[1+\left({\displaystyle \frac{2\delta \omega }{\gamma _{d_1}+\gamma _{d_2}\gamma _p}}\right)^2\right],`$ (7)
together with
$$\mathrm{cot}\mathrm{\Phi }=\frac{2\delta \omega }{\gamma _{d_1}+\gamma _{d_2}\gamma _p}.$$
(8)
Here, $`\mathrm{\Phi }=\theta _{d_1}+\theta _{d_2}\theta _p`$, where the complex amplitude $`A_j`$ may be written as $`A_j=|A_j|e^{i\theta _j}`$. Note that equation (5) is almost identical to the threshold criterion (eq. ) for $`\gamma _p\gamma _d`$ which is the case that concerns us. Here $`\gamma _d=(\gamma _{d_1}+\gamma _{d_2})/2`$ is the characteristic damping rate for the daughter modes. It is also worth mentioning that, in this limit, the parent mode energy is independent of $`\gamma _p`$.
The equilibrium state is a stable attractor for mode triplets with $`|\delta \omega |>\gamma _d`$, and unstable otherwise (Wersinger et al. nl-wersinger80 (1980), Dziembowski nl-dziem82 (1982)). Figures 1a & 1b illustrate these two types of behaviors. Triplets with unstable equilibria undergo a variety of limit cycles. These share a number of common features. The parent mode’s amplitude remains close to its equilibrium (threshold) value, with slow rises on time scale $`\gamma _p^1`$ followed by precipitous drops on time scale $`\gamma _d^1`$. The daughter modes’ amplitudes stay far below their equilibrium values for most of the cycle, but peak with amplitudes comparable to that of the parent mode for a brief interval of length $`\gamma _d^1`$ shortly after the parent mode amplitude reaches its maximum value. During this brief interval, the energy which the parent mode has slowly accumulated is transferred to and dissipated by the daughter modes. So we see that, independent of its stability, the equilibrium state defines the parent mode’s amplitude. This is confirmed by Figure 2a. We employ this result in §3 where we predict upper limits for g-mode amplitudes in pulsating white dwarfs. The time-averaged energies of the daughter modes are also found to be consistent with equations (6)-(7) (Fig. 2b).
## 3. CHOOSING THE BEST DAUGHTER PAIRS
The discussion in the previous section shows that an overstable parent mode’s amplitude saturates at a value close to the threshold for parametric instability. Although each overstable parent mode has many potential daughter pairs, the most important pair is the one with the lowest instability threshold. This section is devoted to identifying these optimal daughter pairs, a task which separates into two independent parts, maximization of $`\kappa ^2`$ and minimization of $`(\delta \omega ^2+\gamma _d^2)`$.
### 3.1. Three-Mode Coupling Coefficients
The 3-mode coupling coefficient characterizes the lowest order nonlinear interactions among stellar modes. A compact form suitable for adiabatic modes under the Cowling approximation is derived in Kumar & Goldreich (nl-kumar89 (1989));
$`\kappa `$ $`=`$ $`{\displaystyle }d^3x{\displaystyle \frac{p}{6}}\{(\mathrm{\Gamma }_11)^2(\mathbf{}𝝃)^3`$ (9)
$`+3(\mathrm{\Gamma }_11)(\mathbf{}𝝃)\xi _{;j}^i\xi _{;i}^j+2\xi _{;j}^i\xi _{;k}^j\xi _{;i}^k\},`$
where $`p`$ is the unperturbed pressure, $`\mathrm{\Gamma }_1`$ is the adiabatic index, $`𝝃`$ is the Lagrangian displacement, the symbol ‘;’ denotes covariant derivative, and the integration is over the volume of the star. This expression for $`\kappa `$ is symmetric with respect to the three modes. Note that the displacements enter only through components of their gradients. The present form is not suitable for accurate numerical computation. We derive a more appropriate version in §A.1 (eq. \[A15\]).
Each eigenmode of a spherical star is characterized by three eigenvalues $`n,\mathrm{},m`$; $`n`$ is the number of radial nodes in the radial displacement eigenfunction, $`\mathrm{}`$ is the spherical degree, and $`m`$ is the azimuthal number.<sup>3</sup><sup>3</sup>3The angular dependence is described by a spherical harmonic $`Y_{lm}(\theta ,\varphi )`$. Integration over solid angle enforces the following selection rules on triplets with non-vanishing $`\kappa `$: $`|\mathrm{}_{d_2}\mathrm{}_{d_1}|\mathrm{}_p\mathrm{}_{d_1}+\mathrm{}_{d_2}`$, $`\mathrm{}_p+\mathrm{}_{d_1}+\mathrm{}_{d_2}`$ even, and $`m_p=m_{d_1}+m_{d_2}`$ (see §A.1). These selection rules guarantee the conservation of angular momentum during nonlinear interactions. The magnitude of $`\kappa `$ is largest when the eigenfunctions of the daughter modes are radially similar in the upper evanescent zone of the parent mode. Radial similarity requires near equality of the vertical components of WKB wavevectors, $`k_z`$, where for gravity modes $`k_z^2(N^2/\omega ^21)\mathrm{\Lambda }^2/r^2`$ (eq. \[A2\] of Paper I) with $`N^2`$ being the Brunt-Väisälä frequency, $`r`$ the radius, and $`\mathrm{\Lambda }^2=\mathrm{}(\mathrm{}+1)`$. We take $`N\omega _p`$ as the major contribution to the peak value of $`|\kappa |`$ comes from the region just above $`z_{\omega _p}`$, the upper boundary of the parent mode’s propagating cavity. The g-mode dispersion relation applicable for $`\omega 10^2\mathrm{s}^1`$ is $`\omega \mathrm{}/n`$ (see Fig. 4 of Paper I). We find that $`\kappa `$ is largest when
$$\frac{n_{d_1}}{n_{d_2}}\frac{\mathrm{\Lambda }_{d_1}\omega _{d_1}}{\mathrm{\Lambda }_{d_2}\omega _{d_2}}\left(\frac{\omega _p^2\omega _{d_1}^2}{\omega _p^2\omega _{d_2}^2}\right)^{1/2}.$$
(10)
Taking $`\mathrm{}_p=1`$, $`\mathrm{}_{d_2}=\mathrm{}_{d_1}+1`$, and evaluating the above relation for $`\mathrm{}_{d_1}<10`$, we locate the peak of $`\kappa `$ to be at
$$n_{d_1}n_{d_2}0.7n_p.$$
(11)
Since a fraction $`n_p^1`$ of the nodes of each daughter mode lie in the region above $`z_{\omega _p}`$, the peak has a width of order $`n_p`$ when measured in $`n_{d_1}n_{d_2}`$. Figure 3 illustrates the behavior of $`\kappa `$ for a variety of combinations of parent and daughter modes.
As we show in §A.2, the maximum value of $`|\kappa |`$ is of order
$$|\kappa |_{\mathrm{max}}\frac{1}{\left(n_p^3\tau _{\omega _p}L\right)^{1/2}}.$$
(12)
Here $`L`$ is the stellar luminosity and $`\tau _{\omega _p}`$ is the thermal timescale at $`z_{\omega _p}`$. This is compared with numerical results in Figure 4. Note that the maximum value of $`\kappa `$ depends entirely upon the properties of the parent mode and not at all upon those of its daughters.
### 3.2. Frequency Mismatch and Damping Rates
Because the maximum value of $`|\kappa |`$ is independent of $`\mathrm{}_{d_1}`$ and $`\mathrm{}_{d_2}`$, we choose these parameters to minimize $`(\delta \omega ^2+\gamma _d^2)`$.
Consider an $`\mathrm{}_p,m_p`$ parent mode. For each choice of $`\mathrm{}_{d_1},\mathrm{}_{d_2}`$, there are of order $`\mathrm{}_{d_1}n_p^2`$ daughter pairs for which $`\kappa `$ is close to its maximum value. The factor $`\mathrm{}_{d_1}`$ arises from the freedom in choosing $`m_{d_1}`$, while the factor $`n_p^2`$ comes from the width of maximum $`|\kappa |`$ at each $`\mathrm{}_{d_1}`$. Relaxing the value of $`\mathrm{}_{d_2}`$ subject to the constraint of the angular selection rules increases the number of pairs by a factor of order $`\mathrm{}_p`$. Now replace $`\mathrm{}_{d_1}`$ by a running variable $`\mathrm{}_{d_1}^{}{}_{}{}^{}`$. The number of pairs with $`\mathrm{}_{d_1}^{}{}_{}{}^{}\mathrm{}_{d_1}`$ is of order $`\mathrm{}_p\mathrm{}_{d_1}^2n_p^2`$. The distribution of the $`\delta \omega `$ values of these pairs is uniform between $`0`$ and $`n_p\omega _{d_1}/n_{d_1}`$.<sup>4</sup><sup>4</sup>4The factor $`n_p`$ arises because the peak in $`\kappa `$ has a width $`|n_{d_1}n_{d_2}|n_p`$. Statistically, the minimum frequency mismatch
$$\mathrm{\Delta }\omega \frac{\omega _{d_1}}{n_{d_1}}\frac{1}{\mathrm{}_p\mathrm{}_{d_1}^2n_p}\frac{\omega _p}{\mathrm{}_{d_1}^3n_p^2},$$
(13)
where the low frequency limit of the dispersion relation, $`\omega \mathrm{}/n`$, is assumed in going from the first to the second relation for $`\mathrm{\Delta }\omega `$. Our estimate for $`\mathrm{\Delta }\omega `$ assumes that rotation lifts $`m`$ degeneracy. If it does not, the minimum frequency mismatch is increased by a factor of $`\mathrm{}_{d_1}`$.
Next we describe how $`\gamma `$ varies with $`\omega `$ and $`\mathrm{}`$. There are two regimes of relevance to this investigation. In the quasiadiabatic limit (cf. §4.4 of paper I),<sup>5</sup><sup>5</sup>5Overstable modes are quasiadiabatic so this estimate for $`\gamma `$ applies to them as well as to damped modes.
$$\gamma \frac{1}{n\tau _\omega }\left(\frac{\mathrm{}}{\omega }\right)^6.$$
(14)
In the strongly nonadiabatic limit (see §B.1),
$$\gamma \frac{\omega }{\pi n}\mathrm{ln}\frac{1}{}\omega ^{0.75}\mathrm{}^{0.2},$$
(15)
where $``$ is the amplitude reflection coefficient at the top of the mode’s cavity. The transition between the quasiadiabatic and strongly nonadiabatic limits is marked by a significant reduction of $``$ by radiative diffusion. The behavior of $`\gamma `$ as a function of $`\omega `$ and $`\mathrm{}`$ is illustrated in Figure 5. Radial similarity of the daughter modes to which a given parent mode couples most strongly (eq.) implies $`\gamma _{d_1}\gamma _{d_2}\gamma _d`$.
The minimum of $`(\delta \omega ^2+\gamma _d^2)`$ is attained for daughter pairs that satisfy $`\delta \omega ^2\gamma _d^2`$. With increasing $`n_p`$, the value of $`\mathrm{}_{d_1}`$ at which this occurs decreases from a few to unity (see Fig. 6). For sufficiently large values of $`n_p`$, $`\mathrm{\Delta }\omega \gamma _d`$ even for $`\mathrm{}_{d_1}=1`$.
## 4. UPPER ENVELOPE OF PARENT MODE AMPLITUDES
Parametric instability provides an upper envelope to the amplitudes of overstable modes. Coupling of an overstable parent mode to a single pair of daughter modes suffices to maintain the parent mode’s amplitude near the threshold value (cf. §2.2). The energy gained by the overstable parent mode and transferred to the daughter modes may be disposed of by linear radiative damping or by further nonlinear coupling to granddaughter modes.
The results presented in this section are obtained from numerical computations and displayed in a series of figures. The general trends they exhibit are best understood in terms of analytic scaling relations. We derive these first in order to be able to refer to them as we describe each figure.
Excited modes of ZZ Ceti stars are usually detected through photometric measurements of flux variations, and in a few cases through spectroscopic measurements of horizontal velocity variations. Hence we calculate surface amplitudes of fractional flux, $`\mathrm{\Delta }F/F`$, and horizontal velocity, $`v_h`$, variations.<sup>6</sup><sup>6</sup>6In this section we drop subscripts on parent mode parameters and denote daughter mode parameters by a subscript $`d`$. Each of these is directly related to the near surface amplitude of the Lagrangian pressure perturbation, $`\delta p/p`$. The threshold value of the latter is obtained by combining equations (4) and (12) with the normalization factor $`(n\tau _\omega L)^{1/2}`$ given by equation (A28) of Paper I. Thus
$$\frac{\delta p}{p}n\left[\left(\frac{\gamma _d}{\omega }\right)^2+\left(\frac{\delta \omega }{\omega }\right)^2\right]^{1/2}.$$
(16)
Adopting relations expressing $`v_h`$ and $`\mathrm{\Delta }F/F`$ in terms of $`\delta p/p`$ from §3 of Paper I, we arrive at
$$v_h\frac{\omega Rn}{[\mathrm{}(\mathrm{}+1)]^{1/2}}\left[\left(\frac{\gamma _d}{\omega }\right)^2+\left(\frac{\delta \omega }{\omega }\right)^2\right]^{1/2},$$
(17)
and
$$\frac{\mathrm{\Delta }F}{F}\frac{n}{\left[1+\left(\omega \tau _c\right)^2\right]^{1/2}}\left[\left(\frac{\gamma _d}{\omega }\right)^2+\left(\frac{\delta \omega }{\omega }\right)^2\right]^{1/2},$$
(18)
In the above, $`R`$ denotes the stellar radius, and $`\tau _c`$ is the thermal time constant describing the low pass filtering action of the convection zone on flux variations input at its base.<sup>7</sup><sup>7</sup>7$`\tau _c3\tau _{\mathrm{th}}`$, where $`\tau _{\mathrm{th}}`$ is evaluated at $`z_b`$.
Figure 7 displays calculated values for amplitudes of overstable modes limited by parametric instability. The rise of $`|\delta p/p|`$ with increasing mode period mainly reflects the corresponding rise of the damping rates of the daughter modes (cf. §3.2) as indicated by equation (16). At the longer periods, the values of $`|\delta p/p|`$ decline with decreasing $`T_{\mathrm{eff}}`$. This is a subtle consequence of the deepening of the convection zone which pushes down the top of the daughter modes’ cavities, thus reducing their damping rates (see §5.2 of Wu & Goldreich nl-paperII (1999), hereafter Paper II). The behavior of $`|v_h|`$ is similar to that of $`|\delta p/p|`$ except that $`|v_h|`$ decreases relative to $`|\delta p/p|`$ with increasing $`\mathrm{}`$ as shown by comparison of equations (16) and (17). A new feature present in the run of $`|\delta F/F|`$ verses mode period is the low pass filtering action of the convection zone as expressed by the factor $`[1+(\omega \tau _c)^2]^{1/2}`$ in equation (18). This factor causes $`|\delta F/F|`$ to rise slightly more steeply than $`|\delta p/p|`$ with increasing mode period. It is also responsible for a more dramatic decrease in $`|\delta F/F|`$ with decreasing $`T_{\mathrm{eff}}`$ at fixed mode period. Since $`|v_h|`$ does not suffer from this visibility reduction, velocity variations may be observable in stars that are cooler than those at the red edge of the instability strip.
Figure 8 reproduces a summary of observational data on mode amplitudes from Clemens (amp-clemens95 (1995)). Each star is represented by a point showing the relation between the V-band photometric amplitude in its largest mode and the amplitude weighted mean period of all its observed modes. The latter quantity is a surrogate for the star’s effective temperature in the sense that longer mean periods correspond to lower effective temperatures (Clemens amp-clemens95 (1995), Paper I). Our theoretical predictions for the amplitudes of $`\mathrm{}=1`$ modes limited by parametric instability are shown by a solid line obtained by interpolation from the numerical results displayed in Figure 7. Two features of this figure are worthy of comment.
For a few low order overstable modes, $`n_p3`$ at $`\mathrm{}_p=1`$, $`\delta \omega `$ is more significant than $`\gamma _d`$ in determining the best daughter pair (Fig. 6). The best daughter pairs for these modes have $`\mathrm{}`$ values of a few and identities which depend sensitively on minor differences among stars. Thus we expect amplitudes of short period overstable modes to show large star to star variations. That the $`n=1,\mathrm{}=1`$ mode is detected in only about one half of the hot DAVs (see Figure 4 of Clemens amp-clemens95 (1995)) is consistent with the expected statistical variations in $`\delta \omega `$.
Overstable $`\mathrm{}=1`$ modes with periods longer than $`800\mathrm{s}`$ ($`n_p20`$) are probably not saturated by parametric instability. Their maximum $`|\kappa |`$ is severely reduced below the adiabatic value by the strong nonadiabaticity of their daughter modes (§B.2). Other mechanisms that contribute towards saturating these modes are discussed in §5.
## 5. DISCUSSION
### 5.1. Turbulent Saturation
Turbulent convection severely reduces the vertical gradient of the horizontal velocity of g-modes in the convection zone. As a result, a shear layer forms at the boundary between the bottom of the convection zone and the top of the radiative interior (Goldreich & Wu nl-paperIII (1999), hereafter Paper III). Kelvin-Helmholtz instability of this layer provides a nonlinear dissipation mechanism for overstable modes. An overstable mode’s amplitude cannot grow beyond the value at which nonlinear damping due to the Kelvin-Helmholtz instability balances its linear convective driving. Equation (44) of Paper III provides an estimate for the value at which this mechanism saturates the surface amplitude of $`(\delta p/p)`$;
$$\left(\frac{\delta p}{p}\right)\frac{0.1}{C_D}\frac{[(\omega \tau _c)^2+1]^{1/2}[(\omega \tau _c)^21]}{\omega \tau _c}\frac{Lz_\omega ^2}{Rz_b}.$$
(19)
Our ignorance of the complicated physics involved in a nonlinear shear layer is covered by the range of possible values of the dimensionless drag coefficient, $`C_D`$. Terrestrial experiments indicate that $`C_D`$ falls between $`10^3`$ and $`10^1`$. The dashed lines in Figure 8 show the effect of including nonlinear turbulent damping in addition to parametric instability on limiting mode amplitudes. Amplitudes of overstable modes saturated by the Kelvin-Helmholtz instability are uncertain because $`C_D`$ is poorly constrained.
### 5.2. Granddaughter Modes
Here we answer the following questions. Under what conditions do daughter modes excite granddaughter modes by parametric instability?<sup>8</sup><sup>8</sup>8In this subsection subscripts $`p`$, $`d`$, and $`g`$ refer to parent, daughter, and granddaughter modes. What are the consequences if they do? In §2.2 we show how parametric instability of linearly damped daughter modes maintains the amplitude of a parent mode close to its equilibrium value. In order to dispose of the energy they receive from the parent mode, the time averaged energies of the daughter modes must be close to their equilibrium values. This raises a worry. Suppose the daughters are prevented from reaching their equilibrium amplitudes by parametric instability of granddaughter modes Then they would not be able to halt the amplitude growth of the parent mode.
To answer the first of these questions, we calculate the ratio, denoted by the symbol $`𝒮`$, between the threshold amplitude for a daughter mode to excite granddaughter modes and its equilibrium amplitude under parametric excitation by the parent mode. The former is obtained from equation (4), and the latter from equations (6)-(7). We make a few simplifying assumptions to streamline the discussion. Resonances between daughters and granddaughters are taken as exact; individual members of daughter and granddaughter pairs are treated as equivalent. Equations (12) and (14) are combined to yield
$$\kappa ^2\frac{\gamma _p}{n_p^2L}.$$
(20)
It is then straightforward to show that
$$𝒮\frac{\omega _p\omega _d}{\omega _g^2}\left(\frac{n_d}{n_p}\right)^2\left(\frac{\gamma _g^2}{\gamma _d^2+\delta \omega ^2}\right)32\left(\frac{\gamma _g^2}{\gamma _d^2+\delta \omega ^2}\right).$$
(21)
The factor $`32`$ is an approximation based on taking $`\omega _g/\omega _d=\omega _d/\omega _p=1/2`$ and $`n_p/n_d=1/2`$.
In general $`𝒮1`$, so the excitation of granddaughter modes requires the daughter modes to have energies in excess of their equilibrium values. However, the equilibrium solution is unstable if the best daughter pair corresponds to $`\gamma _d>|\delta \omega |`$, and then the daughter mode energies episodically rise far above their equilibrium values. At such times, granddaughter modes may be excited by parametric instability and consequentially limit the amplitude growth of the daughter modes. This slows the transfer of energy from parent to daughter modes, but it does not prevent the daughter modes from saturating the growth of the parent mode’s amplitude at the level described by equation (4).
For the few lowest order parent modes, we typically find $`|\delta \omega |\gamma _d`$. This may reduce $`𝒮`$ to below unity with the consequence that the parent mode amplitude may rise above that given by equation (4).
### 5.3. Additional 3-Mode Interactions
Parametric instability sets reasonable upper bounds on the photospheric amplitudes of overstable modes. In a given star this upper bound rises with increasing mode period except possibly for the lowest few modes. However, the observed amplitude distributions are highly irregular. This mode selectivity may arise from 3-mode interactions which involve more than one overstable mode.
We investigate a particular example of this type. It is closely related to parametric instability, the only difference being that the daughter modes of the overstable parent mode are themselves overstable. Acting in isolation, resonant mode couplings tend to drive mode energies toward equipartition. They conserve the total energy, $`\dot{E}_p+\dot{E}_{d_1}+\dot{E}_{d_2}=0`$, and transfer action according to $`\dot{E}_{d_1}/\omega _{d_1}=\dot{E}_{d_2}/\omega _{d_2}=\dot{E}_p/\omega _p`$. In this context it is important to note that the energies of overstable modes limited by parametric instability decline with increasing mode period (Fig. 9). Therefore nonlinear interactions transfer energy from the parent mode to its independently excited daughters. As shown below, this transfer may severely suppress the parent mode’s amplitude.
We start from equations (1)-(3). These may be manipulated to yield
$$\frac{dE_p}{dt}=\gamma _pE_p+3\sqrt{2}\omega _p\kappa (E_pE_{d_1}E_{d_2})^{1/2}\mathrm{sin}\mathrm{\Phi },$$
(22)
where $`\mathrm{\Phi }=\theta _{d_1}+\theta _{d_2}\theta _p`$. For $`E_pE_{d_1}`$ and $`E_{d_2}`$, nonlinear interactions transfer energy from the parent mode to its daughter modes. In particular, if we ignore phase changes in the overstable daughter modes due to their interactions with granddaughter modes, we find that $`\mathrm{\Phi }`$ satisfies
$$\frac{d\mathrm{\Phi }}{dt}=\delta \omega \frac{3}{\sqrt{2}}\kappa (E_pE_{d_1}E_{d_2})^{1/2}\left[\left(\frac{\omega _{d_1}}{E_{d_1}}+\frac{\omega _{d_2}}{E_{d_2}}\right)\frac{\omega _p}{E_p}\right]\mathrm{cos}\mathrm{\Phi },$$
(23)
with a stable solution at $`\mathrm{\Phi }=\pi /2`$ when $`\delta \omega =0`$.
We denote the ratio of the nonlinear term to the linear term in equation (22) by the symbol $`𝒯`$;
$$𝒯\frac{3\sqrt{2}\omega _p\kappa }{\gamma _p}\left(\frac{E_{d_1}E_{d_2}}{E_p}\right)^{1/2}.$$
(24)
Using the magnitudes of $`E_i`$ set by parametric instability of their respective daughters, and adopting the same approximations made in §5.2, we arrive at
$$𝒯\frac{\omega _p\omega _d}{\omega _g^2}\left(\frac{n_d}{n_p}\right)^2\left(\frac{\gamma _g^2}{\gamma _d^2+\delta \omega ^2}\right)32\left(\frac{\gamma _g^2}{\gamma _d^2+\delta \omega ^2}\right).$$
(25)
Comparing equations (21) and (25) we see that $`𝒮=𝒯`$. A little thought reveals that this is not a coincidence. For $`𝒯1`$, overstable daughter modes can suppress a parent mode’s energy below the value set by parametric instability. We expect this suppression to be important in cool ZZ Ceti stars whose overstable modes extend to long periods. It may render their intermediate period modes invisible. In a similar manner, the amplitudes of high frequency overstable modes with $`\mathrm{}=2`$ and $`3`$ may be heavily suppressed by interactions with their lowest $`\mathrm{}`$ overstable daughters.
The irregular amplitude distributions among neighboring modes may be partially accounted for by this type of resonance. Mode variability may also play a role. We explore this in the next subsection.
### 5.4. Mode Variability
Excited g-modes in ZZ Ceti stars exhibit substantial temporal variations. Parametric instability may at least partially account for these variations.
When $`|\delta \omega |<\gamma _d`$, parametric instability gives rise to limit cycles in which the amplitudes and phases of parent and daughter modes vary on time scales as short as $`\gamma _d^1`$. Stable daughter modes may briefly attain visible amplitudes. Temporal amplitude variations may contribute to the irregular mode amplitude distribution seen in individual stars.
Phase variations of a parent mode obey the equation
$$\frac{d\theta _p}{dt}=\omega _p\frac{3}{\sqrt{2}}\omega _p\kappa \frac{|A_{d_1}||A_{d_2}|}{|A_p|}\mathrm{cos}\mathrm{\Phi }.$$
(26)
At the equilibrium given by equations (5)-(8), the parent mode’s frequency is displaced from its unperturbed value such that
$$\omega _p^{}=\frac{d\theta _p}{dt}=\omega _p+\frac{\delta \omega \gamma _p}{\gamma _{d_1}+\gamma _{d_2}\gamma _p}\sqrt{1+\left(\frac{2\delta \omega }{\gamma _{d_1}+\gamma _{d_2}\gamma _p}\right)^2}.$$
(27)
This constant frequency shift is of order $`10^9s^1`$ for the $`n=1`$, $`\mathrm{}=1`$ mode and of order $`10^7s^1`$ for the $`n=2`$, $`\mathrm{}=2`$ mode. Frequency shifts in higher order overstable modes which are involved in limit cycles are predicted to be larger and time variable. During brief intervals of length $`\gamma _d^1`$, when the daughter mode energies are comparable to that of the parent mode, $`|\omega _p^{}\omega _p|\gamma _d`$, which is of order a few times $`10^5s^1`$, or a few $`\mu Hz`$ in angular frequency. These shifts might account for the time-varying rotational splittings reported by Kleinman et al. (kleinmanZZPsc (1998)) provided different $`m`$ components of the overstable modes are involved in different limit cycles.
### 5.5. Miscellany
We briefly comment on two relevant issues.
Gravitational settling produces chemically pure layers between which modes can be partially trapped. Modes that are trapped in the hydrogen layer have lower mode masses, and therefore higher growth rates and larger maximum values of $`|\kappa |`$ than untrapped modes of similar frequency. This implies lower threshold energies for parametric instability. Nevertheless, trapping does not affect predicted photospheric amplitudes of $`\delta p/p`$, and hence $`v_h`$ and $`\delta F/F`$, since these are proportional to $`A`$ divided by the square root of the mode mass.
In circumstances of small rotational splitting, the simple limit-cycles depicted in Figure 1 are unlikely to be realistic. In such cases, different $`m`$ components of an overstable parent mode share some common daughter modes. This leads to more complex dynamics.
We are indebted to Bradley for supplying us with models of DA white dwarfs. Financial support for this research was provided by NSF grant 94-14232.
## Appendix A Three-Mode Coupling Coefficient
We reproduce the expression for the three-mode coupling coefficient presented in equation (9) of §3.1 with one modification:
$$\kappa =d^3xp\left\{\frac{\mathrm{\Gamma }_1(\mathrm{\Gamma }_12)}{6}(\mathbf{}𝝃)^3+\frac{1}{2}\mathrm{\Gamma }_1(\mathbf{}𝝃)\xi _{;j}^i\xi _{;i}^j+\mathrm{Det}|\xi _{;j}^i|\right\},$$
(A1)
where
$$\mathrm{Det}|\xi _{;j}^i|=\frac{1}{6}(\mathbf{}𝝃)^3\frac{1}{2}(\mathbf{}𝝃)\xi _{;j}^i\xi _{;i}^j+\frac{1}{3}\xi _{;j}^i\xi _{;k}^j\xi _{;i}^k.$$
(A2)
Because of strong cancellations among its largest terms, this form is not well-suited for numerical evaluation. We derive a new expression which does not suffer from this defect. Then we estimate the size of $`\kappa `$ and deduce its dependences upon the properties of the three-modes.
### A.1. Simplification
In spherical coordinates $`(r,\theta ,\varphi )`$, the components of the displacement vector may be written as
$$𝝃=\xi ^i𝜺_i=[\xi _r(r),\frac{\xi _h(r)}{r}\frac{}{\theta },\frac{\xi _h(r)}{r}\frac{1}{\mathrm{sin}^2\theta }\frac{}{\varphi }]Y_\mathrm{}m(\theta ,\varphi ),$$
(A3)
where $`Y_\mathrm{}m`$ is a spherical harmonic function, and the $`𝜺_i`$ are covariant basis vectors.
The angular integrations in equation (A1) are done analytically. The following definitions and properties prove useful:
$`T`$ $``$ $`{\displaystyle _0^{2\pi }}𝑑\varphi {\displaystyle _0^\pi }𝑑\theta \mathrm{sin}\theta Y_aY_bY_c<Y_aY_bY_c>,`$ (A4a)
$`F_a`$ $``$ $`<Y_a_\alpha Y_b^\alpha Y_c>=\left(\mathrm{\Lambda }_b^2+\mathrm{\Lambda }_c^2\mathrm{\Lambda }_a^2\right){\displaystyle \frac{T}{2}},`$ (A4b)
$`G_a`$ $``$ $`<_\alpha _\beta Y_a^\alpha Y_b^\beta Y_c>\left[\mathrm{\Lambda }_a^4(\mathrm{\Lambda }_b^2\mathrm{\Lambda }_c^2)^2\right]{\displaystyle \frac{T}{4}},`$ (A4c)
$`S`$ $``$ $`<_\alpha (^\alpha Y_a_\beta Y_b^\beta Y_c)>=G_a+G_b+G_c`$ (A4d)
$`=`$ $`{\displaystyle \frac{1}{2}}(\mathrm{\Lambda }_a^2F_a+\mathrm{\Lambda }_b^2F_b+\mathrm{\Lambda }_c^2F_c)=G_a+\mathrm{\Lambda }_a^2F_a,`$
$`V_a`$ $``$ $`<Y_a_\alpha ^\beta Y_b_\beta ^\alpha Y_c>=\mathrm{\Lambda }_b^2\mathrm{\Lambda }_c^2TF_aS,`$ (A4e)
$`\mathrm{\Lambda }_a^2T`$ $`=`$ $`F_b+F_c,`$ (A4f)
$`\mathrm{\Lambda }_a^2F_a`$ $`=`$ $`G_b+G_c.`$ (A4g)
Here, $`_\alpha `$ is the covariant derivative on a spherical surface; $`\alpha `$ can be either $`\theta `$ or $`\varphi `$. Each angular integration is proportional to $`T`$ which contains all $`m`$ dependences and is of order unity independent of the $`\mathrm{}`$ values of the participating modes. The paramter $`\mathrm{\Lambda }^2\mathrm{}(\mathrm{}+1)`$. Subscripts $`a`$, $`b`$ and $`c`$ denote different modes. The angular selection rules are simply, $`\mathrm{}_c[|\mathrm{}_a\mathrm{}_b|,\mathrm{}_a+\mathrm{}_b]`$, Mod$`[(\mathrm{}_c+\mathrm{}_b+\mathrm{}_c),2]=0`$ and $`m_a+m_b+m_c=0`$.
Angular integration of $`\mathrm{Det}|\xi _{;j}^i|`$ yields
$`{\displaystyle 𝑑\mathrm{\Omega }\mathrm{Det}|\xi _{;j}^i|}`$ $`=`$ $`{\displaystyle \frac{1}{r^2}}\left[T\xi _r^b\xi _r^c\mathrm{\Lambda }_c^2T\xi _r^b\xi _h^c+{\displaystyle \frac{1}{2}}(\mathrm{\Lambda }_b^2\mathrm{\Lambda }_c^2TV_a)\xi _h^b\xi _h^c\right]{\displaystyle \frac{d\xi _r^a}{dr}}`$ (A5)
$`+{\displaystyle \frac{1}{r^2}}\left[S(\xi _r^b\xi _h^b)\xi _h^cF_b\xi _r^b(\xi _r^c\xi _h^c)\right]{\displaystyle \frac{d\xi _h^a}{dr}}.`$
Symmetrizing this expression with respect to modes $`b`$ and $`c`$, employing equations (A4d)-(A4g), we get
$$𝑑\mathrm{\Omega }\mathrm{Det}|\xi _{;j}^i|=\frac{S}{3r^2}\frac{d}{dr}\left(\xi _h^a\xi _h^b\xi _h^c\right)+\frac{(F_a+S)}{2r^2}\frac{d}{dr}\left(\xi _r^a\xi _h^b\xi _h^c\right)\frac{\mathrm{\Lambda }_a^2T}{2r^2}\frac{d}{dr}\left(\xi _h^a\xi _r^b\xi _r^c\right)+\frac{T}{3r^2}\frac{d}{dr}\left(\xi _r^a\xi _r^b\xi _r^c\right).$$
(A6)
For gravity-modes, $`|\xi _h||\xi _r|`$. Thus the four terms in this expression decrease in size from left to right.
In a similar manner we arrive at
$$𝑑\mathrm{\Omega }(\mathbf{}𝝃)^3=T(\mathbf{}𝝃^a)(\mathbf{}𝝃^b)(\mathbf{}𝝃^c),$$
(A7)
and
$`{\displaystyle 𝑑\mathrm{\Omega }(\mathbf{}𝝃)\xi _{;j}^i\xi _{;i}^j}`$ $`=`$ $`{\displaystyle \frac{F_a}{r}}(\mathbf{}𝝃^a){\displaystyle \frac{d}{dr}}\left(\xi _h^b\xi _h^c\right)+{\displaystyle \frac{2F_a}{r}}(\mathbf{}𝝃^a)\xi _r^b{\displaystyle \frac{d\xi _h^c}{dr}}+T(\mathbf{}𝝃^a){\displaystyle \frac{d\xi _r^b}{dr}}{\displaystyle \frac{d\xi _r^c}{dr}}`$ (A8)
$`+{\displaystyle \frac{V_a}{r^2}}(\mathbf{}𝝃^a)\xi _h^b\xi _h^c{\displaystyle \frac{2T\mathrm{\Lambda }_c^2}{r^2}}(\mathbf{}𝝃^a)\xi _r^b\xi _h^c+{\displaystyle \frac{2T}{r^2}}(\mathbf{}𝝃^a)\xi _r^b\xi _r^c.`$
The magnitudes of the terms in the above expression decrease from left to right except that terms 2-4 are of comparable size.
Having disposed of the angular dependences in $`\kappa `$, we turn to the radial integrations. The following relations prove helpful in this context:
$`{\displaystyle \frac{d\xi _r}{dr}}`$ $`=`$ $`(\mathbf{}𝝃)+{\displaystyle \frac{\mathrm{\Lambda }^2}{r}}\xi _h{\displaystyle \frac{2}{r}}\xi _r,`$ (A9)
$`{\displaystyle \frac{d}{dr}}\left[\mathrm{\Gamma }_1p(\mathbf{}𝝃)\right]`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Lambda }^2\rho g}{r}}\xi _h\left(\omega ^2+{\displaystyle \frac{2g}{r}}{\displaystyle \frac{dg}{dr}}\right)\rho \xi _r,`$ (A10)
$`\mathrm{\Gamma }_1p(\mathbf{}𝝃)`$ $`=`$ $`\rho g\xi _rr\rho \omega ^2\xi _h.`$ (A11)
The first equation is the definition of divergence, and the second and third are the radial and horizontal components of the equation of motion written in terms of the Lagrangian displacement.
Curvature terms arise because the directions of the basis vectors depend upon position. For instance,
$$\xi _{;\theta }^r=(\xi _r\xi _h)\frac{Y_\mathrm{}m}{\theta },$$
(A12)
where $`\xi _h`$ appears due to curvature in the coordinate system. The largest terms in equations (A6) and (A8) are curvature terms. However, their radial integrals cancel leaving a much smaller net contribution. Direct numerical evaluation of equation (A1) leads to unreliable results, so it is important to carry out this cancellation analytically. Accordingly, we integrate each term by parts and apply equations (A4d) and (A9) to obtain
$`{\displaystyle _0^R}𝑑r\left[{\displaystyle \frac{Sp}{3}}{\displaystyle \frac{d}{dr}}\left(\xi _h^a\xi _h^b\xi _h^c\right)+{\displaystyle \frac{F_ar\mathrm{\Gamma }_1p}{2}}(\mathbf{}𝝃^a){\displaystyle \frac{d}{dr}}\left(\xi _h^b\xi _h^c\right)\right]`$
$`={\displaystyle _0^R}𝑑r{\displaystyle \frac{F_a\mathrm{\Gamma }_1p}{2}}(\mathbf{}𝝃^a)\xi _h^b\xi _h^c+{\displaystyle _0^R}𝑑r{\displaystyle \frac{F_ar\rho }{2}}\left(\omega _a^2+{\displaystyle \frac{2g}{r}}{\displaystyle \frac{dg}{dr}}\right)\xi _r^a\xi _h^b\xi _h^c.`$ (A13)
This step leads to
$`\kappa `$ $`=`$ $`{\displaystyle _0^R}dr[{\displaystyle \frac{Tr^2\mathrm{\Gamma }_1(\mathrm{\Gamma }_12)p}{6}}(\mathbf{}𝝃^a)(\mathbf{}𝝃^b)(\mathbf{}𝝃^c){\displaystyle \frac{Tr^2\mathrm{\Gamma }_1p}{2}}(\mathbf{}𝝃^a){\displaystyle \frac{d\xi _r^b}{dr}}{\displaystyle \frac{d\xi _r^c}{dr}}`$ (A14)
$`F_ar\mathrm{\Gamma }_1p(\mathbf{}𝝃^a)\xi _r^b{\displaystyle \frac{d\xi _h^c}{dr}}{\displaystyle \frac{(\mathrm{\Lambda }_b^2\mathrm{\Lambda }_c^2TS)\mathrm{\Gamma }_1p}{2}}(\mathbf{}𝝃^a)\xi _h^b\xi _h^c{\displaystyle \frac{(SF_a)\rho g}{2}}\xi _r^a\xi _h^b\xi _h^c{\displaystyle \frac{F_ar\rho }{2}}{\displaystyle \frac{dg}{dr}}\xi _r^a\xi _h^b\xi _h^c`$
$`+\mathrm{\Lambda }_a^2T\rho g\xi _h^a\xi _r^b\xi _r^c+\mathrm{\Lambda }_c^2T\mathrm{\Gamma }_1p(\mathbf{}𝝃^a)\xi _r^b\xi _h^cTp{\displaystyle \frac{d\xi _r^a}{dr}}\xi _r^b\xi _r^cT\mathrm{\Gamma }_1p(\mathbf{}𝝃^a)\xi _r^b\xi _r^c+{\displaystyle \frac{\omega _a^2F_ar\rho }{2}}\xi _r^a\xi _h^b\xi _h^c].`$
Next we systematically eliminate radial derivatives of the displacement vector from the expression for $`\kappa `$. We integrate by parts to dispose of $`d\xi _h/dr`$ and substitute for $`d\xi _r/dr`$ using equation (A9). With the aid of equations (A4d)-(A4f), and using equation (A11) to make the expression symmetric with respect to indexes $`b`$ and $`c`$, we arrive at our final working expression for $`\kappa `$,
$`\kappa `$ $`=`$ $`{\displaystyle _0^R}dr\{{\displaystyle \frac{Tr^2\mathrm{\Gamma }_1(\mathrm{\Gamma }_1+1)p}{6}}(\mathbf{}𝝃^a)(\mathbf{}𝝃^b)(\mathbf{}𝝃^c)+{\displaystyle \frac{\omega _a^2G_ar\rho }{2}}\xi _h^a\xi _h^b\xi _h^c`$ (A15)
$`+{\displaystyle \frac{F_a\rho }{2}}\left(gr{\displaystyle \frac{dg}{dr}}\right)\xi _r^a\xi _h^b\xi _h^c{\displaystyle \frac{\mathrm{\Lambda }_a^2Tr\mathrm{\Gamma }_1p}{2}}\xi _h^a(\mathbf{}𝝃^b)(\mathbf{}𝝃^c)T(3\mathrm{\Gamma }_1+1)p(\mathbf{}𝝃^a)\xi _r^b\xi _r^c`$
$`+2Tr\mathrm{\Gamma }_1p\xi _r^a(\mathbf{}𝝃^b)(\mathbf{}𝝃^c)+{\displaystyle \frac{2Tp}{r}}\xi _r^a\xi _r^b\xi _r^c+{\displaystyle \frac{1}{2}}\left[\left(\omega _a^23\omega _b^23\omega _c^2\right)F_a(2\omega _b^2F_b+2\omega _c^2F_c)\right]r\rho \xi _r^a\xi _h^b\xi _h^c`$
$`+{\displaystyle \frac{1}{2}}[\mathrm{\Lambda }_a^2T(5\rho g+\rho r{\displaystyle \frac{dg}{dr}}{\displaystyle \frac{2p}{r}})\rho (\omega _b^2F_b+\omega _c^2F_c)r]\xi _h^a\xi _r^b\xi _r^c\}`$
All permutations of the 3 modes are to be included when evaluating this expression. However, the largest contribution from each term comes when $`b`$ and $`c`$ are radially similar daughter modes. For high order gravity-modes, the first five terms are of comparable size and much larger than the remaining four. Numerical evaluation of $`\kappa `$ using equation (A15) does not suffer from the errors arising from large cancellations and numerical differentiation that plague attempts using equation (A1).
### A.2. Order-of-Magnitude
Here we estimate the maximum value that $`\kappa `$ can attain for parametric resonances involving a given parent mode. In so doing, we apply results derived in Papers I and II. These include: (1) the scaling relations
$$|\mathbf{}𝝃|\frac{\mathrm{\Lambda }^2}{R}\xi _h\frac{\xi _r}{z_\omega }k_h\xi _h\frac{1}{(n\tau _\omega L)^{1/2}},$$
(A16)
in the evanescent region $`z<z_\omega `$, and
$$|\mathbf{}𝝃|\left(\frac{z_\omega }{z}\right)^{\frac{1}{2}}\frac{\mathrm{\Lambda }^2}{R}\xi _h\frac{\xi _r}{z},$$
(A17)
in the propagating cavity $`z>z_\omega `$; (2) the fact that regions between consecutive radial nodes contribute equally to the following normalization integral
$$\frac{\omega ^2}{2}_0^R𝑑rr^2\rho \left(\mathrm{\Lambda }^2\xi _h^2+\xi _r^2\right)=1.$$
(A18)
As $`\xi _r^2\xi _h^2`$ for g-modes, we have
$$\omega ^2R^2_0^z𝑑z\rho \mathrm{\Lambda }^2\xi _h^2\frac{n^{}}{n}\frac{1}{n}_0^z𝑑zk_z\frac{1}{n}\left(\frac{z}{z_\omega }\right)^{\frac{1}{2}}.$$
(A19)
Here $`n^{}`$ is the number of radial nodes above depth $`z`$ and $`n`$ the total number in the mode.
Equation (A15) yields maximal values for $`\kappa `$ when mode $`a`$ is taken to be the parent mode and modes $`b`$ and $`c`$ to be two radially similar daughter modes. Most of the contribution to the radial integral comes from the region above $`z_{\omega _a}`$: for parametric resonance, $`z_{\omega _a}`$ is much greater than $`z_{\omega _b}z_{\omega _c}`$; the decay and rapid oscillation of the parent mode’s eigenfunction renders insignificant contribution from greater depths. Thus we can take the integrals in equation (A15) to run from $`z=0`$ to $`z=z_{\omega _a}`$ and pull out the parent mode eigenfunctions since they are approximately constant for $`z<z_{\omega _a}`$. This procedure reduces each of the leading terms in $`\kappa `$ and thus their sum to
$$\kappa \left(\frac{z_{\omega _p}}{z_{\omega _d}}\right)^{1/2}\frac{1}{n_d(n_p\tau _{\omega _p}L)^{1/2}}.$$
(A20)
But $`z_{\omega _p}/z_{\omega _d}\mathrm{\Lambda }_d^2/\mathrm{\Lambda }_p^2n_d^2/n_p^2`$ which leads to
$$\kappa \frac{1}{(n_p^3\tau _{\omega _p}L)^{1/2}}.$$
(A21)
The above equation establishes that the maximum value of $`\kappa `$ rises steeply with increasing radial order of the parent mode and is independent of the radial orders and spherical degrees of the daughter modes.
## Appendix B STRONGLY NONADIABATIC DAUGHTER MODES
Strong nonadiabaticity occurs wherever
$$\frac{\omega \tau _{\mathrm{th}}}{(k_zz)^2}1.$$
(B1)
We refer to the depth above which this inequality applies as $`z_{\mathrm{na}}`$. To derive a simple analytic scaling relation for $`z_{\mathrm{na}}`$, we make use of the approximations $`k_z(zz_\omega )^{1/2}`$ and $`\tau _{\mathrm{th}}/\tau _b(z/z_b)^{q+2}`$. Here $`z_b`$ and $`\tau _b`$ are the depth and thermal relaxation time at the bottom of the convection zone, and $`\rho z^q`$ with $`q3.5`$ provides a fit to the density structure in the upper portion of the radiative interior. It then follows that
$$\frac{z_{\mathrm{na}}}{z_b}\left(\frac{1}{\omega \tau _b}\frac{z_b}{z_\omega }\right)^{\frac{1}{(q+1)}},$$
(B2)
Moreover,
$$\left(k_zz\right)_{\mathrm{na}}\left[\frac{1}{(\omega \tau _b)}\left(\frac{z_b}{z_\omega }\right)^{q+2}\right]^{\frac{1}{2(q+1)}}.$$
(B3)
By reducing the effective buoyancy, strong nonadiabaticity lowers the effective lid of a g-mode’s cavity to $`z_{\mathrm{na}}`$. Consequences of this fact are explored in the following subsections.
### B.1. Damping Rates of Strongly Nonadiabatic Modes
As shown in Paper II, the energy dissipation rate for a strongly nonadiabatic mode may be written as
$$\gamma \frac{\omega }{\pi n}\mathrm{ln}^1,$$
(B4)
where $``$ denotes the coefficient of amplitude reflection at the top of the mode’s cavity.
To derive an approximate relation for $``$, we note that the real and imaginary parts of $`k_z`$, $`k_{zr}`$ and $`k_{zi}`$, satisfy
$$\frac{|k_{zi}|}{|k_{zr}|}\frac{(k_{zr}z)^2}{\omega \tau _{\mathrm{th}}},$$
(B5)
provided $`|k_{zi}|/|k_{zr}|1`$. Thus
$$\mathrm{ln}^1_{z_{\mathrm{na}}}^{\mathrm{}}𝑑z|k_{zi}|.$$
(B6)
Evaluating this integral with the aid of equations (B2) and (B3), we obtain
$$\mathrm{ln}^1\left[\frac{1}{(\omega \tau _b)}\left(\frac{z_b}{z_\omega }\right)^{q+2}\right]^{\frac{1}{2(q+1)}}\frac{\mathrm{}^{(q+2)/(q+1)}}{\omega ^{(2q+5)/(2q+2)}}.$$
(B7)
Numerical results for $`\mathrm{ln}^1`$ plotted in the upper panel of Figure 10 confirm this relation.
Because the maximum value of $`\kappa `$ for parametric resonance is independent of the spherical degrees of the daughter modes, the dependence of $`\gamma `$ on $`\mathrm{}`$ at fixed $`\omega `$ is of great significance. Equations (B4) and (B7), together with the relation $`n\mathrm{}`$ at fixed $`\omega `$, establish that $`\gamma \mathrm{}^{1/(q+1)}`$. Since $`\gamma `$ increases with $`\mathrm{}`$ at fixed $`\omega `$, the most important daughter pairs are those with the smallest $`\mathrm{}`$ values subject to the constraint $`\gamma _d>\delta \omega `$.
### B.2. Reduction of the Coupling Coefficient by Strong Nonadiabaticity
The maximum adiabatic coupling coefficient between a parent mode and a pair of daughter modes is $`\kappa _{\mathrm{max}}(n_p^3\tau _{\omega _p}L)^{1/2}`$ (eq. ). Its major contribution comes from the region above $`z_{\omega _p}`$. The factor $`(n_p\tau _{\omega _p}L)^{1/2}`$ is the surface value of the normalized eigenfunction for $`[\mathrm{}(\mathrm{}+1)]^{1/2}\xi _p/R`$. The extra factor $`n_p^1`$ is the fraction of each daughter mode’s nodes that lie above $`z_{\omega _p}`$. The coupling coefficient is reduced compared to equation (12) for daughter modes that are strongly nonadiabatic in the region above $`z_{\omega _p}`$.
Equation (B1) indicates that nonadiabaticity increases with increasing $`\mathrm{}`$ at fixed $`z`$ and $`\omega `$; $`k_z(zz_\omega )^{1/2}`$ and $`z_\omega \omega ^2R^2/g\mathrm{}(\mathrm{}+1)`$, so $`k_z\mathrm{}`$. We define the angular degree of decoupling for a given parent mode, $`\mathrm{}_{\mathrm{dc}}`$, as the smallest spherical degree at which its daughter modes are strongly nonadiabatic all the way down to the top of the parent mode’s cavity; that is, $`z_{\mathrm{na}_\mathrm{d}}z_{\omega _p}`$ at $`\mathrm{}_{\mathrm{dc}}`$. Using equation (B2), we find
$$\mathrm{}_{\mathrm{dc}}\left[\left(\frac{\omega ^2R^2}{gz_b}\right)^{q+2}\frac{\omega \tau _b}{(\mathrm{}_p(\mathrm{}_p+1))^{q+1}}\right]^{1/2}.$$
(B8)
Numerical results displayed in the lower panel of Figure 10 are well-represented by this scaling relation. We find that $`\mathrm{}_{\mathrm{dc}}`$ is relatively independent of stellar effective temperature but decays steeply with mode period. By $`n_p=20`$, we find $`\mathrm{}_{\mathrm{dc}}1`$.
It is plausible that at $`\mathrm{}_d=\mathrm{}_{\mathrm{dc}}`$, $`\kappa _{\mathrm{max}}`$ is reduced by a factor $`n_p/n_d`$ below its adiabatic value (eq. ) because the effective lids of the daughter modes’ cavities are lowered to $`z_{\omega _p}`$. Since $`n_p/n_d\mathrm{}_p/2\mathrm{}_d`$, this is a large reduction for $`\mathrm{}_{\mathrm{dc}}\mathrm{}_p`$. An even more severe reduction of $`\kappa _{\mathrm{max}}`$ is expected for $`\mathrm{}_d>\mathrm{}_{\mathrm{dc}}`$.
### B.3. Parametric Instability for Traveling Waves
In the limit of strong nonadiabaticity, daughter modes are more appropriately described as traveling waves than as standing waves. Thus it behooves us to investigate the parametric instability of traveling waves. Here we demonstrate that the instability criterion for traveling waves is equivalent to that for standing waves (eq. ).
Nonlinear interactions between parent and daughter modes are localized within an interaction region above $`z_{\omega _p}`$. Let us assume that $`z_{\mathrm{na}_\mathrm{d}}z_{\omega _p}`$. Then, in most of the interaction region daughter wave packets may be represented as linear superpositions of adiabatic modes. Propagating at their group velocity, the daughter wave packets pass through the interaction region in a time interval
$$\mathrm{\Delta }T=_0^{z_{\omega _p}}\frac{dz}{v_{gz}}\frac{1}{n_p}\frac{\pi n_d}{\omega _d}.$$
(B9)
Three significant relations involving $`n_p`$ are worth noting: 1) $`n_p^1`$ is the fraction of each daughter mode’s nodes that lie above $`z_{\omega _p}`$, so $`2\mathrm{\Delta }T`$ is a fraction $`n_p^1`$ of the time each daughter wave packet takes to make a round trip across its cavity; 2) approximately $`n_p`$ daughter modes reside within the frequency interval $`\pi /\mathrm{\Delta }T`$, and their relative phases change by less than $`\pi `$ as each daughter wave packet crosses the interaction region; 3) maximal $`\kappa `$ occurs inside an interval of width $`|n_{d_1}n_{d_2}|n_p`$.
Nonlinear interactions between parent and daughter waves within the interaction region are described by equations (1)-(3) with two modifications of the equations governing the time evolution of the daughter modes. The linear damping term is negligible for $`zz_{\mathrm{na}}`$, and the nonlinear term must be multiplied by a factor $`n_p`$. The latter accounts for the number of modes which couple coherently to each daughter mode during the interaction time $`\mathrm{\Delta }T`$. During two passes through the interaction region, the amplitudes of the daughter wave packets grow by a factor $`e^G`$, where the gain, $`𝒢`$, is given by
$$𝒢=\frac{2\mathrm{\Delta }T}{|A_d|}\frac{d|A_d|}{dt}3\sqrt{2}n_d|\kappa ||A_p|.$$
(B10)
For parametric instability to occur,
$$𝒢>\mathrm{ln}^1.$$
(B11)
Combining the relation between $``$ and $`\gamma `$ given by equation (B4) with equation (B11), the threshold condition for parametric instability of traveling waves becomes
$$|A_p|>\frac{\gamma _d}{3\sqrt{2}\omega _d|\kappa |}.$$
(B12)
The above condition is equivalent to equation (4) in the limit that $`\gamma _d\delta \omega `$. Thus the threshold condition for parametric instability of traveling waves reduces to a limiting case of the threshold condition for parametric instability of standing waves.
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# Kaluza-Klein Formalism of General Spacetimes
## I Introduction
It has been known for some time that there is a curious correspondence between (self-dual) Yang-Mills equations and the (self-dual) Einstein’s equations, when the Yang-Mills gauge symmetry is extended to an infinite dimensional symmetry of (volume-preserving) diffeomorphisms of some auxiliary manifold. It is also well-known that the equations of motion of 2-dimensional non-linear sigma models with the target space as the area-preserving diffeomorphism of an auxiliary 2-surface are identical to the the self-dual Einstein’s equations written in the Plebañski form.
These correspondences are most striking for self-dual cases, and indicate an intriguing possibility that we may be able to reconstruct the full Einstein’s general relativity from suitable gauge field theories by replacing the usual finite dimensional gauge symmetry with an infinite dimensional group of the diffeomorphisms of some manifold. If we recall that the gauge symmetry of general relativity is the group of the diffeomorphisms of a 4-dimensional spacetime, this seemingly wild speculation is not totally unreasonable. Recently we have shown that such a description is indeed possible, by rewriting the Einstein-Hilbert action of general relativity of generic 4-dimensional spacetimes in the (2,2)-decomposition . In this approach, the 4-dimensional spacetime is viewed, at least for a finite range of the spacetime, as a locally fibred manifold that consists of a (1+1)-dimensional base manifold $`M_{1+1}`$ and a 2-dimensional fibre space $`N_2`$.
The Yang-Mills gauge fields, which naturally appear in this Kaluza-Klein setting, are defined on the (1+1)-dimensional base manifold $`M_{1+1}`$, and turn out to be valued in the Lie algebra of an infinite dimensional group of the diffeomorphisms of the 2-dimensional fibre space $`N_2`$ (i.e. diff$`N_2`$). This feature is expected to simplify considerably certain issues concerned with the constraints of general relativity. Namely, in Yang-Mills gauge theories, it is well-known that the Gauss-law constraints associated with the Yang-Mills gauge invariance can be made “trivial”, if we consider gauge invariant quantities only. Thus, in principle, one might expect that the problem of solving the constraints of general relativity could be made “trivial”, at least for some of them, if such a gauge theory description is possible. The purpose of this paper is to show explicitly that our variables transform as a tensor field, gauge fields, and scalar fields with respect to the diff$`N_2`$ transformations, and discuss a general spacetime from the 4-dimensional fibre bundle point of view.
This paper is organized as follows. In section II, we shall outline the kinematics of the (2,2)-decomposition of a generic 4-dimensional spacetime, and introduce the Kaluza-Klein (KK) variables without assuming any spacetime isometries. In section III, we shall find the transformation properties of the KK variables with respect to the diff$`N_2`$ transformations, and introduce the notion of the diff$`N_2`$-covariant derivatives. In section IV, we shall write down the Einstein-Hilbert action, and finally, we discuss possible applications of this formalism.
## II Kinematics
Let us decompose a generic 4-dimensional spacetime of the Lorentzian signature from the KK perspective, in which the spacetime under consideration is viewed as a 4-dimensional fibre bundle, consisting of a (1+1)-dimensional base manifold $`M_{1+1}`$ and a 2-dimensional fibre space $`N_2`$. Let the basis vector fields of $`M_{1+1}`$ and $`N_2`$ be $`/x^\mu (=_\mu )`$ and $`/y^a(=_a)`$, respectively, where $`\mu =0,1`$ and $`a=2,3`$. The horizontal vector fields $`\widehat{}_\mu `$, which are defined to be orthogonal to $`N_2`$, can be expressed as linear combinations of $`_\mu `$ and $`_a`$,
$$\widehat{}_\mu =_\mu A_\mu ^a_a,$$
(1)
where the fields $`A_\mu ^a`$ are functions of $`(x^\mu ,y^a)`$. Let us denote by $`\gamma ^{\mu \nu }`$ the inverse metric of the horizontal space spanned by $`\widehat{}_\mu `$, and by $`\varphi ^{ab}`$ the inverse metric of $`N_2`$, respectively. In the horizontal lift basis which consists of $`\{\widehat{}_\mu ,_a\}`$, the metric of the 4-dimensional spacetime can then be written as
$$\left(\frac{}{s}\right)^2=\gamma ^{\mu \nu }\left(_\mu A_\mu ^a_a\right)\left(_\nu A_\nu ^b_b\right)+\varphi ^{ab}_a_b.$$
(2)
In the corresponding dual basis $`\{dx^\mu ,dy^a+A_\mu ^adx^\mu \}`$, the metric becomes
$$ds^2=\gamma _{\mu \nu }dx^\mu dx^\nu +\varphi _{ab}\left(dy^a+A_\mu ^adx^\mu \right)\left(dy^b+A_\nu ^bdx^\nu \right).$$
(3)
Formally the above metric looks similar to the “dimensionally reduced” metric in standard KK theories, but in fact it is quite different. In the standard KK reduction certain isometries are usually assumed, and dimensional reduction is made by projection along the directions generated by these isometries. There the fields $`A_\mu ^a`$ are identified as the KK gauge fields associated with the finite dimensional isometry group. In this paper, we do not assume such isometries: nevertheless, it turns out that the KK idea still works, and as we shall show shortly, the fields $`A_\mu ^a`$ can be identified as the gauge fields valued in the infinite dimensional Lie algebra of the diff$`N_2`$ transformations. Moreover, the fields $`\varphi _{ab}`$ and $`\gamma _{\mu \nu }`$ transform as a tensor field and scalar fields with respect to the diff$`N_2`$ transformations.
## III Diffeomorphisms as a local gauge symmetry
### A Finite transformations
Let us find the transformation properties of the fields $`\varphi _{ab}`$, $`A_\mu ^a`$, and $`\gamma _{\mu \nu }`$ with respect to the diff$`N_2`$ transformations, which are the following coordinate transformations of $`N_2`$, while keeping $`x^\mu `$ constant,
$$y^{}_{}{}^{}a=y^{}_{}{}^{}a(x,y),x^{}_{}{}^{}\mu =x^\mu .$$
(4)
Thus we have
$$dy^a=\frac{y^a}{y^{}_{}{}^{}c}\left\{dy^{}_{}{}^{}c\left(\frac{y^{}_{}{}^{}c}{x^\mu }\right)dx^{}_{}{}^{}\mu \right\},dx^\mu =dx^{}_{}{}^{}\mu .$$
(5)
In the new coordinates the term proportional to $`dx^\mu dy^a`$ in (3) becomes, keeping the $`(x^\mu ,y^a)`$ dependence explicit,
$`2\varphi _{ab}(x,y)A_\mu ^a(x,y)dx^\mu dy^b`$ (6)
$`=`$ $`2\left({\displaystyle \frac{y^a}{y^{}_{}{}^{}c}}\right)\left({\displaystyle \frac{y^b}{y^{}_{}{}^{}d}}\right)\varphi _{ab}(x,y)\left({\displaystyle \frac{y^{}_{}{}^{}d}{y^e}}\right)A_\mu ^e(x,y)dx^{}_{}{}^{}\mu \left\{dy^{}_{}{}^{}c\left({\displaystyle \frac{y^{}_{}{}^{}c}{x^\nu }}\right)dx^{}_{}{}^{}\nu \right\},`$ (7)
where the identity
$$\left(\frac{y^a}{y^{}_{}{}^{}d}\right)\left(\frac{y^{}_{}{}^{}d}{y^e}\right)=\delta _e^a$$
(8)
was used. Also the term proportional to $`dy^ady^b`$ becomes
$`\varphi _{ab}(x,y)dy^ady^b`$ (9)
$`=`$ $`\left({\displaystyle \frac{y^a}{y^{}_{}{}^{}c}}\right)\left({\displaystyle \frac{y^b}{y^{}_{}{}^{}d}}\right)\varphi _{ab}(x,y)\left\{dy^{}_{}{}^{}cdy^{}_{}{}^{}d2\left({\displaystyle \frac{y^{}_{}{}^{}d}{x^\mu }}\right)dy^{}_{}{}^{}cdx^{}_{}{}^{}\mu +\left({\displaystyle \frac{y^{}_{}{}^{}c}{x^\mu }}\right)\left({\displaystyle \frac{y^{}_{}{}^{}d}{x^\nu }}\right)dx^{}_{}{}^{}\mu dx^{}_{}{}^{}\nu \right\}.`$ (10)
After rearranging terms, the metric (3) can be written as, in the new coordinates,
$`ds^2`$ $`=`$ $`\gamma _{\mu \nu }(x,y)dx^{}_{}{}^{}\mu dx^{}_{}{}^{}\nu +\left({\displaystyle \frac{y^a}{y^{}_{}{}^{}c}}\right)\left({\displaystyle \frac{y^b}{y^{}_{}{}^{}d}}\right)\varphi _{ab}(x,y)dy^{}_{}{}^{}cdy^{}_{}{}^{}d`$ (14)
$`+2\left({\displaystyle \frac{y^a}{y^{}_{}{}^{}c}}\right)\left({\displaystyle \frac{y^b}{y^{}_{}{}^{}d}}\right)\varphi _{ab}(x,y)\left\{\left({\displaystyle \frac{y^{}_{}{}^{}d}{y^e}}\right)A_\mu ^e(x,y){\displaystyle \frac{y^{}_{}{}^{}d}{x^\mu }}\right\}dx^{}_{}{}^{}\mu dy^{}_{}{}^{}c`$
$`+\varphi _{ab}(x,y)\{A_\mu ^a(x,y)A_\nu ^b(x,y)2\left({\displaystyle \frac{y^a}{y^{}_{}{}^{}c}}\right)\left({\displaystyle \frac{y^b}{y^{}_{}{}^{}d}}\right)\left({\displaystyle \frac{y^{}_{}{}^{}d}{y^e}}\right)A_\mu ^e(x,y)\left({\displaystyle \frac{y^{}_{}{}^{}c}{x^\nu }}\right)`$
$`+\left({\displaystyle \frac{y^a}{y^{}_{}{}^{}c}}\right)\left({\displaystyle \frac{y^b}{y^{}_{}{}^{}d}}\right)\left({\displaystyle \frac{y^{}_{}{}^{}c}{x^\mu }}\right)\left({\displaystyle \frac{y^{}_{}{}^{}d}{x^\nu }}\right)\}dx^{}_{}{}^{}\mu dx^{}_{}{}^{}\nu ,`$
which must be equal to
$$ds^{}_{}{}^{}2=\gamma _{\mu \nu }^{}(x^{},y^{})dx^{}_{}{}^{}\mu dx^{}_{}{}^{}\nu +\varphi _{ab}^{}(x^{},y^{})\left\{dy^{}_{}{}^{}a+A_\mu ^{}_{}{}^{}a(x^{},y^{})dx^{}_{}{}^{}\mu \right\}\left\{dy^{}_{}{}^{}b+A_\nu ^{}_{}{}^{}b(x^{},y^{})dx^{}_{}{}^{}\nu \right\},$$
(15)
since the line element is invariant under the diff$`N_2`$ transformations. If we compare terms containing $`dy^{}_{}{}^{}ady^{}_{}{}^{}b`$, we find that $`\varphi _{ab}(x,y)`$ transform as
$$\varphi _{ab}^{}(x^{},y^{})=\left(\frac{y^c}{y^{}_{}{}^{}a}\right)\left(\frac{y^d}{y^{}_{}{}^{}b}\right)\varphi _{cd}(x,y).$$
(16)
This shows that $`\varphi _{ab}(x,y)`$ is a tensor field with respect to the diff$`N_2`$ transformations. If we use the equation (16) in (14), the metric becomes
$`ds^2`$ $`=`$ $`\gamma _{\mu \nu }(x,y)dx^{}_{}{}^{}\mu dx^{}_{}{}^{}\nu +\varphi _{cd}^{}(x^{},y^{})dy^{}_{}{}^{}cdy^{}_{}{}^{}d+2\varphi _{cd}^{}(x^{},y^{})\left\{\left({\displaystyle \frac{y^{}_{}{}^{}d}{y^a}}\right)A_\mu ^a(x,y){\displaystyle \frac{y^{}_{}{}^{}d}{x^\mu }}\right\}dx^{}_{}{}^{}\mu dy^{}_{}{}^{}c`$ (18)
$`+\varphi _{cd}^{}(x^{},y^{})\left\{\left({\displaystyle \frac{y^{}_{}{}^{}c}{y^a}}\right)A_\mu ^a(x,y){\displaystyle \frac{y^{}_{}{}^{}c}{x^\mu }}\right\}\left\{\left({\displaystyle \frac{y^{}_{}{}^{}d}{y^b}}\right)A_\nu ^b(x,y){\displaystyle \frac{y^{}_{}{}^{}d}{x^\nu }}\right\}dx^{}_{}{}^{}\mu dx^{}_{}{}^{}\nu ,`$
from which we deduce the following transformation properties of $`A_\mu ^a(x,y)`$ and $`\gamma _{\mu \nu }(x,y)`$
$`A_\mu ^{}_{}{}^{}a(x^{},y^{})=\left({\displaystyle \frac{y^{}_{}{}^{}a}{y^b}}\right)A_\mu ^b(x,y){\displaystyle \frac{y^{}_{}{}^{}a}{x^\mu }}(x,y),`$ (19)
$`\gamma _{\mu \nu }^{}(x^{},y^{})=\gamma _{\mu \nu }(x,y),`$ (20)
under the diff$`N_2`$ transformations.
### B Infinitesimal transformations
It will be instructive to examine the infinitesimal transformations corresponding to the above finite diff$`N_2`$ transformations. The infinitesimal diff$`N_2`$ transformations consist of the following transformations
$$y^{}_{}{}^{}a=y^a+\xi ^a(x,y),x^{}_{}{}^{}\mu =x^\mu (\mathrm{O}(\xi ^2)1),$$
(21)
where $`\xi ^a(x,y)`$ is an arbitrary, infinitesimal, function of $`(x^\mu ,y^a)`$. From this it follows that
$$\frac{y^c}{y^{}_{}{}^{}a}=\delta _a^c\frac{\xi ^c}{y^a}+\mathrm{},$$
(22)
where $`\mathrm{}`$ means terms of $`\mathrm{O}(\xi ^2)`$. If we expand the l.h.s. of the equation (16) in $`\xi ^a`$, it becomes
$$\varphi _{ab}^{}(x^{},y+\xi )=\varphi _{ab}^{}(x,y)+\xi ^c\frac{}{y^c}\varphi _{ab}(x,y)+\mathrm{},$$
(23)
whereas the r.h.s. becomes
$`\left({\displaystyle \frac{y^c}{y^{}_{}{}^{}a}}\right)\left({\displaystyle \frac{y^d}{y^{}_{}{}^{}b}}\right)\varphi _{cd}(x,y)`$ $`=`$ $`\varphi _{ab}(x,y){\displaystyle \frac{\xi ^c}{y^a}}\varphi _{cb}(x,y){\displaystyle \frac{\xi ^c}{y^b}}\varphi _{ac}(x,y)+\mathrm{}.`$ (24)
Thus we have
$`\delta \varphi _{ab}(x,y)`$ $``$ $`\varphi _{ab}^{}(x,y)\varphi _{ab}(x,y)`$ (25)
$`=`$ $`\xi ^c_c\varphi _{ab}(x,y)(_a\xi ^c)\varphi _{cb}(x,y)(_b\xi ^c)\varphi _{ac}(x,y)`$ (26)
$`=`$ $`[\xi ,\varphi ]_{\mathrm{L}ab},`$ (27)
where the subscript <sub>L</sub> denotes the Lie derivative along the vector field $`\xi \xi ^a_a`$, i.e.
$$[\xi ,\varphi ]_{\mathrm{L}ab}=\xi ^c_c\varphi _{ab}+(_a\xi ^c)\varphi _{cb}+(_b\xi ^c)\varphi _{ac}.$$
(28)
It is a straightforward exercise to derive the infinitesimal transformation properties $`A_\mu ^a`$ and $`\gamma _{\mu \nu }`$ from (19) and (20). They are found to be
$`\delta A_\mu ^a(x,y)`$ $`=`$ $`_\mu \xi ^a+[A_\mu ,\xi ]_\mathrm{L}^a`$ (29)
$`=`$ $`_\mu \xi ^a+A_\mu ^c_c\xi ^a\xi ^c_cA_\mu ^a,`$ (30)
$`\delta \gamma _{\mu \nu }(x,y)`$ $`=`$ $`[\xi ,\gamma _{\mu \nu }]_\mathrm{L}`$ (31)
$`=`$ $`\xi ^a_a\gamma _{\mu \nu },`$ (32)
where $`[A_\mu ,\xi ]_\mathrm{L}^a`$ and $`[\xi ,\gamma _{\mu \nu }]_\mathrm{L}`$ are the Lie derivatives of $`\xi ^a`$ and $`\gamma _{\mu \nu }`$ along the vector fields $`A_\mu =A_\mu ^c_c`$ and $`\xi =\xi ^c_c`$, respectively. Notice that the Lie derivative acts on the fibre space index ($`a`$) only. The equations (27), (30), and (32) clearly show that the metric components $`\{\varphi _{ab},A_\mu ^a,\gamma _{\mu \nu }\}`$ transform as a tensor field, gauge fields, and scalar fields under the diff$`N_2`$ transformations, respectively.
### C diff$`N_2`$-covariant derivative
Using the Lie derivative along the diff$`N_2`$-valued gauge fields, the diff$`N_2`$-covariant derivative $`D_\mu `$ can be naturally defined as
$$D_\mu =_\mu [A_\mu ,]_\mathrm{L}.$$
(33)
With this definition, the equation (30) can be written as
$$\delta A_\mu ^a=D_\mu \xi ^a,$$
(34)
which suggests that the diff$`N_2`$-valued field strength $`F_{\mu \nu }^a`$ be defined as
$$[D_\mu ,D_\nu ]_\mathrm{L}\eta =F_{\mu \nu }^a_a\eta $$
(35)
for an arbitrary scalar function $`\eta `$, where $`F_{\mu \nu }^a`$ is given by
$`F_{\mu \nu }^a`$ $`=`$ $`_\mu A_\nu ^a_\nu A_\mu ^a[A_\mu ,A_\nu ]_\mathrm{L}^a`$ (36)
$`=`$ $`_\mu A_\nu ^a_\nu A_\mu ^aA_\mu ^c_cA_\nu ^a+A_\nu ^c_cA_\mu ^a.`$ (37)
Similarly, the diff$`N_2`$-covariant derivative of $`\varphi _{ab}`$ is defined as
$`D_\mu \varphi _{ab}`$ $`=`$ $`_\mu \varphi _{ab}[A_\mu ,\varphi ]_{\mathrm{L}ab}`$ (38)
$`=`$ $`_\mu \varphi _{ab}A_\mu ^c_c\varphi _{ab}(_aA_\mu ^c)\varphi _{bc}(_bA_\mu ^c)\varphi _{ac}.`$ (39)
It remains to show that $`F_{\mu \nu }^a`$ and $`D_\mu \varphi _{ab}`$ transform covariantly under the infinitesimal diff$`N_2`$ transformations (21). Let us consider $`D_\mu \varphi _{ab}`$ first. The infinitesimal transformation of $`D_\mu \varphi _{ab}`$ becomes
$$\delta (D_\mu \varphi _{ab})=_\mu \left([\xi ,\varphi ]_{\mathrm{L}ab}\right)+[A_\mu ,[\xi ,\varphi ]_\mathrm{L}]_{\mathrm{L}ab}+[D_\mu \xi ,\varphi ]_{\mathrm{L}ab},$$
(40)
where we used the equations (27) and (30), and the Lie brackets are
$`[A_\mu ,[\xi ,\varphi ]_\mathrm{L}]_{\mathrm{L}ab}=A_\mu ^c_c\left([\xi ,\varphi ]_{\mathrm{L}ab}\right)+(_aA_\mu ^c)[\xi ,\varphi ]_{\mathrm{L}bc}+(_bA_\mu ^c)[\xi ,\varphi ]_{\mathrm{L}ac},`$ (41)
$`[D_\mu \xi ,\varphi ]_{\mathrm{L}ab}=(D_\mu \xi ^c)(_c\varphi _{ab})+_a(D_\mu \xi ^c)\varphi _{bc}+_b(D_\mu \xi ^c)\varphi _{ac}.`$ (42)
Using the Leibniz rule of the derivative $`_\mu `$
$$_\mu \left([\xi ,\varphi ]_{\mathrm{L}ab}\right)=[_\mu \xi ,\varphi ]_{\mathrm{L}ab}+[\xi ,_\mu \varphi ]_{\mathrm{L}ab},$$
(43)
and the properties of the Lie bracket
$`[D_\mu \xi ,\varphi ]_{\mathrm{L}ab}=[_\mu \xi ,\varphi ]_{\mathrm{L}ab}[[A_\mu ,\xi ]_\mathrm{L},\varphi ]_{\mathrm{L}ab},`$ (44)
$`[A_\mu ,[\xi ,\varphi ]_\mathrm{L}]_{\mathrm{L}ab}=[\xi ,[\varphi ,A_\mu ]_\mathrm{L}]_{\mathrm{L}ab}[\varphi ,[A_\mu ,\xi ]_\mathrm{L}]_{\mathrm{L}ab},`$ (45)
we find that the equation (40) becomes
$`\delta (D_\mu \varphi _{ab})`$ $`=`$ $`[\xi ,_\mu \varphi ]_{\mathrm{L}ab}+[\xi ,[A_\mu ,\varphi ]_\mathrm{L}]_{\mathrm{L}ab}`$ (46)
$`=`$ $`[\xi ,D_\mu \varphi ]_{\mathrm{L}ab},`$ (47)
which shows that $`D_\mu \varphi _{ab}`$ transforms covariantly under the diff$`N_2`$ transformation.
Similarly, the infinitesimal transformation $`\delta F_{\mu \nu }^a`$ becomes
$$\delta F_{\mu \nu }^a=_\mu \left([A_\nu ,\xi ]_\mathrm{L}^a\right)+[D_\mu \xi ,A_\nu ]_\mathrm{L}^a(\mu \nu ).$$
(48)
Using the following identities
$`_\mu \left([A_\nu ,\xi ]_\mathrm{L}^a\right)=[_\mu A_\nu ,\xi ]_\mathrm{L}^a+[A_\nu ,_\mu \xi ]_\mathrm{L}^a,`$ (49)
$`[D_\mu \xi ,A_\nu ]_\mathrm{L}^a=[A_\nu ,D_\mu \xi ]_\mathrm{L}^a=[A_\nu ,_\mu \xi ]_\mathrm{L}^a+[A_\nu ,[A_\mu ,\xi ]_\mathrm{L}]_\mathrm{L}^a,`$ (50)
we find that
$`\delta F_{\mu \nu }^a`$ $`=`$ $`[_\mu A_\nu _\nu A_\mu ,\xi ]_\mathrm{L}^a+[A_\nu ,[A_\mu ,\xi ]_\mathrm{L}]_\mathrm{L}^a[A_\mu ,[A_\nu ,\xi ]_\mathrm{L}]_\mathrm{L}^a`$ (51)
$`=`$ $`[\xi ,F_{\mu \nu }]_\mathrm{L}^a,`$ (52)
where we used the Jacobi identity
$$[A_\nu ,[A_\mu ,\xi ]_\mathrm{L}]_\mathrm{L}^a=[A_\mu ,[\xi ,A_\nu ]_\mathrm{L}]_\mathrm{L}^a[\xi ,[A_\nu ,A_\mu ]_\mathrm{L}]_\mathrm{L}^a.$$
(53)
Therefore it follows that
$$\delta F_{\mu \nu }^a=[\xi ,F_{\mu \nu }]_\mathrm{L}^a,$$
(54)
which shows that $`F_{\mu \nu }^a`$ is indeed the diff$`N_2`$-valued field strength.
It must be marked here that, in the (2,2)-KK formalism, the Lie derivative, rather than the covariant derivative, appears naturally. The appearance of an infinite dimensional symmetry such as diff$`N_2`$ is not surprising, since in general relativity the underlying gauge symmetry is the infinite dimensional group of the diffeomorphisms of a 4-dimensional spacetime. The point is that it is the diff$`N_2`$ symmetry, the subgroup of the diffeomorphisms of a 4-dimensional spacetime, that shows up as a local gauge symmetry of the Yang-Mills type. This implies that the (2,2)-KK formalism can be made a viable method of studying general relativity from the standpoint of the (1+1)-dimensional Yang-Mills gauge theory with the diff$`N_2`$ symmetry as a local gauge symmetry.
## IV The Action
The Einstein-Hilbert action in this KK formalism is given by
$`I`$ $`=`$ $`{\displaystyle }d^2xd^2y\sqrt{\gamma }\sqrt{\varphi }[\gamma ^{\mu \nu }\widehat{\mathrm{R}}_{\mu \nu }+\varphi ^{ac}\mathrm{R}_{ac}+{\displaystyle \frac{1}{4}}\gamma ^{\mu \nu }\gamma ^{\alpha \beta }\varphi _{ab}F_{\mu \alpha }^aF_{\nu \beta }^b`$ (57)
$`+{\displaystyle \frac{1}{4}}\gamma ^{\mu \nu }\varphi ^{ab}\varphi ^{cd}\left\{(D_\mu \varphi _{ac})(D_\nu \varphi _{bd})(D_\mu \varphi _{ab})(D_\nu \varphi _{cd})\right\}`$
$`+{\displaystyle \frac{1}{4}}\varphi ^{ab}\gamma ^{\mu \nu }\gamma ^{\alpha \beta }\{(_a\gamma _{\mu \alpha })(_b\gamma _{\nu \beta })(_a\gamma _{\mu \nu })(_b\gamma _{\alpha \beta })\}]+{\displaystyle }d^2xd^2y(_AS^A).`$
Let us summarize the notations:
1. The curvature tensors $`\widehat{\mathrm{R}}_{\mu \nu }`$ and $`\mathrm{R}_{ac}`$ are defined as
$`\widehat{\mathrm{R}}_{\mu \nu }=\widehat{}_\mu \widehat{\mathrm{\Gamma }}_{\alpha \nu }^\alpha \widehat{}_\alpha \widehat{\mathrm{\Gamma }}_{\mu \nu }^\alpha +\widehat{\mathrm{\Gamma }}_{\mu \beta }^\alpha \widehat{\mathrm{\Gamma }}_{\alpha \nu }^\beta \widehat{\mathrm{\Gamma }}_{\beta \alpha }^\beta \widehat{\mathrm{\Gamma }}_{\mu \nu }^\alpha ,`$ (58)
$`\mathrm{R}_{ac}=_a\mathrm{\Gamma }_{bc}^b_b\mathrm{\Gamma }_{ac}^b+\mathrm{\Gamma }_{ad}^b\mathrm{\Gamma }_{bc}^d\mathrm{\Gamma }_{db}^d\mathrm{\Gamma }_{ac}^b,`$ (59)
$`\widehat{\mathrm{\Gamma }}_{\mu \nu }^\alpha ={\displaystyle \frac{1}{2}}\gamma ^{\alpha \beta }\left(\widehat{}_\mu \gamma _{\nu \beta }+\widehat{}_\nu \gamma _{\mu \beta }\widehat{}_\beta \gamma _{\mu \nu }\right),`$ (60)
$`\mathrm{\Gamma }_{ab}^c={\displaystyle \frac{1}{2}}\varphi ^{cd}\left(_a\varphi _{bd}+_b\varphi _{ad}_d\varphi _{ab}\right).`$ (61)
2. The last term in (57) is a surface integral, where $`S^A=(S^\mu ,S^a)`$ is given by
$`S^\mu =\sqrt{\gamma }\sqrt{\varphi }j^\mu ,`$ (62)
$`S^a=\sqrt{\gamma }\sqrt{\varphi }\left(A_\mu ^aj^\mu +j^a\right),`$ (63)
$`j^\mu =\gamma ^{\mu \nu }\varphi ^{ab}D_\nu \varphi _{ab},`$ (64)
$`j^a=\varphi ^{ab}\gamma ^{\mu \nu }_b\gamma _{\mu \nu }.`$ (65)
One can easily recognize that this action is in a form of a (1+1)-dimensional field theory action. In geometrical terms the above action can be understood as follows. $`\gamma ^{\mu \nu }\widehat{\mathrm{R}}_{\mu \nu }`$ can be interpreted as the “gauged” scalar curvature of $`M_{1+1}`$, since the diff$`N_2`$-valued gauge fields are coupled to $`\gamma _{\mu \nu }`$ and $`\widehat{\mathrm{\Gamma }}_{\mu \nu }^\alpha `$ in the formulae (58) and (60). $`\varphi ^{ac}\mathrm{R}_{ac}`$ is the scalar curvature of $`N_2`$, which is proportional to the Euler characteristics $`\chi `$ when integrated over $`N_2`$.
The remaining terms in (57) are the extrinsic terms, telling us how $`M_{1+1}`$ and $`N_2`$ are embedded into the enveloping 4-dimensional spacetime. Each term in (57) is manifestly diff$`N_2`$-invariant, and the $`y^a`$-dependence of each term is completely “hidden” in the Lie derivatives. In this sense we may view the fibre space $`N_2`$ as a kind of “internal” space as in Yang-Mills theory. Thus, the Einstein-Hilbert action is describable as a (1+1)-dimensional Yang-Mills type gauge theory interacting with scalar fields and (1+1)-dimensional non-linear sigma fields of generic types, with couplings to curvatures of two 2-surfaces. The associated Yang-Mills gauge symmetry is the diff$`N_2`$ symmetry.
## V Discussions
In this paper, we presented the KK formalism of general relativity of generic 4-dimensional spacetimes, viewing the spacetime as a local product of the (1+1)-dimensional base manifold and the 2-dimensional fibre space. Within this framework, we made a decomposition of a given 4-dimensional spacetime metric into sets of fields which transform as a tensor field, gauge fields, and scalar fields under the group of the diffeomorphisms on $`N_2`$.
In connection with issues of quantum gravity, this KK approach has the following aspects which deserve further remarks. For instance, solving the Einstein’s constraint equations or constructing the gauge invariant physical observables is known to be one of the most important problems in quantum general relativity. In our formalism, the diffeomorphisms of the 2-dimensional space $`N_2`$ plays the role of a local gauge symmetry exactly as in Yang-Mills theory. Therefore the two constraint equations associated with the diff$`N_2`$ transformations can be “automatically” solved, using the diff$`N_2`$-invariant scalars. However, there are two additional constraint equations which require further studies in order to fully take care of the four Einstein’s constraint equations.
It should be also stressed that the Lie derivative appears naturally in this formalism, via the minimal couplings to the diff$`N_2`$-valued gauge fields. In the standard (3+1)-formalism, the natural derivative operator is the metric-compatible covariant derivative, which requires the metric be non-degenerate. The Lie derivative, on the other hand, can be defined even when the metric is degenerate. For instance, at null infinity $`^+`$ of the asymptotically flat spacetimes, the natural derivative operator is the Lie derivative, rather than the covariant derivative, because the metric on $`^+`$ is degenerate with the signature $`(0,+,+)`$. Therefore, the KK formalism, based on the notion of the Lie derivative, should be extendable to spacetimes where the metric is degenerate, which would be difficult to describe in conventional approaches.
Finally, it will be a challenging problem to try to reinterpret the exact solutions of the Einstein’s equations from this gauge theory point of view. This seems very interesting, for there are a number of exact solutions of the Einstein’s equations which do not permit sensible physical interpretations from the 4-dimensional spacetime perspective.
Acknowledgments
The author thanks the referee for informing him the related work of L.J. Mason and E.T. Newman . This work is supported in part by Korea Science and Engineering Foundation (95-0702-04-01-3).
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# Contents
## 1 Introduction
Our present understanding of QCD is based on the widely accepted idea that the confining regime of Yang-Mills theories should be described by some kind of effective string model .
This conjecture has by now a very long history. It originates from two independent observations.
* The first one is of phenomenological nature, and predates the formulation of QCD. It is related to the observation that the linearly rising Regge trajectories in meson spectroscopy can be easily explained assuming a string-type interaction between the quark and the antiquark. This observation was at the origin of a large amount of papers which tried to give a consistent quantum description of strings.
* The second one comes from the lattice regularization of pure gauge theories (LGT in the following) and was realized right after the formulation of QCD. In the LGT framework one can easily study non-perturbative phenomena, like those involved in the conjectured string description of $`YM`$ theories and it is easy to see that in the strong coupling limit of pure LGTs the interquark potential rises linearly, and that the chromoelectric flux lines are confined in a thin “string like” flux tube .
Some clear indications were later found that the vacuum expectation value of Wilson loops could be rewritten as a string functional integral even in the continuum . This led to conjecture that there exists an exact duality between gauge fields and strings .
However, despite these results, in the following years all the attempts to explicitly construct the conjectured string description of QCD failed. In fact it was realized at the beginning of the eighties that the strong coupling approximation for the lattice description of interquark potential is plagued by lattice artifacts which make it inadequate for the continuum theory (this is the famous “roughening transition” that we shall discuss in sect. 3.5.1).
Since then, while impressive results were obtained by means of Montecarlo simulations, very little progress have been achieved with analytic techniques. Lattice gauge theories can be solved exactly in two dimensions for any gauge group, but become unaffordably complex in more than two dimensions, even in absence of quarks. Moreover, most of the approximation techniques which are usually successful in dealing with simpler statistical mechanical systems, like (suitably improved) mean field methods or strong coupling expansions turn out to be less useful in the case of LGT.
Several proposals were made during the eighties to overcome these difficulties. In particular, two of them led to rather interesting results.
Effective string theory
The first proposal was to assume a milder version of the conjecture. This milder version only requires that the behavior of large Wilson loops is described in the infrared limit by an effective two-dimensional field theory (hence not a true string theory) which accounts for the string-like properties of long chromoelectric flux tubes. We shall refer to such a $`2d`$ field theory in the following as the “effective string theory”. Several interesting results can be obtained in this framework. We shall discuss them in detail in sect. 3.5 below. Let us anticipate here that the most interesting feature of this approach is that it leads to predictions which are in very good agreement with the results of Montecarlo simulations of QCD. Its major drawback is that it is not consistent at the quantum level. It is not clear how to extend this effective string description to the ultraviolet regime, i.e. how to relate it with some kind of “fundamental” string which is consistent at the quantum level.
Large $`N`$ limit
The second proposal was to study the large $`N`$ limit of $`SU(N)`$ gauge theories instead of the phenomenologically relevant $`SU(3)`$ model . It was shown that this large $`N`$ approximation is able to keep the whole complexity of the finite $`N`$ models. Unfortunately, even in the large $`N`$ limit (despite the fact that some major simplifications occur) it is not possible to give exact solution (the so called “Master Field”) to the Lattice SU$`(N)`$ model and essentially no improvement was made in these last fifteen years also in this direction.
Recently this situation drastically changed thanks to a new, original proposal based on the Maldacena conjecture which relates the large $`N`$ expansion of certain supersymmetric gauge theories to the behaviour of string theory in a non-trivial geometry. Witten’s extension of this conjecture to non-supersymmetric gauge theories, led to the hope of a possible non-perturbative description also for large $`N`$ QCD in four dimensions. In fact in these last months several attempts have been made to extract predictions for the string tension and the glueball spectrum of large $`N`$ QCD. These predictions have some appealing features, but also raised serious criticism. All the authors agree that some independent test of the applicability of the AdS/CFT correspondence to large $`N`$ QCD is needed. This is indeed possible thanks to the impressive progress of montecarlo simulations of LGT which have by now reached stable estimates both of string tension and glueball spectrum for finite $`N`$ and allow reliable extrapolations to the large $`N`$ limit.
The aim of these lectures is to allow the reader (which is assumed $`not`$ to be an expert of LGT) to understand how these estimates were obtained and to test their reliability. We shall also compare the “effective string model” mentioned above with the string theory which is at the basis of AdS/CFT model. The goal is to be able not only to accept or reject the AdS/CFT predictions on the basis of the LGT results but also, if possible, to gain some insight in the AdS/CFT proposal itself.
To this end we shall devote the first part of these lectures (sect. 2 and the first part of sect. 3) to an elementary introduction to the lattice regularization of QCD, starting form the very beginning. Then in the second part of the lectures (from sect. 3.4 to 3.7) we shall jump to our main object of interest and study in some detail the LGT results for the string tension and the glueball spectrum.
Unfortunately we shall have to skip several important and interesting hot topics of LGT, like the issue of improved (and “perfect”) actions, that of chiral fermions, topological observables, deconfinement transition …. We leave the interested reader to the books and review articles listed at the end of the bibliography - (we tried to make the list as complete as possible) which summarize the present state of the art in LGT.
The last part of these lectures (sect.s 4 and 5) is devoted to the AdS/CFT correspondence and in particular to discuss its predictions for the non-perturbative behaviour of non-supersymmetric Yang-Mills theories. Let us stress in this respect that this set of lectures is not intended as an introduction to the AdS/CFT correspondence for which we refer to the other lectures delivered at this school and to two thorough reviews which already exist on the subject ,. In this lectures we shall assume that the reader is already acquainted with the topic and shall only remind some basic informations at the beginning of sect. 4.
## 2 Introduction to Lattice Gauge Theories.
### 2.1 Quantum Field Theories and Statistical Mechanics.
The modern approach to Quantum Field Theories (QFT in the following) is based on Feynman’s path integral formulation. Using path integrals impressive results have been obtained in the last fifty years in perturbative QFT However these methods require the existence of one (or more) weak coupling parameters in which the theory can be expanded perturbatively. As such they are not suited for the analysis of phenomena governed by intrinsically large coupling constants, or even worse, with a non-analytic behaviour at the origin, in the space of complex coupling parameters. This is exactly the case of QCD, at least as far as dimensional observables are concerned. To overcome these difficulties a different regularization was proposed almost thirty years ago by K. Wilson . Wilson suggested to formulate the theory on a discrete lattice of points in Euclidean space-time. Such proposal has some very important advantages:
* The path integral becomes a collection of well defined ordinary integrals at the lattice sites.
* The lattice spacing becomes an ultraviolet cut-off.
* As far as the number of sites of the lattice is kept finite all the ultraviolet divergences are removed and all quantum averages are given by mathematically well defined expressions, for any value of the coupling constant.
* A QFT in $`d`$ space and 1 time dimensions regularized on a lattice becomes equivalent to an equilibrium statistical mechanics model in $`(d+1)`$ space dimensions. As a consequence one can study the model with all the tools which are typical of statistical mechanics like strong coupling expansion or Montecarlo simulations.
Obviously the lattice regularization is not a magic wand and all the problems which have been overcome appear again, in some other form, when we take the continuum limit (we shall deal with this very delicate issue in sect. 2.3). However the main feature of the regularization, i.e. the fact that it is intrinsically non-perturbative survives in the limit and allows one to obtain results which could never be obtained with standard perturbative expansions .
In this section we shall discuss in details the two main steps which allow one to construct (and extract results from) a lattice regularization of QFT:
a\] The translation from Minkowski to Euclidean Quantum Field Theory
b\] The connection between Euclidean QFT and Statistical Mechanics in the canonical Ensemble.
The starting point for a path integral formulation of QFT is the vacuum to vacuum amplitude (also called the generating functional) in presence of an external source $`J`$
$$Z[J]=𝑑\varphi e^{\frac{i}{\mathrm{}}{\scriptscriptstyle 𝑑td^3x(\varphi )}+J\varphi }.$$
(1)
Correlation functions can be obtained from this in the standard way by differentiating $`\mathrm{log}Z[J]`$ with respect to $`J`$, for example
$$0\left|T[\varphi (x)\varphi (0)]\right|0=\frac{\delta }{\delta J(x)}\frac{\delta }{\delta J(0)}\mathrm{log}Z[J]|_{J=0}.$$
(2)
Looking at eq.(1) we immediately recognize the analogy with the usual expression of the partition function in the canonical ensemble. The only difference (which is however of great importance) is that in the exponential of eq.(1) we have an oscillatory term while the argument in the exponent of a Boltzmann weight is real. We can bridge this difference with an analytic continuation of eq.(1) to imaginary values of the time. This is the well known “Wick rotation”.
$`x_0`$ $``$ $`tix_4i\tau ,`$
$`p_0`$ $``$ $`Eip_4.`$ (3)
In this way we obtain
$$Z[J]=𝑑\varphi e^{𝒮_E+J\varphi }$$
(4)
where $`𝒮_E`$ denotes the Euclidean action. We shall discuss in detail its form in the Yang Mills case in the next section. At this point we may well interpret $`𝒮_E`$ as the Hamiltonian of a static model in four space dimensions and $`Z[J]`$ with the corresponding partition function. It is far from obvious that we can perform a Wick rotation without problems. On the continuum this is granted by the good analyticity properties of the propagator, but on the lattice it imposes some strict constraint on the form of the discretized action. These constraints are known as reflection positivity conditions (and also as “Osterwalder and Schrader positivity conditions”).
The connection between QFT and Statistical mechanics is a crucial issue of modern quantum field theory. It has deep, far reaching, consequences in several physical contexts. Its main implications are summarized in tab. 1
### 2.2 Lattice discretization of pure Yang-Mills theories
The goal of this section is the explicit construction of a lattice regularization of a gauge theory with gauge group $`𝐆`$. To this end we need first of all a Wick rotated, Euclidean formulation of the theory (sect. 2.2.1). Then for its lattice discretization three main ingredients are needed: the lattice structure (sect. 2.2.2), the lattice definition of the gauge variables (sect. 2.2.3) and the action (sect. 2.2.4). In each of these steps we have a great amount of freedom. We shall always choose the simplest option, and leave to exercise 1 the discussion of possible alternative choices. We shall then check in sect. 2.2.5 that the proposed action gives in the “naive” continuum limit the expected gauge invariant expression and add some further remarks on the integration measure (sect. 2.2.6), on the fermionic sector (sect. 2.2.7) and on the constraints which must be imposed to obtain a finite temperature version of the theory (sect. 2.2.8).
#### 2.2.1 Euclidean Yang-Mills theories
In the following we shall be interested in the lattice formulation of $`YM`$ theories. The continuum limit of these models is the Euclidean version of $`YM`$ theories. Let us briefly remind its expression. The building blocks are the field $`A_\mu ^i`$, $`i=1\mathrm{}N`$ where $`N`$ is the number of the generators of the gauge group. The indices $`i,j,\mathrm{}`$ run in the space of the generators of the gauge group. The Yang Mills action is:
$$S_{YM}=\frac{1}{4}d^4xF_{\mu \nu }^iF^{i\mu \nu }$$
(5)
with
$$F_{\mu \nu }^i=_\mu A_\nu ^i_\nu A_\mu ^i+gf_{ijk}A_\mu ^jA_\nu ^k$$
(6)
where $`g`$ is the coupling constant and $`f_{ijk}`$ are the structure constants of the gauge group defined by
$$[\tau _i,\tau _j]=if_{ijk}\tau _k.$$
(7)
#### 2.2.2 The lattice
Let us choose the simplest possible lattice structure: a four dimensional hypercubic lattice $`𝚲`$ of size $`L`$ in the four directions. Let us denote the sites of the lattice with $`n(n_0,n_1,n_2,n_3)`$ and with $`\widehat{\mu }`$ the unit vector in the $`\mu `$ direction $`(\mu =0,1,2,3)`$. We shall often call in the following the $`0`$ direction as time-like and the other three as space-like, but let us stress that, due to the Wick rotation, all four directions are exactly on the same ground (we shall make use of this symmetry in the following). $`𝚲`$ contains $`N_s=L^4`$ sites, $`N_l=4L^4`$ links and $`N_p=6L^4`$ plaquettes. We shall denote with $`n_\mu `$ the link starting form $`n`$ and pointing in the positive $`\mu `$ direction i.e the link joining the two sites $`n`$ and $`n+\widehat{\mu }`$. Similarly $`n_{\mu ,\nu }`$ denotes the plaquette joining the four sites $`n`$, $`n+\widehat{\mu }`$, $`n+\widehat{\mu }+\widehat{\nu }`$, $`n+\widehat{\nu }`$. We shall denote with $`a`$ in the following the “lattice spacing” i.e. the separation between two nearby sites.
#### 2.2.3 The gauge field.
There is a standard recipe to choose the lattice analog of the bosonic fields of a continuum QFT. One must define the scalars on the sites, the vectors on the links and the two index tensors on the plaquettes of the lattice. This recipe can be simply understood if we rewrite the QFT in terms of differential forms. The discrete analog of a p-form is a p-simplex. In whole generality this rule can be stated as follows:
“The lattice analog of a tensor field of degree $`k`$ is a function which takes values on the $`k`$ simplexes of the lattice”
Thus the vector potential $`A_\mu (x)`$ must be defined on the links of the lattice. The simplest choice, which ensures both gauge invariance and a smooth continuum limit is to put in each link $`n_\mu `$ an element of the gauge group: $`U_\mu (n)𝐆`$ such that
$$U_\mu (n)=e^{cA_\mu (n)}$$
(8)
where $`c`$ is a suitable constant which we shall discuss below.
A nice, intuitive, way to understand this choice is the following: Imagine that in each site $`n`$ does exist an internal space $`E`$ (on which the gauge group $`𝐆`$ acts in a non trivial way). Let us assume that the reference frame $`E_n`$ on $`E`$ changes from site to site. Then we can interpret $`U_\mu (n)`$ as the transformation which relates the two nearby reference frames $`E_n`$ and $`E_{n+\widehat{\mu }}`$. An immediate consequence of this picture is that we must impose, for consistency:
$$U_\mu (n+\widehat{\mu })=U_\mu ^1(n)$$
(9)
In this framework a gauge transform is simply an arbitrary rotation, site by site of the reference frames $`E_n`$. Let us denote these transformations as $`V(n)𝐆`$. It is clear that the effect on $`U_\mu (n)`$ of such transformations is the following
$$U_\mu (n)V(n)U_\mu (n)V^1(n+\widehat{\mu })$$
(10)
Thus, as expected, the single variable $`U_\mu (n)`$ is not gauge invariant. The simplest way to construct, out of the gauge variables $`U_\mu (n)`$, gauge invariant observables, is to choose a closed path $`\gamma `$ on the lattice and then construct
$$W(\gamma )=\mathrm{Tr}\underset{n_\mu \gamma }{}U_\mu (n)$$
(11)
(where the product is assumed to be ordered along the path $`\gamma `$). This observable is usually called “Wilson loop”.
#### 2.2.4 The action
Obviously, the main requirement that we must impose on the discretized version of the action is that it must be gauge invariant. Among all the possible Wilson loops the simplest one is the product of the four link variables around a plaquette. If we sum these elementary Wilson loops over all the plaquettes of the lattice we obtain an expression which is invariant with respect to the discrete subgroups of the translational and rotational symmetries which survive in the lattice discretization. This was the original Wilson proposal for the lattice discretization of the gauge invariant action. Such an expression defines a perfectly consistent gauge invariant model for any group $`𝐆`$ on the lattice and is a good candidate to define a translational and rotational invariant gauge theory also in the continuum limit. At this stage we have no constraint on $`𝐆`$ which can well be a finite, discrete group. For instance several interesting results have been obtained in the case in which $`𝐆=Z_2`$ (the “gauge Ising model”). However since we are interested in constructing the lattice version of Yang-Mills theories we must concentrate on the case in which $`𝐆`$ is continuous Lie group. Even if the physically interesting case is $`𝐆=SU(3)`$, we shall study in the following the general case $`𝐆=SU(N)`$ with $`N2`$. This extension essentially does not add any further complication and allows to study the limit $`N\mathrm{}`$ which is essential if we aim to compare our lattice results with the AdS/CFT predictions.
Let us see it in detail. Let us define with $`U_{\mu \nu }(n)`$ the product of the four link variables around the plaquette $`n_{\mu \nu }`$ (see fig. 1)
$$U_{\mu \nu }(n)=\left\{U_\mu (n)U_\nu (n+\widehat{\mu })U_\mu (n+\widehat{\mu }+\widehat{\nu })U_\nu (n+\widehat{\nu })\right\}$$
(12)
The Wilson action is
$$S_W=\frac{\beta }{2N}\underset{n,\mu \nu }{}\mathrm{Re}\mathrm{Tr}\left\{U_{\mu \nu }(n)\right\},$$
(13)
where we have introduced the $`\frac{\beta }{2N}`$ constant in front of the action for future convenience. Notice that since the sum over $`\mu `$ and $`\nu `$ is unrestricted all the plaquettes of the lattice are counted twice. We shall check in the next section that this proposal indeed gives in the continuum limit the correct gauge invariant action.
* Exercise 1: discuss some possible generalizations of the lattice discretization of $`SU(N)`$ $`YM`$ theories. In the above derivation we chose at each step the simplest possible option. However there are infinitely many different lattice regularizations which lead to the same continuum limit of eq.(25). Discuss some of these possible generalizations.
#### 2.2.5 “Naive” continuum limit.
Let us call $`𝒢`$ the Lie algebra associated with $`𝐆`$. Let us assume that $`𝒢`$ has $`N`$ generators $`\tau _1\mathrm{}\tau _N`$. Then we can define:
$$U_\mu (n)e^{iB_\mu (n)}$$
(14)
with
$$B_\mu (n)ag\tau _iA_\mu ^i(n)$$
(15)
$`a`$ being the lattice spacing and $`g`$ a suitable coupling constant. As we shall see the $`A_\mu ^i`$ functions will become in the continuum limit the standard vector potential fields of the Yang Mills theory.
We may expand the fields appearing in the Wilson action (keeping only the first order in $`a`$) as follows.
$`B_\nu (n+\widehat{\mu })`$ $``$ $`B_\nu (n)+a_\mu B_\nu (n)`$
$`B_\mu (n+\widehat{\mu }+\widehat{\nu })`$ $``$ $`B_\mu (n+\widehat{\nu })B_\mu (n)a_\nu B_\mu (n)`$
$`B_\nu (n+\widehat{\nu })`$ $``$ $`B_\nu (n)`$ (16)
where we have denoted with
$$_\nu f(n)\frac{f(n+\widehat{\nu })f(n)}{a}$$
(17)
the finite difference on the lattice, whose continuum limit is the partial derivative:
$$_\nu f(n)_\nu f(x).$$
(18)
From the above expansions we obtain:
$$U_{\mu \nu }(n)e^{iB_\mu (n)}e^{i(B_\nu (n)+a_\mu B_\nu (n))}e^{i(B_\mu (n)+a_\nu B_\mu (n))}e^{iB_\nu (n)}$$
(19)
Let us use at this point the Baker-Hausdorff formula, which at the first order is:
$$e^xe^y=e^{x+y+\frac{1}{2}[x,y]}$$
(20)
Keeping in the expansion only terms up to $`O(a^2)`$ (remember that $`B_\mu `$ is of order $`a`$) we find
$$U_{\mu \nu }e^{\left\{ia(_\mu B_\nu _\nu B_\mu )[B_\mu ,B\nu ]\right\}}e^{ia^2gF_{\mu \nu }}$$
(21)
where we have neglected for simplicity the argument $`n`$ and we have defined:
$$F_{\mu \nu }_\mu A_\nu _\nu A_\mu +ig[A_\mu ,A_\nu ]$$
(22)
and $`A_\mu \tau _iA_\mu ^i`$.
Let us insert this result in eq.(13), and expand in powers of $`a`$. we find:
$$S_W\frac{\beta }{2N}\underset{n,\mu \nu }{}\mathrm{Re}\mathrm{Tr}\left(\mathrm{𝟏}+ia^2gF_{\mu \nu }\frac{1}{2}a^4g^2F_{\mu \nu }^2\right)$$
(23)
We can always parameterize the $`SU(N)`$ generators so as to have
$$\mathrm{Tr}(\tau _i)=0,\mathrm{Tr}(\tau _i\tau _j)=\frac{1}{2}\delta _{ij}.$$
(24)
while $`\mathrm{Tr}\mathrm{𝟏}`$ only gives an irrelevant constant which can be neglected. In this way we obtain
$$S_W=\frac{\beta a^4g^2}{8N}\underset{n,\mu \nu }{}F_{\mu \nu }^iF^{i\mu \nu }+O(a^5)$$
(25)
In the naive limit $`a0`$ we can set $`_n\frac{d^4x}{a^4}`$ (for a four-dimensional lattice) and we see that the dominant term in the $`a0`$ limit becomes the standard expression of the pure $`YM`$ action if we set
$$\beta =\frac{2N}{g^2}$$
(26)
From this position we see that $`\beta `$ is proportional to the inverse of $`g^2`$. It is also easy to see, from dimensional arguments that in four dimensions the coupling constant $`g`$ is adimensional while for 3d $`YM`$ theories $`g^2`$ has the dimensions of a mass.
#### 2.2.6 The partition function.
As mentioned above, our real interest is not in the gauge action but in the partition function (or generating functional in the QFT language) $`Z`$. In constructing $`Z`$ we must address the problem of the integration measure (and the related problem of gauge fixing this integration). Here we see one of the most interesting advantages of the lattice regularization. The integrals involved in the construction of $`Z`$ are site by site (or better, in the case of a gauge theory, link by link) ordinary integrals. We have a natural choice for the integration measure: the Haar measure (i.e. the invariant measure over the group manifold $`dU_\mu (n)`$). We have:
$$Z=\underset{n,\mu }{}dU_\mu (n)e^{S_W}$$
(27)
and similarly, for a generic expectation value we have
$$𝒪=\frac{_{n,\mu }dU_\mu (n)𝒪(U_\mu (n))e^{S_W}}{Z}$$
(28)
A remarkable consequence of these definitions is that, since the integration is made, link by link, over the whole group manifold all the gauge equivalent configurations (see eq.(10)) are automatically included in the sum with the correct weight. Contrary to the continuum case, on the lattice, the integration over the pure gauge degrees of freedom does not make quantum averages ill defined. In the lattice regularization there is no need to fix the gauge: all the quantum averages are by construction gauge invariant.
#### 2.2.7 Fermions.
In the previous sections we studied the lattice discretization of pure $`YM`$ theories. In full QCD we should also take into account the quarks. However putting fermions on the lattice is a rather non trivial issue. Moreover, once a consistent discretization of the theory is obtained if we try to integrate out the fermions from the Lagrangian we obtain a determinant of the gauge fields which is very difficult to handle in Montecarlo simulations. In view of this consideration it has become a common habit to organize lattice gauge theories in three levels of increasing complexity.
* In the first level we find pure $`YM`$ theories. These models are by now rather well understood and precise Montecarlo results exist for the continuum limit of several quantities, among these the most interesting ones are the glueballs since we may rather safely expect that their mass should not be too affected by the absence of quarks.
* The second level is the so called “quenched QCD” in which the quarks are explicitly added into the game, thus giving the possibility to explore several new observables (for instance the meson spectrum) which are of great interest for the phenomenology, but they are kept quenched, i.e. the determinant mentioned above is simply neglected. This means that in the partition function the quarks are treated as classical quantities or equivalently that we are neglecting quark loops in our calculations. Also for quenched QCD stable and reliable result have been obtained from Montecarlo simulations. However there are effects, like the string breaking at large distance, which are a typical consequence of quarks loops and that cannot be observed in the quenched approximation. Moreover it is by now clear that several quantities of great physical interest (like the meson spectrum or the temperature of the chiral phase transition) are heavily affected by such approximation.
* The third level is that in which full QCD, with dynamical fermions is discretized on the lattice. In this last case Montecarlo simulations are still at a preliminary stage and a few years (and the next generation of supercomputers) will be needed before we may reach stable and reliable results also in this case.
Luckily enough, the predictions of the AdS/CFT correspondence, whose comparison with the lattice results is the main goal of these lectures, refer to the pure $`YM`$ theories. Thus on the lattice side of the comparison we may rely on very stable and trustable numbers. Moreover in pure $`YM`$ theories extensive simulations exist in three dimensions also for values of $`N`$ larger than 3 and reliable large $`N`$ limits for various quantities of interest can be obtained . In four dimensions the results of these large $`N`$ limits are still preliminary , but nevertheless they are stable enough to allow for the comparison with the AdS/CFT results. In the following, both in the lattice sections and in the AdS/CFT ones we shall make an effort to keep well distinct those results which refer to pure Yang Mills theories from those which refer to full QCD.
#### 2.2.8 Finite temperature LGT.
We shall see in sect. 4 that in the framework of the AdS/CFT correspondence a very important role is played by the “temperature” of the theory. It turns out that it is only at high temperature that we can get rid of the supersymmetry and obtain non-supersymmetric Yang-Mills like theories. It is thus important to understand how can we describe finite temperature on the lattice.
The model that we have defined and studied in the previous sections describes $`YM`$ theories at strictly zero physical temperature. The parameter which in the statistical mechanics counterpart of the model plays the role of the temperature becomes in LGT the coupling constant of theory. The discussion of sect. 2.1 can be easily extended so as to implement a non-zero temperature also in LGT. The main ingredient is that periodic boundary conditions must be imposed in the “time direction”. Then it can be shown that the inverse of the lattice size in this direction is proportional to the temperature of the theory.
In general, for technical reason, one imposes periodic boundary conditions in all the $`d`$ directions, but one also takes care to choose the lattice length much larger that the typical correlation length of the theory, so as to make negligible the effect of the boundary conditions. On the contrary if we want to see the effects of the finite temperature in the theory, the lattice size in the “time” direction must be chosen of the same order of the correlation length. If we increase the lattice size in the “time” direction then we smoothly reach the zero temperature limit, which is effectively reached when the effects of the periodic boundary condition become negligible. Thus typically a finite temperature discretization requires asymmetric lattices, with one direction much shorter than the others. As a consequence the original equivalence of space and time directions which is a typical feature of Euclidean QFT is lost. In finite temperature LGT (FTLGT in the following), we have a very precise notion of “time” which is the compactified direction proportional to the inverse of the temperature.
Due to the presence of periodic boundary conditions a new class of observables exists in FTLGT i.e. the loops which close winding around the compactified time direction. These are usually called Polyakov loops. The discussion of FTLGT (apart from a few issues that we shall address in sect. 3) is beyond the scope of these lectures. We refer to the reviews in the bibliography for further details.
### 2.3 Continuum limit
It is clear that by simply sending $`a0`$ as in the previous section we do not obtain a meaningful continuum limit. In particular, all the dimensional quantities, which will be proportional to a non-zero power of $`a`$ will go to zero or infinity. This is the meaning of the word “naive” used above.
Let us study this problem in more detail. Let us take a physical observable $`𝒪`$ of dimensions $`d_𝒪`$ in units of the lattice spacing. Let us assume that we have in some way calculated the mean value of $`𝒪`$ in the lattice regularized version of the theory. The result of this calculation will take the form:
$$𝒪=a^{d_𝒪}f_𝒪(g)$$
(29)
where $`f_𝒪(g)`$ is a suitable function of the coupling constant of the theory (in general of all the parameters of the model if they are more than one). For instance $`𝒪`$ could be one of the correlation lengths of the model (i.e. the inverse of the mass of one of the states of the theory). In this case $`d_𝒪=1`$ and $`f_𝒪(g)`$ measures the correlation length in units of the lattice spacing, for the particular value $`g`$ of the coupling constant. It is now clear that if we simply send $`a0`$ we obtain the trivial result $`𝒪=0`$. In order to have a meaningful continuum limit we must change $`g`$ at the same time as $`a`$ is set to zero in such a way as to make the observable to approach a well defined finite value in the limit. In the example this means that we must tune $`g`$ to a critical value $`g_c`$ in which the correlation length measured in units of the lattice spacing goes to infinity.
From a statistical mechanics point of view this implies that at $`g=g_c`$ the system must undergoes a continuous phase transition. This is a mandatory requirement for a non-trivial continuum limit. From the point of view of QFT we have a nice interpretation of this constraint. While the lattice discretization gives a way to regularize the theory. The process of tuning $`g`$ to its critical value, thus removing the cut-off, corresponds to the renormalization of the theory.
This process is highly non trivial, since for all the physical quantities $`𝒪_i`$ that we may define in the model we must require a well defined continuum limit. If $`f_i(g)`$ is the function which measures the value of $`𝒪_i`$ in units of the lattice spacing, then we must require that as $`gg_c`$ all the functions $`f_i(g)`$ go to zero or infinity (depending on the sign of $`d_{𝒪_i}`$) in such a way that the same rate of approach of $`g`$ to $`g_c`$ which makes the correlation length $`\xi `$ tend to a constant value also makes all the observables $`𝒪_i`$ tend to constant values. This stringent requirement is better understood in the framework of the Renormalization Group approach and is commonly summarized, by saying that the critical point $`g_c`$ must be a scaling critical point.
It is clear from the above discussion that a meaningful continuum limit requires a precise functional relationship between $`a`$ and $`g`$. However in principle we can even ignore such a dependence. The notion of scaling defined above allows to reach a well defined continuum limit even if we do not know how to fix $`g`$ as a function of $`a`$. If we know the physical value of one of the observables of the theory, say the mass $`m_e`$ of a particle “e” which is easily accessible from the experiments, then we may set the overall scale of the theory as follows
$$m_e=\frac{1}{a(g)\xi _e(g)}$$
(30)
where $`\xi _e(g)`$ is the particular correlation length related to the particle “e”. Then the continuum limit value of any other dimensional quantity of the theory, like for instance the masses $`m_i`$ of other particles, can be obtained as
$$m_i=\underset{a0}{lim}\frac{1}{a(g)\xi _i(g)}=\underset{gg_c}{lim}\frac{\xi _e(g)}{\xi _i(g)}m_e$$
(31)
This is the power of scaling!
In some cases it may happen that, on top of the above relations, we also have some independent way to fix asymptotically (i.e. in the vicinity of the critical point) the relationship between $`g`$ and $`a`$. This is the exactly the case for non-abelian gauge theories, for which asymptotic freedom tells us that $`g_c=0`$ and perturbative methods can be used. This leads to the following well known expression, for a pure $`SU(N)`$ gauge theory in four dimensions:
$$a=\frac{1}{\mathrm{\Lambda }}f(g)$$
(32)
with
$$f(g)=(g^2\beta _0)^{\beta _1/(2\beta _0)}e^{1/(2\beta _0g^2)}(1+O(g^2))$$
(33)
where
$$\beta _0=\frac{11N}{48\pi ^2},\beta _1=\frac{34}{3}\left[\frac{N}{16\pi ^2}\right]^2$$
(34)
are the first two coefficients of the Callan-Symanzik function $`\beta (g)`$ and $`\mathrm{\Lambda }`$ is a scale parameter, which does not have a direct physical meaning and in general depends on the renormalization scheme that we have chosen (for instance it may depend on the type of lattice that we have chosen).
One usually refers to this relation (and to the procedure of taking the continuum limit following eq.(33)) as “asymptotic scaling” to stress the fact that one is using more informations than those simply implied by the scaling property.
Since it will play an important role in the following it is worthwhile to discuss in detail how one should use eq.(33). To this end let us continue with the above example, and let us assume again that $`𝒪`$ is the correlation length of the theory (i.e the inverse of the lowest glueball). Let us measure it on the lattice for various values of $`g`$ in the scaling region in units of the lattice spacing $`a`$ and let us call these numbers (which at this point are pure adimensional numbers) $`\xi (g)`$. We have
$$𝒪=a\xi (g)$$
(35)
Let us now insert eq.(32), we find
$$𝒪=\frac{f(g)\xi (g)}{\mathrm{\Lambda }}$$
(36)
Our goal is to find a finite continuum limit value (let us call it $`\frac{\xi _0}{\mathrm{\Lambda }}`$) for $`𝒪`$. This implies that $`\xi (g)`$ must scale as:
$$\xi (g)=\frac{\xi _0}{f(g)}$$
(37)
If this condition is fulfilled by the data then we may say that the lowest glueball has a mass in the continuum limit whose values is
$$m=\frac{1}{\xi }_0\mathrm{\Lambda }$$
(38)
if we are able to fix the value of $`\mathrm{\Lambda }`$ in $`MeV`$ (and this can be done, for instance, by comparing the string tension evaluated on the lattice with the physical value of the string tension obtained from the spectroscopy of the heavy quarkonia) then eq.(38) will give us the value in $`MeV`$ of the lowest glueball. The precision of this prediction will only be limited by the uncertainty in the determination of $`\mathrm{\Lambda }`$ in $`MeV`$ and by the statistical and systematic errors which affect the estimate of the amplitude $`\xi _0`$ from a fit to the data $`\xi (g)`$ according to eq.(37).
Let us conclude this section by noticing that the functional dependence on $`g`$ of eq.s(32),(33) is a direct consequence of the fact that the coupling $`g^2`$ in this case is adimensional. In fact the scaling behaviour for $`SU(N)`$ gauge theories in three dimensions is completely different. In this case it is $`g^2`$ itself (which has the dimensions of a mass) which sets the overall scale for the theory. This means that near the continuum limit a physical observable with the dimensions of a mass like for instance the inverse of the correlation length that we studied above, can be written as a series in powers of $`g^2a`$.
$$\frac{1}{a\xi (g)}=mg^2(1+m_1g^2a+m_2g^4a^2+\mathrm{})$$
(39)
where $`m`$ is the continuum limit value of the mass and the constants $`m_1,m_2\mathrm{}`$ measure the finite $`a`$ corrections to the scaling behaviour.
## 3 Extracting physics from the lattice.
In order to extract some physical results from the lattice regularized model we need three steps.
* First, we must define the lattice version of the quantities in which we are interested. In general there is not a unique prescription to do this. As in the previous section we shall always choose the simplest lattice realization and shall then comment on the possible generalizations.
* Second, we must use some non perturbative technique to extract the expectation value of these operators for a (possibly wide) set of values of the parameters of the model (in the simplest case only the coupling constant $`g`$). In some very exceptional situation (like 2d $`YM`$ theories) exact results can be obtained with analytic techniques, but in general some approximate method is needed. The most popular tool is the Montecarlo simulation, however in some situations also strong coupling expansions can give reliable results.
* Third, one must test that the $`g`$-dependence of the measured quantities scales according to eq.(33). If this condition is fulfilled then one can set $`a0`$ and extract in the continuum limit the value of the observable under study, in units of the physical scale $`\mathrm{\Lambda }`$.
Let us see these steps in detail
### 3.1 Lattice observables.
There are several quantities which are accessible on the lattice. For the reason discussed in the introduction we shall concentrate in the following only on two of them: the string tension and the glueball spectrum.
#### 3.1.1 The string tension $`\sigma `$
Phenomenologically we know that the quark and the antiquark in a meson are tied together by a linearly rising potential. The simplest way to describe such a behaviour is to assume that the infrared regime of QCD is described by an “effective” string (which, as we shall see, is very different from the one which lives in AdS) which joins together quark and antiquark. This is the origin of the name“string tension” to describe the strength of the rising potential (we shall discuss in detail this effective string picture in sect. 3.5).
In the real world the best set up to extract experimental informations on the string tension $`\sigma `$ is represented by the spectrum of the heavy quarkonia where, thanks to the large masses of the quarks, the quark-antiquark pair can be studied with non-relativistic techniques. Suitable potential models can be used to fit the spectrum and in this way an experimental estimate for $`\sigma `$ can be extracted.
On the lattice the simplest way to mimic the quark-antiquark pair is to study the mean value of a large rectangular Wilson loop (say of sizes $`R\times T`$). Let us denote it as $`W(R,T)`$ (see fig. 2). Let us also assume that T is a segment in the “time” direction (remember however that the notion of time direction is purely conventional in our Euclidean framework) from the time $`t_0`$ to the time $`t_0+T`$, and $`R`$ a segment in any one of the space directions.
The physical interpretation of $`W(R,T)`$ is that it represents the variation of the free energy due to the creation at the time $`t_0`$ of a quark and an antiquark which are instantaneously moved at a distance $`R`$ from each other, keep their position for a time T and finally annihilate at the instant $`t_0+T`$. According to this description we expect for large $`T`$
$$W(R,T)e^{TV(R)}$$
(40)
where $`V(R)`$ is the potential energy of the quark-antiquark pair. The interquark potential is thus given by:
$$V(R)=\underset{T\mathrm{}}{lim}\frac{1}{T}\mathrm{log}W(R,T)$$
(41)
If we assume that $`V(R)`$ is dominated by the linear term $`\sigma R`$ then we end up with the celebrated “area law” for the Wilson loop:
$$<W(R,T)>=e^{\sigma RT+p(R+T)+k}.$$
(42)
The area term is responsible for confinement while the perimeter and constant terms are non universal contributions related to the discretization procedure. When one takes the $`T\mathrm{}`$ limit to extract the potential (see eq.(41)) the constant term disappears and the perimeter one gives a normalization constant $`V_0`$ which can be easily fixed from a fit. The physically important quantity is the coefficient of the area term which represents the lattice estimate of the string tension which we were looking for.
We shall see below that the real expression for $`W(R,T)`$ is slightly more complicated, but in any case the dominant term is the string tension, which is given by
$$\sigma =\underset{T\mathrm{}}{lim}\frac{1}{RT}log(W(R,T))$$
(43)
It is easy to see <sup>1</sup><sup>1</sup>1All the dimensional quantities discussed in this section $`R,T,\sigma `$ and $`V(R)`$ are measured in units of $`a`$, but here and in the following we shall omit the factors of $`a`$ (which can be easily deduced from a dimensional analysis) to avoid a too heavy notation. that $`\sigma `$ has dimensions of $`a^2`$ and as a consequence its expected scaling behaviour in $`d=4`$ is, according to the discussion at the end of sect. 2.3,
$$\sigma (g)=\sigma _0f^2(g)$$
(44)
where $`f(g)`$ is given by eq.(33). From the numerical value of $`\sigma _0`$ we may then obtain the continuum limit value $`\sigma _0\mathrm{\Lambda }^2`$ (in units of $`MeV^2`$ once we have fixed the value of $`\mathrm{\Lambda }`$) for the string tension.
#### 3.1.2 Glueballs.
As it is well known, pure gauge theories have a rich spectrum of massive states which are called glueballs. This is one of the most remarkable effects of quantization since classical gauge field theories do not contain any mass terms and are scale-invariant. Moreover it is a truly non-perturbative effect: in perturbation theory the gluon propagator remains massless to all orders.
The lattice offers a perfect framework to study the glueball spectrum and in fact on the lattice massive states of pure gauge models arise in a very natural way. One must select two elementary ”space-like” plaquettes at two different positions in the “time” direction. Looking at the large distance (i.e. large time) behaviour (with a Montecarlo simulation or with a strong coupling expansion) of the connected correlator of the two plaquettes one immediately recognizes the exponential decay which denotes the presence of a massive state. From this behaviour one can extract the correlation length, whose inverse is the mass of the lowest glueball (the $`0^{++}`$ state).
$$\mathrm{Tr}(U_{ij}(𝐱,t_1))\mathrm{Tr}(U_{ij}(𝐱,t_2))e^{M|t_1t_2|}$$
(45)
where $`(U_{ij}(𝐱,t_1))`$ is the plaquette (with spacelike indices $`(i,j)`$) located in the point $`(𝐱,t_1)`$ of the lattice, and $`M`$ is the mass of the lowest glueball.
It is nice to observe that the exponential decay is already visible (without the need of a MC simulation) in the very first order of a strong coupling expansion, a very simple calculation that we shall perform below as an exercise on strong coupling expansion. However for the non abelian gauge theories which are of physical relevance, the SC series are limited to rather few terms and allow to obtain only a rough estimate of the spectrum. If one is interested in comparing the lattice with possible experimental candidates it is mandatory to use MC simulations.
In order to extract good estimates for higher states of the spectrum one must study connected correlators (in the time direction) of more complicated combinations of space-like Wilson loops of suitable shapes. These combinations are chosen so as to match the symmetry properties of the glueball (which is encoded in the notation $`J^{PC}`$). This is a rather subtle issue since on the lattice the rotation group is broken to its cubic subgroup. This has two relevant consequences:
1\] Some of the irreducible representations of the rotation group are not any more irreducible with respect to the cubic subgroup and hence split in several components which on the lattice, in principle, correspond to different massive states. These masses must coincide as $`a0`$ if we want to recover the correct continuum theory. This represents a non trivial consistency test of the continuum limit extrapolation. For finite values of $`a`$ the splitting between these “fragments” of the same state gives an estimate of the relevance of lattice artifacts.
2\] Conversely, since there is only a finite number of irreducible representations in the cubic group, any one of them must contain infinitely many irreps of the rotation group. In particular all values of $`J`$ (mod 4) coincide in the cubic subgroup (they can be recognized as metastable states following the suggestions of ref ).
Constructing the correct identifications between the continuum states and their lattice realizations turns out to be a non-trivial (and very instructive) exercise of group theory. We shall discuss as an example in exercise 2 the construction in the case of the $`SU(2)`$ theory in (2+1) dimensions. Even if this is the simplest possible situation it is already complex enough to show all the subtleties of the problem. The generalization to other values of $`N`$ and to the (3+1) dimensional case can be found, for instance, in . It is also possible to disentangle glueballs with the same $`J^{PC}`$ but different radial quantum numbers. This is very important since in general the first state above the $`0^{++}`$ is its first radial excitation and not a glueball of different angular momentum. We can summarize the discussion of this section (and the analysis performed in exercise 2) as follows
* For any glueball state of quantum numbers $`J^{PC}`$ with $`J<4`$ it does exist a lattice representation in terms of suitable combinations of spacelike Wilson loops of various shapes. These combinations can be constructed by using group theory. Let us call them $`W(J^{PC})`$
* each $`W(J^{PC})`$ gives a lattice realization for a whole family of glueball states with angular momentum $`J`$ (mod 4). Looking at the large distance (in the “time” direction) behaviour of the connected correlator of $`W(J^{PC})`$ we may extract the one of lowest mass. Higher glueballs can be seen as exponentially suppressed corrections in the correlator or as metastable states.
* This representation is not unique, in general there are infinitely many combinations of spacelike Wilson loops with the same symmetry properties. By using some variational method we can select those combinations which enhance the particular higher mass state in which we are interested (say the first radial excitation of the $`0^{++}`$) and thus measure its mass.
* Exercise 2: group theoretical analysis of the glueball states for the $`SU(2)`$ LGT in $`d=3`$.
### 3.2 Strong Coupling Expansions
In sect. 2.1 we have shown that there is a correspondence between QFT and Statistical Mechanics. In particular we can interpret the lattice regularized $`YM`$ theories as a peculiar statistical model in which $`g^2`$ plays the role of the temperature. One of the most powerful tools to study statistical models are the high temperature expansions (i.e. perturbative expansions in the inverse of the temperature). The main ingredient in this game is the expansion of the Boltzmann factor of the model on the character basis. In such a basis it becomes very simple to perform, order by order, the sum over all the possible configurations of the model which appear in the partition function and in the correlators. Moreover, by using the orthogonality properties of the characters, a set of rules can be constructed which greatly simplify the terms in the expansion. The final result can be written as a series in powers of $`\beta `$ i.e. perturbative in $`1/T`$ as desired.
It is easy to export this technique to LGT. The important point is that, thanks to the identification between temperature in statistical mechanics and coupling constant $`g^2`$ in LGT, the high $`T`$ expansion becomes in LGT a strong coupling expansion, i.e. an expansion in the inverse of the coupling constant. But this is exactly what we need to study the non-perturbative physics of $`YM`$ theories, and in fact in the strong coupling limit all the features that we expect to find in the theory (and have never been able to proof), like the linear confinement and a nonzero mass gap for the glueballs can be explicitly shown.
In principle, if we could push the strong coupling expansion (which is centered in $`g^2=\mathrm{}`$) up to the scaling region (small values of $`g^2`$) we would reach the long sought continuum limit description of the non-perturbative physics of $`YM`$ theories. Unfortunately this seems a too difficult task. In some cases, like for the Wilson loop that we shall discuss in detail below, it can be shown that the task is actually impossible and that the scaling region is separated form the strong coupling region by a phase transition, the roughening transition, which cannot be overcome. In some other cases (like for the glueball masses) the problem is that too many orders in the strong coupling expansion are needed to reach the scaling limit<sup>2</sup><sup>2</sup>2 It is important to stress that this is only a technical and not a conceptual problem. In fact, for instance, for the simplest possible gauge theory, i.e. the $`Z_2`$ gauge theory in three dimensions, the SC expansion has been pushed to so high levels that it gives results for the lowest masses of the spectrum which are comparable in precision with those obtained with MC simulations ..
A second reason of interest, which is particularly important from the point of view of the comparison that we have in mind in these lectures, is that in the framework of SC expansions a string description of LGT arises in a very natural way. In fact both the partition function and the correlators of gauge invariant operators become in the SC limit sums over suitably chosen surfaces. This is to be contrasted with the case of ordinary (not gauge invariant) field theories regularized on the lattice where the SC expansion becomes a sum over paths instead of surfaces.
In order to clarify the above statements let us study in more detail how the strong coupling expansion works in LGT. We need first of all to expand the Wilson action in the character basis of $`SU(N)`$. This is a standard problem in group theory and we shall discuss it in the exercise 3 below. The next step consists in substituting this expansion in each plaquette and then perform the group integrations over the links. One easily sees that, due to the orthogonality relations of characters, the only terms which survive in the expansions are those in which the plaquettes are “glued” together to form a closed surface. If we are interested in the partition function this is the end of the story. The partition function becomes a weighted sum over all possible closed surfaces that we can construct on the lattice. As anticipated above, this sum strongly resembles the discretized version of some (unknown) string-type theory. If we are instead interested in the expectation value of some gauge invariant operator described by a closed contour $`\mathrm{\Gamma }`$ it is easy to see that the first contribution in the strong coupling limit is given by the minimal surface bounded by $`\mathrm{\Gamma }`$. Further terms in the expansion will come from the fluctuations around this minimal surface. Again, this result strongly suggests a string like description for these observables.
As an example we report in the exercises 4 and 5 the calculation of the first term in the SC expansion for the Wilson loop and for the lowest glueball. The results are (see eq.s (E4.3) and (E5.3)):
$$\sigma =\mathrm{log}D_f(\beta )$$
(46)
$$M(0^{++})=4\mathrm{log}D_f(\beta )$$
(47)
where $`f`$ denotes the fundamental representation and $`D_f(\beta )`$ is given, for a generic value of $`N`$ by eq.(E3.10)
Let us see two examples which are particularly relevant for us: the $`SU(2)`$ case which is the simplest possible non abelian $`YM`$ theory and the large $`N`$ limit which is the limit in which the results obtained using the AdS/CFT correspondence are expected to hold.
For $`SU(2)`$ we have (see eq.(E3.16)
$$D_f(\beta )D_{\frac{1}{2}}(\beta )=\frac{I_2(\beta )}{I_1(\beta )}$$
(48)
where $`I_1`$ and $`I_2`$ are modified Bessel functions of integer order.
In the large $`N`$ limit we find first of all that a consistent limit can only be obtained by sending also $`\beta \mathrm{}`$ and keeping $`\beta /N`$ fixed (in agreement with the ’t Hooft prescription that we shall discuss in sect. 3.4). In this limit we find (see eq.(E3.17))
$$D_f(\beta /N)=\frac{\beta }{2N}.$$
(49)
The discussion of this section only gives a very short account of all the richness and complexity of SC expansions in LGT. The standard reference for further readings is where a very detailed and complete discussion of the subject can be found.
* Exercise 3: Character expansion for the $`SU(N)`$ group.
Construct the character expansion of the Wilson action eq.(13).
* Exercise 4: Evaluate the first order in the strong coupling expansion of the Wilson loop in $`SU(2)`$ $`YM`$ theory.
* Exercise 5: Evaluate the first order in the strong coupling expansion of the lowest glueball mass in $`YM`$ theories.
### 3.3 Scaling.
Once we have obtained with some non-perturbative method the value of a dimensional physical quantity for some fixed value of $`\beta `$ we face the problem of extracting a continuum limit estimate out of these numbers. To this end one must first check that the values that we have measure scale as a function of $`\beta `$ according to the expected asymptotic scaling behaviour. However it is often much simpler to test the behaviour of adimensional ratios of different observables. The reason is that very often the deviations from the asymptotic scaling behaviour (due for instance to irrelevant operators) cancel in the ratio. As a general rule the adimensional ratios are more stable than the single observables.
Notwithstanding this trick, one has in general to face rather large deviations from the expected scaling behaviours. The obvious solution to this problem would be to study very large values of $`\beta `$. However both SC expansions and MC simulations cannot be easily pushed upto these regions. SC expansions are centered in $`\beta =0`$ and very high orders are needed to obtain stable results at large $`\beta `$. MC simulations suffer of the so called “slowing down” problem. As the correlation length increases it becomes more and more difficult to obtain statistically independent configurations. Thus, practically, MC simulations are confined to not too large values of $`\beta `$. It is thus very important to have a good control of the systematic (not statistical!!) errors involved in extrapolating toward the continuum limit the MC results. There is by now a well developed technology to play this game. However one should always consider the results obtained from MC simulations with some caution.
A completely different problem is represented by the possible presence of phase transitions in the phase diagram of the model. If the range of $`\beta `$ values that we can study is separated from the continuum limit $`\beta =\mathrm{}`$ by a phase transition (which cannot be overcome by a suitable modification of the action), then there is no hope to be able to obtain reliable continuum estimates of the physical quantities.
The most important example of such a situation is represented by the SC expansion for the string tension. For a finite value $`\beta _r`$ of the coupling the Wilson loop undergoes a phase transition (the well known roughening transition) which does not allow to extend the SC series up to $`\beta =\mathrm{}`$. For this reason in studying the string tension we must only resort on MC simulations.
### 3.4 Large $`N`$ limit and the loop equations.
We have seen in the previous sections that the main advantage of the lattice discretization is that it is a truly non-perturbative regularization of QCD. The price that one has to pay is the introduction of the lattice spacing and the difficult part of the game becomes the elimination of this new scale so as to reach the correct continuum limit of the theory. It would be of great importance to have some kind of non-perturbative insight of the theory directly in the continuum. The large $`N`$ limit of ’t Hooft represents the most concrete proposal in this direction.
’t Hooft’s proposal starts from the observation that in non-abelian gauge theories another dimensionless quantity exists besides the bare coupling constant $`g`$. It is the number of colours $`N`$.
The main problem with QCD is that $`g^2`$ is not a good expansion parameter for the theory since (as we have seen in sect. 2.3) it runs with the cutoff. In fact the correct way to deal with $`g^2`$ is to trade it and the cutoff for the Renormalization Group invariant scale $`\mathrm{\Lambda }_{QCD}`$ (see eq.(32)). ’t Hooft was able to show that $`1/N`$ is indeed a much better expansion parameter than $`g^2`$ and that in the large $`N`$ limit the theory drastically simplifies. Before discussing these simplifications let us concentrate on the limit itself.
If we look at eq.(33) we see that $`g^2`$ always appears multiplied by $`\beta _0`$ which is proportional to $`N`$, thus if we want to keep $`\mathrm{\Lambda }_{QCD}`$ fixed as we take the $`N\mathrm{}`$ limit we must simultaneously take $`g0`$ keeping the so called ’t Hooft coupling $`\lambda g^2N`$ fixed. In this limit the Feynman graphs are well defined and can be organized in a perturbative expansion in powers of $`1/N`$. Now the seminal observation of ‘t Hooft was that this perturbative expansion is actually an expansion in the $`topology`$ of the Feynman graphs. To understand this statement let us first of all explain what do we mean for “topology of a graph” .
Let us assume the simplest possible definition of graph, i.e. a collection $`(V,E)`$ of Edges $`E`$ and Vertices $`V`$ in which all the edges are of the same type. Then it is possible to associate a genus to each graph by noticing that each graph can be unambiguously embedded in a 2d Riemann surface and hence can be characterized by its genus. For instance the graphs with genus 0 are the planar graphs which, in fact, are exactly those graphs which can be “drawn” on a sphere. It is far from obvious that the Feynman graphs of a non-abelian gauge theory (with different propagators for quarks and gluons) fall into the above definition, but ’t Hooft was able to give a set of recipes to allow this identification.
Let us now study as an example the $`1/N`$ expansion of the free energy. It turns out that the first term in the expansion is proportional to $`N^2`$ (this was to be expected since if we keep $`\lambda `$ fixed then the Lagrangian itself becomes of order $`N^2`$) and that only even powers of $`1/N`$ appear. Thus the expansion can be written as:
$$F=\underset{𝐠=0}{\overset{\mathrm{}}{}}N^{22𝐠}F_𝐠(\lambda )$$
(50)
where the $`F_𝐠`$ are complicated unknown functions of $`\lambda `$ which, in QCD-like theories are better expressed as functions of $`\mathrm{\Lambda }_{QCD}`$.
The remarkable result of ’t Hooft is that $`F_𝐠(\lambda )`$ only contains Feynman graphs of genus $`𝐠`$. Thus the expansion obtained in this way strongly resembles the loop expansion in string theory if we identify $`N`$ with the inverse string coupling $`1/g_s`$. This is one of the strongest indications in favour of a string-like description of QCD.
The above observations have some important consequences. Let us discuss them in detail.
* The planar limit
In the large $`N`$ limit only the planar graphs survive ($`𝐠=0`$) and the original theory greatly simplifies. In two dimensions the planar theory can be solved exactly and thus, at least in this case, an exact solution for QCD in the large limit can be obtained.
* The Master Field.
The most interesting implication of the large $`N`$ limit is the idea of the so called ”master field”. The starting point is the observation that in the large $`N`$ limit disconnected diagrams are in general dominating. This implies that if we study the vacuum expectation value of a collection of (gauge invariant) operators $`𝒪_i`$, all the propagators (i.e. connected correlators) joining together two of the operators disappear in the large $`N`$ limit and the VEV becomes the product of the VEV’s of the single operators.
$$𝒪_1𝒪_2\mathrm{}𝒪_n=𝒪_1𝒪_2\mathrm{}𝒪_n$$
(51)
This means that in the large $`N`$ limit the functional integral defining the above correlation function must be dominated by a single field configuration, which is usually called master field .
This important result gave the hope that it could be possible to explicitly solve QCD in the large $`N`$ limit and was the starting point of a large number of papers discussing the peculiar properties of the master field and the equations which it must satisfy.
These methods were successfully applied to some simple problems for which indeed a master field could be explicitly found. However in the most interesting cases of QCD in three and four dimensions none of these approaches has led so far to an explicit expression for the master field. One of the reasons of interest in the AdS/CFT correspondence is that it gives the first example of a master field solution for a set of non trivial, interacting, gauge theories in more than two dimensions.
* The loop equations.
A related result is that in the large $`N`$ limit the Wilson loop satisfies a closed set of equations, which are called “loop equations” . These equations can be derived in a rigorous way in the framework of the lattice regularization and it can be shown that they are solved by the master field of the theory. They can also be derived, at least formally, for the continuum theory where it can be shown that they are formally satisfied order by order in the perturbative expansion (for a review see ).
Unfortunately, despite many efforts, these equation could be solved explicitly only in the case of 2d $`YM`$ theories and it turned out to be impossible to extend this solution also to the $`d>2`$ case. However, even if they cannot be solved exactly, these equations are a very interesting object themselves. In fact they are, so to speak, an intrinsic, defining, property of $`YM`$ theories. In the lattice version of the theory the loop equations hold for any value of the coupling, both near the continuum limit and deep in the strong coupling region. Thus they are a perfect tool to test also in the strong coupling regime if a theory which we hope could be identified with QCD displays the correct large $`N`$ behaviour.
* Eguchi-Kawai models.
The most remarkable feature of the lattice version of the loop equations is that they allow in the large $`N`$ limit to reduce the whole lattice model to a much simpler one plaquette model, while keeping the full physical content of the original model. This idea was first proposed by Eguchi and Kawai and subsequently perfected by several authors and it is based on the observation that in the large $`N`$ limit a suitably twisted<sup>3</sup><sup>3</sup>3The twist consists in a suitable phase factor belonging to the center of SU$`(N)`$ that multiplies each plaquette variable in the action. lattice gauge theory on a lattice consisting of just one site and one link variable for each space-time direction generates the same set of loop equations as a theory defined on a large lattice, typically consisting of $`N^{(d+1)/2}`$ sites. Hence twisted one plaquette models can be used to describe lattice gauge theories on large lattices, by essentially mapping space-time degrees of freedom into internal degrees of freedom. A general review of the Eguchi–Kawai model, can be found in .
Let us conclude by noticing that in the framework of the AdS/CFT correspondence in order to obtain non-supersymmetric $`YM`$-like theories, one must look at the finite temperature behaviour of the large $`N`$ model in which the “time” direction is compactified. Trying to perform the same construction in Eguchi–Kawai type models turns out to be a rather non-trivial task due to the interplay between the twists needed to define the EK model and those induced by the periodic boundary conditions in the time direction. We refer the reader to for a discussion of this problem and a review of the remarkable properties of the large $`N`$ limit in finite temperature LGTs.
### 3.5 The effective string picture.
#### 3.5.1 The roughening transition
It is important to stress that the roughening transition is $`not`$ a phase transition of the model itself. At the roughening point the LGT partition function is regular, and the correlation length of the model (the inverse of the lowest glueball mass) is finite. The roughening point is instead the point in which the expectation value of one particular observable: the Wilson loop becomes singular. This means that for all the observables different from the Wilson loop (and in particular for instance for those related to the glueball states) there is no obstruction (i.e. no phase transition in between) to reach the continuum limit starting from the strong coupling phase.
On the contrary, as far as the Wilson loop is concerned, the confining regime of LGTs contains (in general) two phases: the strong coupling phase and the rough phase. The two are separated by the roughening transition which is the point in which the strong coupling expansion of the Wilson loop ceases to converge . These two phases are related to two different behaviors of the quantum fluctuations of the flux tube around its equilibrium position . In the strong coupling phase, these fluctuations are massive, while in the rough phase they become massless and hence survive in the continuum limit. The inverse of the mass scale of these fluctuations (which is completely different from the glueball mass scale and only appears in the model if we study the expectation value of the Wilson loop) can be considered as a new correlation length of the model. It is exactly this new correlation length which goes to infinity at the roughening point and determine the singular behaviour of the Wilson loop. This fact has several consequences:
(a) The flux–tube fluctuations can be described by a suitable two-dimensional massless quantum field theory, where the fields describe the transverse displacements of the flux tube. This quantum field theory is expected to be very complicated and will contain in general non renormalizable interaction terms . However, exactly because these interactions are non-renormalizable, their contribution becomes negligible in the infrared limit (namely for large Wilson loops). In this infrared limit the QFT becomes a conformal invariant field theory (CFT) (See e.g. chapter 9 of Ref. for a comprehensive review on CFTs).
(b) The massless quantum fluctuations delocalize the flux tube which acquires a nonzero width, which diverges logarithmically as the interquark distance increases .
(c) The quantum fluctuations give a non-zero contribution to the interquark potential, which is related to the partition function of the above $`2d`$ QFT. Hence if the $`2d`$ QFT is simple enough to be exactly solvable (and this is in general the case for the CFT in the infrared limit) also these contributions can be evaluated exactly.
(d) In the simplest case, this CFT is simply the two dimensional conformal field theory of $`(d2)`$ free bosons ($`d`$ being the number of spacetime dimensions of the original gauge model); its exact solution will be discussed in exercise 6.
#### 3.5.2 Finite Size Effects: the Lüscher term.
The feature of the effective string description which is best suited to be studied by numerical methods is the presence of finite–size effects.
Wilson loops in the confining phase are classically expected to obey the area law (see eq.(42)). This law is indeed very well verified in the strong coupling regime (before the roughening transition), but it is inadequate to describe the Wilson loop in the rough phase. In this phase the strong coupling expression must be multiplied by the partition function of the $`2d`$ QFT describing the quantum fluctuations of the flux tube. This QFT in the infrared limit becomes a 2d CFT whose partition function $`Z_q(R,T)`$ can be in some cases evaluated exactly. We shall discuss in exercise 6 an example of this type of calculations.
Eq. (42) in the rough phase becomes:
$$<W(R,T)>=e^{\sigma RT+p(R+T)+k}Z_q(R,T).$$
(52)
In general, even if one cannot give the exact expression of $`Z_q(R,T)`$ it is always possible to express its dominant contribution to the interquark potential, (i.e. in the limit $`T>>R`$) which turns out to be:
$$\underset{T\mathrm{}}{lim}\frac{1}{T}\mathrm{log}Z_q(R,T)=\frac{c\pi }{24R},$$
(53)
where $`c`$ is the central charge of the CFT. In the simplest possible case, namely when the CFT describes a collection of $`n`$ free bosonic fields, we have $`c=n`$. Thus for the free boson realization of the effective string theory, we find $`c=d2`$. This is the result obtained by Lüscher, Symanzik and Weisz in .
The interquark potential is thus given (neglecting an irrelevant constant) by:
$$V(R)=\underset{T\mathrm{}}{lim}\frac{1}{T}\mathrm{log}W(R,T)=\sigma R\frac{c\pi }{24R}.$$
(54)
The $`1/R`$ term in the potential is the finite size effect mentioned above; it is completely due to the quantum fluctuations of the flux tube and, if unambiguously detected, it represents a strong evidence (the strongest we have) in favor of the effective string picture discussed above. Moreover if the measurement is precise enough we can in principle extract numerically the value of $`c`$ and thus select which kind of effective string model describes the infrared regime of the LGT under examination.
Unfortunately, if one tries to evaluate the $`1/R`$ contribution in $`SU(2)`$ or $`SU(3)`$ gauge theories in (3+1) dimensions one faces a non trivial problem. In LGTs in (3+1) dimensions with continuous gauge groups the interquark potential has another contribution of $`1/R`$ type which has a completely different origin. It is due to the one gluon exchange. It can be evaluated perturbatively, and it exists only in the ultraviolet regime, namely for small Wilson loops. Even if it holds only in the perturbative regime, we cannot fix a sharp threshold after which it disappears, so it could well be that, in the set of large Wilson loops from which we extract our data we find a superposition of the two terms. There are two ways to avoid this problem:
* Study LGT in three dimensions where the perturbative term has a $`\mathrm{log}R`$ form instead of $`1/R`$, and does not mix up with the string contribution.
* Study Wilson loops with comparable values of $`T`$ and $`R`$. In this case, the whole functional form of the two interaction terms becomes important. These are completely different and thus can be separated.
Since the beginning of eighties several numerical works have been done to study this problem. The main results can be summarized as follows:
(a) A $`1/R`$ term exists in the potential. In the case of (3+1) LGT with continuous gauge group it can be observed also at very large distances, thus it is unlikely that it can be only due to the one gluon exchange.
(b) A similar $`1/R`$ term has been observed in various (2+1) models. In these cases the string interpretation is unambiguous.
(c) The same correction is found in very different LGTs, ranging from the $`3d`$ Ising gauge model to the $`4d`$ $`SU(3)`$ model. This remarkable universality is an important feature of these finite size effects of the effective string description.
(d) The central charge has been measured with rather good precision. The numbers are in good agreement with the $`c=d2`$ prediction of Lüscher Symanzik and Weisz for the (2+1) dimensional theories. They slightly differ in the (3+1) case. It is well possible that this is only due to the superposition of Coulomb potential.
(e) In the case of the simplest possible gauge theory, i.e. the gauge Ising model in 3 dimensions a high precision test of eq.(52) has been performed . Not only the central charge, but the whole functional form of the $`Z_q`$ correction was tested and full agreement with the effective string predictions was found.
#### 3.5.3 String Universality
We have seen (point (c) of the previous section) that the same effective string corrections have been found in all the LGT which have been studied up to now. As a matter of fact not only the string corrections, but also other features of the infrared regime of LGTs in the confining phase display a high degree of universality, namely they seem not to depend on the choice of the gauge group. This is the case for instance of the ratio between the critical temperature and the square root of the string tension, or the behavior of the spatial string tension above the deconfinement transition. All these examples show a substantial independence on the gauge group and a small and smooth dependence on the number of spacetime dimensions.
This “experimental fact” has a natural explanation in the context of an effective string model: even if in principle different gauge models could be described by different string theories, in the infrared regime, as the interquark distance increases all these different string theories flow toward the common fixed point which is not anomalous and corresponds, in the simplest case discussed above, to the two dimensional conformal field theory of $`(d2)`$ free bosons. Also the small dependence on the number of spacetime dimensions of the theory is well predicted by the effective string theory.
It is well possible that this string universality is only due to the fact that we are addressing with our simulations the simplest possible gauge theories and that looking to more complicated models a whole spectrum of effective string theories could appear, similarly to what happens in standard 2d conformal field theories. However it is interesting in this respect to notice that by resorting only to the basic property of Osterwalder-Schrader positivity (which must be true in the most general unitary LGT) one can obtain a constraint on the possible values of the coefficient of the $`1/R`$ correction (and hence of the central charge of the effective string theory). It turns out that the interquark potential must be both monotonic increasing and concave thus implying that the central charge of the effective string must be non-negative. We shall come back to this observation when discussing the AdS/CFT results for the interquark potential.
* Exercise 6: The effective string contribution to a rectangular Wilson loop.
Construct the effective string contribution to a rectangular Wilson loop (the $`Z_q(R,T)`$ term in eq.(52)) assuming a simple Nambu-Goto action for the string.
### 3.6 The string tension.
#### 3.6.1 (3+1) dimensions
The best way to discuss the present status of the lattice results on the interquark potential is to look at fig. 3 where the interquark potential for the (3+1) dimensional $`SU(3)`$ model in the quenched approximation is displayed.
The figure is taken from (to which we refer for a thorough discussion of the interquark potential) and is a a compilation of data reported in Refs .
Let us briefly comment this figure. This will also give us the opportunity to explain how LGT results are usually presented in the literature.
* Both the potential and the interquark distance are measured in units of $`r_0`$. This scale is obtained by looking at the intermediate range in the interquark potential. While the large-$`r`$ part of the potential is characterized by the string tension $`\sigma `$, one can characterize its behaviour at intermediate distances by the distance $`r_0`$ at which the force, $`F`$, has a particular value. It has become customary to use the particular definition $`r_0^2F(r_0)=1.65`$ (which corresponds to a value that can be calculated with precision on the lattice and which can be estimated with some reliability from the observed spectrum of heavy quark systems). In physical units this corresponds to $`r_00.5`$ fm.
* The important consequence of this choice is that in this way all the physical quantities ($`r`$ and $`V(r)`$) are measured in physical units and not in terms of the lattice spacing. The whole complexity of the scaling function eq.(33) is hidden in $`r_0`$ and the two combinations $`r/r_0`$ and $`V(r)r_0`$ are adimensional ratios, in the sense discussed in sect. 3.3. They are renormalization group invariant quantities and must keep the same value as the cutoff is changed (or, that is the same, as $`\beta `$ is changed) if we are in the scaling region. Thus we have an immediate and very effective test of scaling: data taken at different values of $`\beta `$ must overlap in the figure.
* This is indeed the case for the data reported in the figure which correspond to three samples of data (denoted by squares, triangles and circles respectively), obtained with MC simulations performed at three different values of $`\beta `$ (see the inset in the figure). The perfect overlap of the data is telling us that, at least for this observable, the scaling region is reached already at $`\beta =6.0`$.
* By using the scaling function (and the value $`r_00.5`$ fm) we may obtain the value, in physical units, of the lattice spacing for the three samples in the figures. They correspond to $`a0.094`$ fm, 0.069 fm and 0.051 fm, respectively. This gives an idea of the size of the “grid” of our lattice approximation.
* Looking at the figure we see that the maximum interquark distance that we one can study is about 1.5 fm (recall that we are in the quenched approximation, so the interquark string cannot break). If we tried to push the quark and antiquark pair further apart we would have to face two types of problems. First we would have to fight against increasing statistical errors (denoted in the figure by the errors bars) due to the fact that as the Wilson loops become larger and larger, since they are exponentially depressed due to the area law, the signal to noise ratio becomes smaller and smaller and too long runs are needed to obtain statistically significant results. Second, one must take into account the systematic errors due to the finite size of the lattice in which the Wilson loops are immersed. The lattice size must be much larger than the Wilson loop size to allow one to neglect these systematic errors, but again larger lattices require much more time to obtain statistical independent configurations.
* The data are fitted with the so called “Cornell potential”, which is essentially eq.(54) in which the coefficient of the $`1/R`$ term is kept as a free parameter:
$$V(r)=V_{self}+\sigma r\frac{e}{r}$$
(55)
The result of the two parameter fit<sup>4</sup><sup>4</sup>4The additive self-energy contribution, (associated with the perimeter term in the area law) is eliminated from the fit by the parametrisation-independent normalization of the data to $`V(r_0)=0`$. is plotted in the figure as a continuous line. One can directly see that the data agree very well with the proposed function. The best fit value for $`e`$ is $`e=0.295`$ which is slightly higher than the bosonic string prediction. This could be due to the interplay with the one gluon exchange contribution or to the fact that in the Cornell approximation one is neglecting the subleading (log type) contributions of the effective string. However it could also be the signature that the effective string description of the model is more complicated than the simple free bosonic model.
* From the fit we also obtain a best fit estimate for the string tension. This is the value that we shall use in the next subsection as a scale to measure the glueball masses. We report in tab. 2 the result for both $`SU(2)`$ and $`SU(3)`$ in units of $`\mathrm{\Lambda }_{\overline{MS}}`$. We also report in the same table for comparison the string tension in units of $`r_0`$ and $`T_c`$ (the deconfinement temperature).
* Looking carefully at the figure one can see that at small distances the data points lie somewhat above the curve, indicating a weakening of the effective coupling. This is a signature of the onset of asymptotic freedom at short distances.
The next step is now to study the large $`N`$ limit of the string tension. We shall address this problem in the simpler case of (2+1) dimensional theories
#### 3.6.2 (2+1) dimensions
The analysis is similar to that discussed in the previous section, but in this case it is possible to perform simulations also for larger values of $`N`$ in particular, in results for $`N=4`$ and $`N=5`$ were obtained. It turns out that these values are already large enough to perform a reliable large $`N`$ limit. Recall that in this case the coupling constant has the dimensions of a mass and thus can be used to set the mass scale of the whole theory.
It is thus natural in this case to express the string tension in units of $`g^2`$. The results for the various groups are:
$$\frac{\sqrt{\sigma }}{g^2}=\{\begin{array}{cc}0.3353(18)\hfill & \text{SU(2)}\hfill \\ 0.5530(20)\hfill & \text{SU(3)}\hfill \\ 0.7581(40)\hfill & \text{SU(4)}\hfill \\ 0.9657(54)\hfill & \text{SU(5)}\hfill \end{array}$$
(56)
It is easy to see that these values increase linearly as a function of $`N`$. This agrees with the discussion of sect. 3.4 on the large $`N`$ limit where it was shown that the natural coupling in this limit is the ’t Hooft combination $`g^2N`$.
If we try to fit the data of eq.(56) keeping also into account the first subleading term in $`1/N`$ we see that it is proportional to $`1/N^2`$. This too is a prediction of the large $`N`$ analysis which is perfectly confirmed by the simulations.
The fit with two free parameters to the equation:
$$\frac{\sqrt{\sigma }}{g^2N}=c_0+\frac{c_1}{N^2}$$
(57)
gives as a result:
$$c_0=0.1975(10),c_1=0.119(8).$$
(58)
with a very good confidence level (see for the details). The value of $`c_0`$ obtained in this way represents the first example of a non-perturbative result in the large $`N`$ limit $`SU(N)`$ gauge theories.
Let us conclude with two comments on this result
* The fit gives an acceptable confidence level even if the $`SU(2)`$ result is taken into account, this means that the large $`N`$ limit analysis (taking into account also the first $`1/N^2`$ correction) holds all the way down to $`N=2`$
* The fact that the data show the correct $`N`$ dependence is a highly non trivial test of ’t Hooft analysis, since it comes from a truly non-perturbative regularization and is completely independent from the weak coupling arguments of sect. 3.4
#### 3.6.3 The space-like string tension in finite temperature LGT.
It is important to stress that the results discussed in the two previous sections strictly refer to the zero-temperature version of $`SU(N)`$ LGT. In finite temperature LGT the interquark potential is obtained from the connected correlator of Polyakov loops. Recall that in FTLGT the lattice is asymmetric and we can distinguish between time-like and space-like Wilson loops. The spacelike Wilson loops are those orthogonal to the compactified time direction. If periodic boundary conditions in the time directions are imposed a timelike Wilson loop whose length in the $`T`$ directions is larger that the lattice size naturally becomes a pair of Polyakov loops, and this is the true order parameter for confinement. On the contrary the spacelike string tension is no longer an order parameter of the theory. In general it is different from zero even in the deconfined phase and (contrary to the naive expectation) increases as the temperature is increased! (for a discussion of this issue see for instance ). It will be important to remember this fact when we shall look at the Wilson loops in the framework of the AdS/CFT correspondence for non-supersymmetric gauge theories. It will turn out that they are actually spacelike Wilson loops of the underlying supersymmetric gauge theories. This explains why, while the supersymmetric theories, being conformally invariant, are not confining the spacelike Wilson loops (which will be interpreted as ordinary Wilson loops of the non-supersymmetric theory) are indeed confining.
### 3.7 The glueball spectrum.
#### 3.7.1 (3+1) dimensions.
In this section we summarize our present knowledge of the glueball spectrum in (3+1) dimensions in the quenched approximation from Montecarlo simulations. The quantum numbers are presented with the standard convention $`J^{PC}`$ while the asterisks refer to the radial excitations. The most precise results have been obtained in the case of the $`SU(2)`$ and $`SU(3)`$ groups. They are reported in tab. 3 for $`SU(2)`$ and in tab. 4 for $`SU(3)`$. We take this opportunity to show some of the ways in which these data are usually presented in the lattice literature. In tab. 3 the glueball masses for the $`SU(2)`$ model are reported in units of the string tension (second column) and in units of the lowest glueball (last column). Notice the absence of any $`C=`$ states in the $`SU(2)`$ case. These values are taken from ref. and the quoted errors take care both of the statistical and the systematic uncertainties.
In tab. 4 we report the glueball spectrum for $`SU(3)`$ first (in the second column) in units of the scale $`r_0`$ (see the discussion in sect. 3.6.1) and then (last column) in physical units (MeVs). These values are taken from ref. . The $`SU(3)`$ data are also plotted in fig. 4 (taken again from Ref. ). The width of the states in the figure corresponds to the combined statistical and systematic uncertainties of the estimates.
Let us stress an important non-trivial feature of the spectrum. Contrary to the naive expectations the $`2^{++}`$ state has a mass lower than that of the $`1^+`$. This is not an accident, it happens in all the model studied up to now (both in (2+1) and (3+1) dimensions, both with continuous and discrete gauge groups) and can be considered as a fingerprint of $`YM`$ theories. We shall come back to this point when discussing the glueball spectrum in the AdS/CFT framework.
Let us now study the large $`N`$ limit of the glueball spectrum. As we did in the case of the string tension we shall first address the problem in the (2+1) dimensional case where everything is much simpler and under control. We shall then try in sect. 3.7.3 below to perform the same analysis in the more interesting (3+1) dimensional case.
#### 3.7.2 (2+1) dimensions.
We report in tab. 5 the glueball spectrum in units of the square root of the string tension for $`N=2,3,4`$ and $`5`$. The table is taken from . As anticipated above we observe again the inversion between the states of the $`J=2`$ family and those of the $`J=1`$ one. The fact that in this case also $`N=5`$ data exist allows to make a reliable large $`N`$ limit analysis. Thus, following the discussion of sect. 3.6.2 let us fit these data with
$$\frac{M(J^{PC})_N}{g^2N}=M(J^{PC})_{\mathrm{}}+\frac{d_1}{N^2}$$
(59)
where we denote with $`M(J^{PC})_N`$ the mass of the glueball of quantum numbers $`J^{PC}`$ in the $`SU(N)`$ theory.
For all the $`J^{PC}`$ values the fits turn out to have good confidence levels. The large $`N`$ limit results are reported in tab. 6 (in units of the large $`N`$ string tension). In the upper part of the table we have listed the glueball states with $`C=+`$ for which also $`SU(2)`$ data exist. In the lower part of the table we report the $`C=`$ states which are obtained by fitting only data with $`N>2`$. These are the numbers with which we shall compare the AdS/CFT predictions for (2+1) dimensions which we shall discuss in the next section.
#### 3.7.3 Large $`N`$ limit in 3+1 dimensions
In order to perform a reasonable large $`N`$ limit also in (3+1) dimensions we need at least few informations also on the $`SU(4)`$ theory. In Tab. 7 are reported some of the existing data for $`SU(4)`$ taken from . They deal only with the lowest three states of the spectrum, but at least for them, they allow a tentative large N extrapolation.
In fact we see from tables 3,4 and 7 that the physical properties of $`SU(2)`$, $`SU(3)`$ and $`SU(4)`$ gauge theories are very similar. Assuming, as in the $`2+1`$ dimensional case that we are already close to the $`N=\mathrm{}`$ limit even with $`N=2,3,4`$ we can extrapolate the glueball masses (in units of $`\sqrt{\sigma }`$) to the $`N=\mathrm{}`$ limit using
$$\frac{m}{\sqrt{\sigma }}|_N=\frac{m}{\sqrt{\sigma }}|_{\mathrm{}}+\frac{c}{N^2}$$
(60)
In this way we obtain the values displayed in table 8 (see for the details). While these results are slightly less stable than the (2+1) dimensional ones they represent nevertheless the first non-perturbative results in the large $`N`$ limit of $`SU(N)`$ gauge theories in (3+1) dimensions. As such they are of the greatest importance. We shall use them to discuss the validity of the AdS/CFT approach in the (3+1) dimensional case in the next section.
## 4 AdS/CFT.
As we mentioned in the introduction we shall assume in the following that the reader is already acquainted with the theory behind the AdS/CFT correspondence. Some good reviews exist on the subject where the interested reader can find a thorough discussion of the correspondence and all the needed background material. The aim of this section is to provide the reader with the necessary information to compare the physical picture which emerges in the framework of the AdS/CFT correspondence with the results discussed in the previous sections following the lattice approach. For this reason we have organized this section in two parts: in the first one (sect. 4.1) we shall state the conjecture and discuss a few basic results (in particular on its finite temperature version) which will be needed in the following. In the second part (sect. 4.2 and 4.3) we shall review those results which are relevant for a comparison with the lattice. In particular, in sect. 4.2 we shall only deal with the finite temperature (i.e. non-supersymmetric) realization of the correspondence but, in this restricted field, we shall try to keep our review as complete as possible. In sect. 4.3 we shall mention a few results concerning the supersymmetric theory (i.e. the zero temperature case) which, due to their generality, could be (despite the presence of supersymmetry) of some importance for the comparison that we are discussing.
### 4.1 The AdS/CFT correspondence.
It is very important to stress that in going from string theory to QCD along the lines that we are discussing now, two distinct steps are needed.
The first one is the AdS/CFT correspondence, based on the Maldacena conjecture , and further specified in the works of Witten and Gubser, Klebanov and Polyakov <sup>5</sup><sup>5</sup>5In particular in the precise relation between the supergravity effective action on one side and the correlation functions of the CFT on the other side was formulated for the first time. Both the results of and those of are based on a set of earlier works .. This correspondence relates string theories on suitably chosen AdS manifolds with conformally invariant field theories whose symmetries depend on the internal manifold.
The second step is the breaking of conformal invariance and (if present) of supersymmetry, in order to obtain a candidate for a QCD-like theory. In these lectures we shall follow the proposal of Witten , in which a QCD-like theory is obtained by compactifying the original theory with suitable boundary conditions. This proposal has several appealing features and originated a large amount of papers, but it is not the unique possible choice. We shall briefly comment on this further freedom below.
Let us now discuss in more detail these two steps.
* The Maldacena conjecture relates the $`M`$ theory (or, depending on the case, one of its superstring limits) in the $`AdS_{d+1}\times 𝐗`$ background to the large $`N`$ limit of a $`d`$ dimensional conformal field theory. $`𝐗`$ is an Einstein manifold whose particular form depends on the type of field theory which we are interested to describe. In particular, by suitably choosing $`𝐗`$ it is possible to induce the presence on the field theoretic side of the correspondence of supersymmetry, or of a $`SU(N)`$ gauge symmetry. It is important to stress that, independently from the choice of the background, the resulting field theory is always conformally invariant (this explains the abbreviation AdS/CFT which is used to denote this correspondence) and, as a consequence, is not confining.
In the following we shall only discuss this correspondence in two cases
1\] In the first case we choose on the string side a type IIB superstring in the $`AdS_5\times 𝐒^\mathrm{𝟓}`$ background. The corresponding field theory turns out to be the large $`N`$ limit of the $`SU(N)`$ $`𝒩=4`$ supersymmetric gauge theory in 4 dimensions. This will be the starting point to obtain, following Witten’s suggestion a candidate for a three dimensional non-supersymmetric YM theory
2\] In the second case we choose to study, on the string side, $`M`$ theory on a $`AdS_7\times 𝐒^\mathrm{𝟒}`$ background. This theory is mapped by the Maldacena conjecture to the large $`N`$ limit of a a six dimensional $`SU(N)`$ type $`(2,0)`$ theory which is again supersymmetric and conformally invariant. By compactifying the theory in two directions according to Witten’s proposal we shall then obtain a candidate for a four dimensional non-supersymmetric YM theory
It is important to stress that this AdS/CFT correspondence is formally only a conjecture. As a matter of fact one can recognize three levels of this conjecture. Let us discuss them in the most studied example of type IIB string in the $`AdS_5\times 𝐒^\mathrm{𝟓}`$ background (case 1 in the above list)
a\] In this case the “weak” statement is that the correspondence only holds between supergravity on $`AdS_5\times 𝐒^\mathrm{𝟓}`$ and the strong coupling limit of large $`N`$ $`SU(N)`$ $`𝒩=4`$ supersymmetric gauge theory in 4 dimensions (this is the “supergravity limit” that we shall discuss below). This level of the conjecture has obtained by now so many confirmations (see for instance for a thorough discussion of all these checks) that it is commonly accepted as a firmly established result.
b\] The “normal” level of the conjecture is the one that we have stated at the beginning of this section. It extends the relation from the supergravity limit to the whole type IIB superstring in the $`AdS_5\times 𝐒^\mathrm{𝟓}`$ background, which is related with the large $`N`$ limit of the $`SU(N)`$ $`𝒩=4`$, with no constraint on the gauge coupling. This means that, with respect to the “weak” interpretation we are now pushing the correspondence outside the strong coupling limit. If we want to to obtain from the AdS/CFT correspondence a QCD-like theory in the weak coupling limit (which is the ultimate goal of these lectures), we must at least invoke this level of the conjecture. This is usually implicitly assumed in most of the papers that we shall discuss below. However essentially no check exists of the Maldacena conjecture at this level.
c\] The “strong” level consists in assuming that the conjecture holds also for the string theory at an arbitrary order in the loop expansion. From the field theory side this would imply that we have informations not only in the $`N\mathrm{}`$ limit but to any order in the $`1/N`$ expansion. As we have seen in sect. 3, this is actually not so important since we have by now a rather good control of the large $`N`$ limit on the Lattice side.
It is worthwhile to stress that (independently from the possible applications to QCD) the AdS/CFT correspondence is of great theoretical interest in itself. In some sense it represents the first nontrivial case in which we have been able to find the master field solution of a $`d>2`$ dimensional gauge theory in the large $`N`$ limit.
At the same time it is the first explicit description of a $`d>2`$ gauge theory in terms of a string theory. It seems somehow paradoxical that this remarkable result has been obtained for the first time in the case of a gauge theory which is not confining while the string description of gauge theories has been based, from the very beginning, on the intuitive picture of a string configuration spanning the minimal area of a confining Wilson loop. The mechanism behind this apparent contradiction is very instructive. The intuitive picture is indeed correct and also in the present case the string is spanning the minimal area of the Wilson loop. However due to the peculiar properties of the $`AdS`$ space the world-sheet of the string, in order to minimize its area must wander deep into the extra dimensions. This destroys the linear confining potential and leads to an effective $`1/R`$ behaviour.
* If we aim to reach a description of real QCD-like theories it is mandatory to break the conformal invariance discussed above so as to recover a well behaved confining potential. At the same time (if needed) we must somehow break the supersymmetry of the theory <sup>6</sup><sup>6</sup>6In principle the breaking of supersymmetry is not a compelling requirement, since various proposals exist for confining supersymmetric theories which are good candidates to describe the phenomenology of strong interactions. However in these lectures we shall follow a conservative attitude and look for non-supersymmetric candidates for QCD. This is almost mandatory if we want to compare the results with those obtained on the lattice where it is very difficult to implement supersymmetry (for a recent discussion of this very delicate issue see for instance ref. ).. There are several possible ways to obtain these two results and we refer to and to a discussion of these options. In the rest of these lectures we shall concentrate on the proposal suggested by Witten in . The main appealing feature of this proposal is its simplicity, however one must always keep in mind that it is not the unique possibility. In principle some of the problems that we shall discuss below and that seem to make impossible a successful comparison with standard YM theories could be avoided following other routes.
Following , we can break the conformal invariance of the theory by compactifying the theory in one (or more) direction(s). The presence of a new scale (the compactification radius) in the problem automatically breaks conformal invariance. If we then choose antiperiodic boundary conditions for the fermions in (one of) the compactified directions we also break supersymmetry. In fact as a consequence of the antiperiodic boundary conditions the gauginos and and the adjoint scalars acquire a nonzero mass. The important point in all these steps is that the Maldacena conjecture can be extended also to the compactified version of the theory thus allowing to have an insight in the infrared regime of the resulting gauge theory also in this case. It is also important to stress, so as to avoid confusion, that the theory obtained in this way is good candidate for a pure Yang Mills theory. The term QCD which is often used in the literature is from this point of view rather misleading.
Let us now study the two interesting examples of $`YM_3`$ and $`YM_4`$. We choose to study first the case of 3d $`YM`$ which, for some technical reason turns out to be simpler, we shall later generalize the results to the more interesting case of 4d $`YM`$.
#### 4.1.1 The simplest example: $`YM_3`$
In this case we must start by studying type IIB superstring in the $`AdS_5\times 𝐒^\mathrm{𝟓}`$ background. The Maldacena conjecture allows then to relate this theory with the large $`N`$ limit of the $`SU(N)`$ $`𝒩=4`$ supersymmetric gauge theory in 4 dimensions. The pattern suggested by Witten to break conformal invariance and supersymmetry is very simple in this case. Both these goals can be reached, by compactifying only one direction. There is a nice physical interpretation of this recipe. If we choose to compactify the manifold in the time direction then Witten’s proposal is equivalent to study the original system at a nonzero temperature $`T`$, proportional to the inverse of the compactification radius $`R_0`$. For this reason we shall often call in the following the original $`SYM`$ theories as $`T=0`$ theories and the non-supersymmetric compactified ones as $`T>0`$ theories. In the $`R_00`$ (hence $`T\mathrm{}`$) limit we then obtain a three dimensional effective theory which has several features in common with large $`N`$ 3d $`YM`$ and could hopefully be identified (at least in some limit) with it.
Few comments are in order at this point:
Coupling constants.
It is important to follow the coupling constant identifications that emerge from the two above steps. Let us define the coupling constant of the $`𝒩=4`$ $`SU(N)`$ theory as $`g_{YM,(𝒩=4)}^{(4)}`$. The Maldacena conjecture tells us that $`(g_{YM,(𝒩=4)}^{(4)})^2`$ is proportional to the string coupling $`g_s`$. We have seen in sect. 3.4 that the large $`N`$ limit must be taken keeping the ’t Hooft coupling $`\lambda Ng_{YM}^2`$ finite. In the present case the ’t Hooft coupling is $`\lambda N(g_{YM,(𝒩=4)}^{(4)})^2`$. The coupling constant of the three dimensional compactified theory can be expressed in terms of the four dimensional one as follows (this is a standard result in finite temperature gauge theories):
$$N(g_{YM}^{(3)})^2=\frac{N(g_{YM,(𝒩=4)}^{(4)})^2}{R_0}$$
(61)
Thus while the original coupling $`N(g_{YM,(𝒩=4)}^{(4)})^2`$ was dimensionless the 3d one has the dimensions of a mass, which completely agrees with the expected behaviour of $`YM_3`$ (see sect. 2.3). In the rest of this review we shall adopt the following convention to distinguish among the various ’t Hooft couplings. We shall denote with $`\stackrel{~}{\lambda }_d`$ the couplings which refer to the original $`T=0`$ supersymmetric $`YM`$ theories, where $`d`$ refers to the dimension of the theory, while we shall denote with $`\lambda _d`$ the coupling of the compactified non-supersymmetric theories. In this last case $`d`$ will denote the number of uncompactified dimensions. Thus in the present case:
$$\stackrel{~}{\lambda }_4N(g_{YM,(𝒩=4)}^{(4)})^2,\lambda _3N(g_{YM}^{(3)})^2.$$
(62)
So that eq.(61) becomes
$$\lambda _3=\frac{\stackrel{~}{\lambda }_4}{R_0}$$
(63)
The supergravity limit.
A crucial point for the following discussion is that we are actually unable to study the string theory on the AdS manifold in its full complexity. As a matter of fact we are bound to study the so called supergravity limit, in which the string excitations are negligible and the string theory reduces to supergravity. From the AdS/CFT correspondence one can see that this region corresponds to the $`\stackrel{~}{\lambda }_4>>1`$ regime of the $`𝒩=4`$ $`SU(N)`$ theory. In this limit the AdS/CFT correspondence is rather well understood: it essential amounts to a correspondence between supergravity fields on one side and local operators of the gauge theory on the other side. The problem is that in this limit we may only have informations on the strong coupling sector of the gauge theory. If we now move to the compactified theory we face exactly the same problem. By using supergravity we may only have informations on the strong coupling regime of the theory that we hope to identify with $`YM_3`$.
Kaluza-Klein states.
Another relevant problem is represented by the fact that in compactifying the $`S_5`$ part of the original ten dimensional supergravity on $`AdS_5\times S_5`$ a lot of Kaluza-Klein (K-K in the following) states are generated. Most of them have no counterpart in ordinary $`YM`$ theories and are expected to decouple in the $`\lambda _30`$ limit. However in the limit in which the calculations can be performed there is still no evidence of such a decoupling. We shall come back to this problem when dealing with the glueball spectrum below.
Beyond supergravity.
We have already said (and will repeat several time in the following) that in order to test in a reliable way the predictions obtained from the AdS/CFT correspondence we would need to extend them to the weak coupling regime of the theory. The quest for such an extension will appear in all the tests that we shall discuss below. However for small values of $`\stackrel{~}{\lambda }_4`$ the background geometry develops a singular behaviour and the supergravity approximation breaks down. In this regime one has to study the string theory on the AdS manifold in its full complexity. This means in particular that one should be able to study the string theory with background Ramond-Ramond charge in a singular background geometry. Despite several efforts few progress have been made in this direction up to now. However it is important to stress that this seems to be only a technical and not a conceptual obstacle and that it is well possible that in future this barrier could be overcome.
#### 4.1.2 Extension to $`YM_4`$
One can follow a procedure similar to the one outlined above to obtain a non-supersymmetric four dimensional gauge theory which could hopefully be in the same universality class of $`YM_4`$. This time one must start by looking at the $`M`$ theory in the $`AdS_7\times S_4`$ background. This theory is mapped by the Maldacena conjecture to the large $`N`$ limit of a a six dimensional $`SU(N)`$ type $`(2,0)`$ theory which is supersymmetric and conformally invariant. By compactifying the theory in two directions we then reach the desired four dimensional $`SU(N)`$ gauge theory. As a consequence of the compactification both supersymmetry and conformal invariance are lost, as it happened in the three dimensional case. However this time the relationship between the six dimensional gauge coupling and the four dimensional one is much more subtle. Let us see this correspondence in more detail. Let us denote with $`R_1`$ and $`R_2`$ the two compactification radii. In order to obtain a reduced theory with the same features of $`YM_4`$, supersymmetry must be broken only in one of the two directions, let us choose it to be the one with radius $`R_2`$, while $`R_1`$ will be the radius of the supersymmetry preserving circle. Then the four dimensional gauge coupling constant $`g_{YM}^{(4)}`$ is given by:
$$(g_{YM}^{(4)})^2=\frac{R_1}{R_2}$$
(64)
which is adimensional, as in $`YM_4`$. In order to reach the four dimensional theory that we hope to identify with $`YM_4`$ both the radii must be sent to zero, however their role is very different. Since in the large $`N`$ limit we want to keep the ’t Hooft coupling $`\lambda _4=N(g_{YM}^{(4)})^2`$ to be finite, $`R_1`$ must go to zero in the large $`N`$ limit much faster than $`R_2`$. The remaining scale $`R_2`$ plays the role of an ultraviolet cutoff for the four dimensional theory. Thus the gauge coupling $`g_{YM}^{(4)}`$ must be thought of as the bare coupling at distances of order $`R_2`$, and all dimensional quantities must be measured in units of $`R_2`$. Remarkably enough the situation is exactly the same that we have in LGT, with $`R_2`$ playing the same role of the lattice spacing in LGT.
If we aim to identify the theory that we have found with $`YM_4`$, we must require that, in the $`R_20`$ limit, $`\lambda _4`$ scales as follows:
$$\lambda _4\frac{b}{\mathrm{log}(\mathrm{\Lambda }_{QCD}R_2)}$$
(65)
with $`b`$ a suitable constant dictated by the Callan Symanzik equation.
However, exactly as in the three dimensional case discussed above we are only able to study the large $`\lambda _4`$ regime of the theory and any test of a behaviour like that of eq.(65) is well beyond our present control of theory. Again, we are bound to study the strong coupling regime of the theory.
There are at this point two possible options: the first one is to try to infer the small $`\lambda _4`$ behaviour of theory from the strong coupling informations that we have. The second is to try to extrapolate real $`YM`$ to the strong coupling limit and then compare with our findings. In both these approaches the comparison with the lattice results plays a crucial role.
### 4.2 Review of the results for the non-supersymmetric theories
All the attempts which have been made up to now to compare the results obtained in the framework of the finite temperature version of AdS/CFT correspondence with $`YM`$ theories dealt with essentially only two topics: the glueball spectrum and the string tension. The present status of these calculations is rather controversial. While a substantial agreement on the general pattern of both the glueball spectrum and the string tension has been achieved, some conflicting results still exist and the whole issue is still evolving. In particular, the agreement with the LGT results which was claimed at the beginning has been lost in the most recent analyses. For this reason we shall avoid to collect in a table a tentative list of mass values as we did when reviewing the LGT results. Instead we shall devote the next two sections to a review of the various attempts and results together with the open problems. We shall then conclude by listing the general features on which a consensus has been reached. Let us finally mention that in the following we shall write the coupling dependence of all the dimensional quantities as a function of $`\lambda _4`$ in the four dimensions case and as a function of $`\stackrel{~}{\lambda }_4`$ in the three dimensional one. The reason of this asymmetric choice is that both $`\lambda _4`$ and $`\stackrel{~}{\lambda }_4`$ are adimensional, and this greatly simplifies the discussion of the results.
#### 4.2.1 Glueball spectrum.
In principle the calculation of the glueball spectrum, at least for the $`0^{++}`$ state, is rather simple. In the supergravity limit the $`0^{++}`$ glueball is mapped by the Maldacena conjecture into the dilaton field of the corresponding supergravity description. Its mass is then obtained by solving the dilaton wave equation. This calculation was first performed in where the spectrum of the first excited states of the $`0^{++}`$ and $`0^{}`$ states in $`d=3`$ and of the $`0^{++}`$ and $`0^+`$ glueballs in $`d=4`$ was obtained. The result had the correct dependence on the ultraviolet cutoff ($`R_0`$ and $`R_2`$ respectively) and showed no explicit dependence on $`\stackrel{~}{\lambda }_4`$ or $`\lambda _4`$.
The numerical values of the lowest states turned out to be of order unity if measured in units of $`1/R_0`$ (or $`1/R_2`$). These values were then compared with LGT results <sup>7</sup><sup>7</sup>7Since there was no possibility to set the mass scale in units of some other physical quantity as in LGT (we shall discuss below the problems involved in the use of the string tension as a reference scale), the authors actually compared the ratios of higher mass glueballs with respect to the $`0^{++}`$ one. and a good agreement was claimed. However it later appeared that such a claim was probably not justified. In the mass of the $`2^{++}`$ glueball (both in $`3`$ and in $`4`$ dimensions) was obtained and turned out to be degenerate with the $`0^{++}`$ one, a result which certainly disagree with the LGT estimate reported in tab.s 6 and 8.
At the same it was realized in that in the three dimensional case the $`0^{++}`$ state associated with the graviton (in the supergravity limit) has a mass smaller than the one associated with the dilaton and hence must be considered as the lowest glueball state. If the various glueball masses are measured in units of this new fundamental mass, then the quantitative agreement with the LGT spectrum is definitely lost. However a “qualitative” agreement with lattice results is still present. In the authors also evaluated the glueball state with quantum numbers $`1^+`$ and it turns out that<sup>8</sup><sup>8</sup>8Notice that the degeneration between the $`2^{++}`$ and the $`0^{++}`$ states found in and further confirmed in refers to the “dilaton” $`0^{++}`$ glueball and not to the “graviton” one..
$$m(0^{++})<m(2^{++})<m(1^+)$$
(66)
which is exactly the same pattern which emerges in LGT. Let us stress that this is a rather non-trivial result. As mentioned in sect. 3.7 the fact that the $`2^{++}`$ state has a mass lower than that of the $`1^+`$ one is a fingerprint of QCD.
The main problem of all these calculations is the presence of the unwanted Kaluza-Klein states discussed above. It turns out that in the supergravity limit the masses of these states are of the same order of those of the glueballs. As mentioned above this is nothing else that another signature that we are actually looking at the strong coupling regime of the theory. The approach to the problem followed in the papers discussed above was to simply neglect these states assuming that they should eventually decouple if one would be able to reach the weak coupling limit. However it was noticed in that, at least in the first order in the string corrections, such a decoupling is not evident and the masses of the K-K states remain of the same order of that of the true glueball states.
It was thus proposed to study some suitable deformation of the supergravity theory which could eliminate right from the beginning these states. The idea is that there is not a unique realization of the strong coupling theory which we hope to identify with $`YM`$, but a whole family of theories which depend on one or more free parameters. Thus in principle we could tune these parameters so as make the theory as similar as possible to the weak coupling one. This is in some sense the same philosophy of the so called ”improved actions” in LGT. In this framework, the elimination of the K-K modes is certainly a step in the right direction.
This program was pursued in and more recently in in , Unfortunately only part of the K-K spectrum could be eliminated in this way. The glueball spectrum obtained in this framework depends in the most general case on three free parameters in $`d=3`$ and two parameters in $`d=4`$ , however it turns out that the mass ratios are very stable as a function of these parameters. This interesting phenomenon, which points toward the presence of some kind of universal behaviour has been discussed in .
Let us conclude with a last, positive, observation. If one were able to extrapolate these mass gap calculations up to the weak coupling limit then one should observe the scaling behaviour of eq.(33). In a first step has been made in this direction, by looking at the first string correction of the mass spectrum. The authors found that the corrections are negative. This means that, for a fixed ultraviolet cutoff the masses decrease as the ’t Hooft coupling $`\lambda _4`$ is decreased, in agreement with the expected behaviour of eq.(33).
Let us summarize the main results.
* Even if there is no quantitative agreement with the lattice estimates, the pattern of the glueball masses (at least in $`d=3`$) is correctly reproduced.
* The leading string corrections to the masses have the correct sign.
* The (dual) supergravity description of $`YM`$ can be generalized so as to incorporate two (or three) free parameters. The mass ratios show a negligible dependence on these parameters
#### 4.2.2 String Tension.
As we discussed in the introduction to LGT (sect. 3.1.1) in order to estimate the string tension one must be able to evaluate the expectation value of Wilson loops of large size. This problem is completely different from the one discussed in the previous section and requires different tools. A method to compute these Wilson loops in Super Yang Mills theories via supergravity was suggested in . These ideas were then applied to the compactified theories in which we are interested in and (as already predicted in ) a confining, linearly rising potential was indeed found. It is important at this point to recall the discussion made in sect. 3.6.3 .
The Wilson loops which are studied in (and also in all the other papers that we shall discuss in this section) are the equivalent of what we called in sect. 3.6.3 “spacelike” Wilson loops. As such they do not give informations on the potential of the original theory (which in fact, as we know, is not confining) but only on the dimensionally reduced one. Since the original theory is not confining we expect a transition (or a smooth crossover) between the two behaviours as the compactification radius is shrunk to zero (i.e. as the temperature is increased).
Let us call $`L`$ the size of the Wilson loop (i.e. the distance between the quark and the antiquark), let us study first the case of the $`d=4`$ $`𝒩=4`$ theory compactified to three dimensions. We expect that a confining potential appears in the limit $`\frac{L}{R_0}>>0`$, i.e. when the distance between the quarks is much larger than the compactification radius. In the opposite limit $`\frac{L}{R_0}<<0`$ on the contrary we expect to recover the Coulomb like behaviour of the $`𝒩=4`$ $`SYM`$ in $`d=4`$. Indeed this is exactly the behaviour which was found in . In the large $`L`$ limit it is thus possible to extract the string tension whose value turns out to be
$$\sigma =\sqrt{\pi \stackrel{~}{\lambda }_4}\frac{\pi }{R_0^2}$$
(67)
A similar analysis can be performed also in four dimensions, leading to the following expression for the string tension:
$$\sigma =\frac{8\pi }{27}\frac{\lambda _4}{R_2^2}$$
(68)
Both eq.(67) and (68) show that the string tension has the correct dimensions of $`(mass)^2`$, however its dependence on the coupling constant shows that there is a serious problem in the whole calculation. Moreover it is rather puzzling the fact that there is no signature of a $`1/L`$ type term which in LGT arises from the quantum fluctuations of the effective string. Let us discuss these two problems in more detail.
The $`\lambda `$ dependence of $`\sigma `$
Let us address this problem directly in the four dimensional case. We saw in the previous section that the lowest glueball masses were of order unity if measured in units of $`1/R_2`$ and showed no dependence on $`\lambda _4`$. On the contrary the square root of the string tension measured in units of $`1/R_2`$ is proportional to $`\sqrt{\lambda _4}`$ which must be much larger than unity in the supergravity limit, which in turn is the only regime in which we can trust this solution. This has two unwanted consequences.
First of all it completely disagrees with what we expect from both $`YM_3`$ and $`YM_4`$ in the continuum limit (see sect. 3.7) where the ratio $`\sqrt{\sigma }/M(0^{++})`$ is of order unity.
Second, it is telling us that we are actually testing the QCD string at very short distances, much shorter than the compactification radius, i.e. in a regime in which in ordinary $`YM`$ theories we do not expect to observe a “string-like” behaviour which is instead a peculiar feature of the large distance infrared regime of Wilson loops.
As stressed in the fact that $`M(0^{++})`$ and $`\sqrt{\sigma }`$ are of the same order of magnitude and have the same dependence on the bare coupling constant is an unavoidable consequence of the existence of a string-like description for the theory in which the glueballs come from closed strings. It is amazing that this property does not hold if we describe $`YM`$ theories in the framework of the AdS/CFT correspondence which is explicitly constructed to obtain a string description of $`YM`$ theories. As we have seen above, this can be interpreted as a consequence of the fact that with the AdS/CFT results we are actually probing the short range regime of the the Wilson loop and that the confining regime that we observe has little to do with the real large distance potential of the theory. Thus we may hope that also this problem will be solved when we shall be able to overcome the supergravity limit. An interesting proposal in this direction has been recently suggested in . The main point is that a log term appears in the the string tension if the corrections induced by the quantum fluctuations of the string are taken into account (we shall discuss in detail these corrections below, when dealing with the Lüscher term). The string tension becomes:
$$\sigma =\frac{8\pi }{27}\frac{\lambda _4}{R_2^2}+\frac{4\pi }{R_2^2}\mathrm{log}(R_2^2\mu ^2)+O(1/\lambda _4)$$
(69)
where $`\mu `$ is an arbitrary scale which is introduced to regulate the sum over the modes of the string fluctuations. In principle we may use this additional term to eliminate the unwanted $`\lambda _4`$ dependence in $`\sigma `$ by suitably choosing the dependence of $`\lambda _4`$ on the ultraviolet cutoff $`R_2`$. If we want to have a string tension
$$\sigma =\frac{c^2}{R_2^2}$$
(70)
with some fixed constant $`c`$, such that the ratio $`c/M(0^{++})`$ agrees with the results from LGT we must impose:
$$\lambda _4=\frac{27c^2}{8\pi }27\mathrm{log}(R_2\mu )$$
(71)
This behaviour is compatible with the other constraints on $`R_2`$ and $`\lambda _4`$ if $`\mu <<\frac{1}{R_2}`$. In this case the log term dominates over the constant and we recover the large $`\lambda _4`$ regime in which the result of eq.(68) was obtained. It is important to stress that eq.(71) imposes a constraint on the behaviour of the coupling constant $`\lambda _4`$ as a function of the ultraviolet cutoff which is different from the one due to asymptotic freedom (see eq.(65)). It can be shown that these two constraints have somehow a symmetric role. They determine the behaviour of the coupling $`\lambda _4`$ as a function of the cutoff in the strong and in the weak coupling regime respectively.
The Lüscher term
We have seen in sect. 3.5 that the area law is only the first dominant term of the potential and that besides it we expect an universal subleading correction: the Lüscher term, which is due to the quantum fluctuations of the string. In the framework of the AdS calculations that we are discussing, in the supergravity limit, there is no signature of such a term. This problem was noticed and discussed in .
Once again a possible solution to the problem is that such a term could appear if higher order string-like corrections are taken into account. This type of calculations are very delicate since they require a careful treatment of the boundary conditions for the Wilson loop which in the finite temperature case is a rather non-trivial issue (see the comments in this respect in and ). A tentative in this direction was performed in , and and a term with the desired $`\frac{1}{L}`$ behaviour was found, but with the wrong sign! <sup>9</sup><sup>9</sup>9In with a different calculation, a Lüsher term with the correct sign, was found. However in , by mimicking the type of calculation that we discussed in the section on the effective string picture in LGT (see sect. 3.5), only the quantum fluctuations of the transverse degrees of freedom were taken into account..
In view of what we were saying before, i.e. of the fact that we are actually probing the short range regime of the Wilson loops, the lack of a Lüscher term is not surprising. The same would happen also in LGT, where the $`1/L`$ correction manifests itself only at distances much larger than the ultraviolet cutoff.
However the fact that an $`1/L`$ term is present, but with the wrong sign rises a different problem, which on the contrary seems to be rather serious. We have seen at the end of sect. 3.5.3 that a very general requirement for the potential is that it must be a concave function of the interquark distance and a Lüscher term with the wrong sign violates this requirement . This problem was studied in detail in where the authors discuss which conditions must be imposed on a AdS type theory so as to fulfill the concavity requirement in the induced gauge theory.
### 4.3 A few results on the supersymmetric case.
String fluctuations in $`AdS_5\times S^5`$. As we mentioned above, the calculations on the Lüscher term discussed in the previous section are particularly delicate since they require a careful treatment of the boundary conditions for the Wilson loop. Recently some progress has been made in this respect in the zero temperature case. In particular in a careful discussion of the semiclassical fluctuations of strings in $`AdS_5\times S^5`$ based on the Green-Schwarz formalism can be found. In this paper the authors also study, among other examples, the string corrections to the expectation value of the Wilson loop in the $`𝒩=4`$ Super Yang Mills theory in $`d=4`$.
Loop equations.
We have seen in sect. 3.4 that a basic feature of the large N limit of $`SU(N)`$ gauge theories is the fact that they satisfy the so called loop equations. The solution of these loop equations would be the sought for master field of the theory. Since these equations hold also in the strong coupling phase of the theory they are a perfect testground for the validity of the conjecture. This program was recently addressed in where the zero temperature case i.e. the original $`𝒩=4`$ theory was studied<sup>10</sup><sup>10</sup>10The loop equations may be derived in a rigorous form only in the framework of the lattice discretization of the theory. Since for the moment there is no satisfactory formulation of supersymmetric theories on the lattice strictly speaking we cannot be sure that the loop equations still hold for these theories. However they can be derived, at least formally also in the continuum theory, and in this case they can be extended also to the supersymmetric case.. The loop equations have been indeed shown to hold, at least in all the cases studied by the authors.
Interaction between Wilson loops.
Let us finally mention that some interesting results have also been obtained in the study of the interaction between Wilson loops . These studies could offer new possibilities of comparison with LGT where similar studies have been also performed.
## 5 Comparison between LGT and AdS/CFT results.
We have seen in the previous sections that $`YM`$ theories regularized on the lattice have several features in common with the $`YM`$-like theories which one obtains in the framework of the AdS/CFT correspondence.
Let us summarize the results of such comparison. For each of the following items we shall first compare the AdS/CFT results with LGT in the continuum limit, then we shall also mention the difference between strong and weak coupling LGT.
* Glueball spectrum.
Both theories have a mass gap. The qualitative features of the spectrum are the same in the two theories, but there is no agreement at a quantitative level. LGT calculations in the strong coupling limit disagree with the continuum limit results (they also disagree with the AdS/CFT ones), but this disagreement becomes less significative as higher orders in the strong coupling expansion are added. There is at least one example ($`Z_2`$ theory in three dimensions) in which the expansion can be pushed to so high order that the continuum limit results are correctly recovered.
* String tension.
Both theories are confining, but the string tension which one obtains in the AdS/CFT framework disagrees in many respects with the continuum limit LGT results. First, the ratio $`\sqrt{\sigma }/M(0^{++})`$ has the wrong dependence on the coupling constant. Second the Lüscher term has (most probably) the wrong sign. A similar situation also happens if one looks at strong coupling LGT. Also in this case the ratio $`\sqrt{\sigma }/M(0^{++})`$ has the wrong scaling behaviour. The Lüscher term is exactly zero and it is well known that, as far as the string tension is concerned, the strong coupling phase is separated from the continuum limit by a phase transition: the roughening transition.
* Loop equations.
The loop equations hold both in strong coupling and weak coupling LGT. They can be defined, at least formally, also in the supersymmetric theories which appear in the AdS/CFT correspondence. A preliminary analysis shows that they hold also in this case.
* String picture.
The string description which is at the basis of the AdS/CFT approach is very different from the LGT effective string discussed in sect. 3.5 . While the first one fluctuates in the complementary space the second one originates by an attempt to describe the string fluctuations in the transverse dimensions of the physical space. However in principle it is possible that the second one could emerge as a large scale effective description from the first one. An example of such a behaviour in a different context is the 3D gauge Ising model, in which the effective string discussed in sect. 3.5 emerges at large scale from the dynamics of the Peierls contours at the level of the lattice spacing . These Peierls contours (with a suitable choice of the 3D lattice) are self-avoiding surfaces of very high genus which are probably described by some unknown string theory of which the model discussed in sect. 3.5 is a large scale effective description.
* Phase transitions.
It is by now clear that in the phase diagram of both $`SU(2)`$ and $`SU(3)`$ LGT in four dimensions with the Wilson action there is no phase transition separating the strong coupling regime from the continuum limit. For larger (but finite) values of $`N`$ or for different actions some lines of phase transitions may appear in the phase diagram, but they do not represent an obstruction to reach the continuum limit. Particular observables can undergo phase transitions (like the roughening one for the Wilson loop) in which some other correlation length of the theory (in the case of the Wilson loop the inverse of the stiffness of the surface bordered by the loop) goes to infinity without affecting the true correlation length of the theory (i.e. the inverse of the $`0^{++}`$ mass). In the large $`N`$ limit the original Eguchi-Kawai model shows a phase transition which can be avoided by introducing suitable ”twists” in the boundary conditions thus obtained the so called twisted Eguchi-Kawai model. The possible presence of phase transitions in the AdS/CFT approach is an important open problem, for which no result has been obtained up to now. The fact that the qualitative features of the glueball spectrum are correctly predicted by the theory may be considered as a hint that also in this case there is no phase transition which forbids to reach the weak coupling regime.
### 5.1 Concluding remarks
Let us try to extract the relevant outcomes of the above analysis.
It is clear that we may think of the AdS/CFT approach as a new non-perturbative regularization, alternative to the lattice, of $`YM`$ theories, in which the compactification radius in the extra dimensions plays the role of the lattice spacing $`a`$ in the lattice regularization. What is new with respect to the lattice approach is that in the AdS/CFT approach the ultraviolet cutoff, unlike the lattice spacing, does not destroy the Lorentz symmetry of the theory. On top of this, if we study the theory at the scale of the cutoff we see a higher dimensional theory, with a much larger symmetry group and a clear string interpretation. What is missing with respect to the lattice is that we lack a method to get rid of the ultraviolet cutoff and reach the weak coupling limit. This would require a better understanding of string theory on AdS manifolds, and seems for the moment a too difficult task. Any progress in this direction would play in the AdS/CFT context the same role which was played by Montecarlo simulations in LGT.
Indeed the present status of the AdS/CFT physics strongly resembles the first years of LGT, before the advent of Montecarlo simulations, when one was not even sure that the continuum limit could be reached without finding some phase transition in between, which could destroy all the nice properties found in the strong coupling limit. The danger of the possible existence of such a phase transition in the AdS/CFT approach has been stressed in .
Notwithstanding this analogies strong coupling LGT and strong coupling AdS/CFT theory show very different behaviour. This had to be expected since the fixed point, which (thanks to universality) would justify a common behaviour, is too far away. Thus it is meaningless to try to compare the two strong coupling regimes. It is much better to compare the AdS/CFT results directly with the weak coupling limit of LGT and see which of the various predictions seems to be less affected by the presence of the cutoff. In this respect the qualitative agreement of the glueball spectrum as a function of the angular momentum is a remarkable result. On the contrary, it seems that all the physics concerning the Wilson loop and the string tension is, at least at the present status of the analysis, definitely different from the weak coupling expectations.
However one should not care too much of these difficulties, because the goal is certainly worthwhile. As a matter of fact the advent of MC simulations in the lattice community, besides the obvious advantages, also had the serious drawback that people felt less urged to reach a theoretical understanding of the nonperturbative physics of LGT. In these last years progress in this direction has been much less significative than twenty years ago. In this respect the AdS/CFT approach represents a new fascinating idea and could help also people working in other areas to have a fresh look to old problems.
Acknowledgements
I am deeply indebted with M. Bianchi, M. Hasenbusch, I. Pesando, P. Provero, A. Zaffaroni and K. Zarembo for several useful suggestions and a careful reading of an earlier version of these notes. I would also like to thank M. Billó, A. D’Adda, F. Gliozzi, U. Magnea, K. Pinn and S.Vinti with whom I wrote some of the works discussed in these lectures. G. Bali and C. Morningstar are acknowledged for granting permission to reproduce their figures. This work was partially supported by the European Commission TMR programme ERBFMRX-CT96-0045.
## Appendix A Exercise Solutions.
### A.1 Exercise 1: discuss some possible generalizations of the lattice discretization of $`SU(N)`$ $`YM`$ theories.
The Wilson action can be generalized in three main directions (which may obviously be combined together):
Different representations.
In this class of generalizations the basic variable, i.e. the group element on the elementary plaquette, is unchanged, but we change the trace to a more general real function on the group. We know from group theory that if we require invariance under the gauge transformation of eq.(10) the most general function must be a linear combination of the group characters (see exercise 3 for the definitions). Hence the most general for for the action is
$$S=\underset{p}{}\underset{r}{}c_rRe\chi _r(U_p)$$
(E1.1)
where $`U_p`$ is a shorthand notation to denote the gauge variable of the $`p`$ plaquette and the $`c_r`$ are generalized coupling constants associated to the various possible representations.
Notice that at this level of generality we add no further complexity if we expand in the character basis the Boltzmann factor, thus the action $`S`$ is often presented as
$$e^S=\underset{p}{}\underset{r}{}t_r\chi _r(U_p)$$
(E1.2)
where the $`t_r`$ are in general complicated functions of the couplings $`c_r`$, but their explicit form is irrelevant and in this formulation they are usually taken as the free parameters of the theory.
Among all the possible choices of $`t_r`$ a very interesting one is the so called “heat kernel” action in which all the $`t_r`$ depend on a single parameter $`\beta `$ as follows
$$t_r=d_re^{\frac{C_r^{(2)}}{\beta N}}$$
(E1.3)
where $`d_r`$ is the dimension of the representation and $`C_r^{(2)}`$ is the quadratic Casimir invariant for the representation $`r`$. Apart from several interesting mathematical properties of this action, its major reason of interest is that for any gauge group $`SU(N)`$, in the large $`\beta `$ limit the coefficients $`D_r`$ introduced in eq.(E3.10) become equivalent to the heat kernel ones:
$$\underset{\beta \mathrm{}}{lim}D_r(\beta )=\mathrm{e}^{\frac{C_r}{2N\beta }},$$
(E1.4)
Extended plaquette actions.
The Wilson action can also be generalized by using loops larger than the elementary plaquette. To each term we may associate a coupling constant. Its reason of interest is that by suitably choosing these coupling constant one can improve the scaling behaviour of the action.
Different lattices.
Another obvious generalization is that of using different lattices, they can be both random lattices or regular lattices with different elementary Brillouin cells. In this last case the rotation symmetry is broken to subgroups different from the cubic one and, again, this can help to keep under better control the lattice artifacts.
### A.2 Exercise 2: group theoretical analysis of the glueball states for the $`SU(2)`$ LGT in $`d=3`$.
The various glueball masses are labelled by their angular momentum. Thus, in order to distinguish the various states of the spectrum one must construct operators with well defined angular momentum with respect of the two-dimensional rotations group (remind that we are interested in spacelike loops). Since we are working on a cubic lattice, where only rotations of multiples of $`\pi /2`$ are allowed, we must study the symmetry properties of our operators with respect to a finite subgroup of the two-dimensional rotations. Let us first ignore the effect of the lattice discreteness and deal with the peculiar features which, already in the continuum formulation, the (2+1) $`SU(2)`$ spectrum has with respect to the (3+1) dimensional $`SU(3)`$ spectrum.
a\] For the $`SU(2)`$ model, we cannot define a charge conjugation operator. The glueball states are thus labelled only by their angular momentum $`J`$ and by their parity eigenvalue $`P=\pm `$. The standard notation is $`J^P`$
b\] In (2+1) dimensions it can be shown that all the states with angular momentum different from zero are degenerate in parity. Namely $`J^+`$ and $`J^{}`$ (with $`J0`$) must have the same mass.
On the cubic lattice the group of two dimensional rotations and reflections becomes the dihedral group $`D_4`$. This group is non abelian, has eight elements and five irreducible representations. Four of these are one-dimensional irreps, the last one has dimension two. The group structure is completely described by the table of characters which we have reported in tab. 9 .
In the top row of tab. 9 are listed the invariant classes of the group, and in the first column the irreducible representations. We followed the notations of to label classes and representation (with the exception of the class containing the identity which we have denoted with $`\mathrm{𝟏}`$ instead of the usual $`E`$ to avoid confusion with the two-dimensional representation). The entries of the table allow to explicitly construct the various representations and hence also the lattice operators which we are looking for. The relationship of these operators with the various glueball states immediately follows from the group structure. In particular one can show that:
a\] Only operators with angular momentum $`J(mod(4))`$ can be constructed. This is a common feature of all cubic lattice regularizations. It means that glueball states which in the continuum have values of $`J`$ higher than 3 appear on the lattice as secondary states in the family of the corresponding $`J(mod(4))`$ lattice operator.
b\] The four one-dimensional irreps are in correspondence with the even $`J`$ states. More precisely:
$$0^+A_1,0^{}A_2,2^+B_1,2^{}B_2$$
This means that the discreteness of the lattice splits the degeneracy between $`2^+`$ and $`2^{}`$ which we discussed above. The splitting between these two states gives us a rough estimates of relevance of the breaking of the full rotational group due to the lattice discretization. Precise Montecarlo data have shown that in the scaling region this splitting is essentially zero within the errors, in agreement with our expectation that approaching the continuum limit the full continuum symmetries should be recovered. Notice however that this is a very non-trivial result since the operators associated to $`2^+`$ and $`2^{}`$ on the lattice turn out to be very different.
c\] All the odd parity states are grouped together in the two-dimensional irreducible representation $`E`$. This means that we cannot distinguish among them on the basis of the lattice symmetries. We can conventionally assume, say, that the $`J=1`$ states have a mass lower than the $`J=3`$ ones, and that the $`J=3`$ thus appear as secondaries in the $`J=1`$ family. In agreement with the above discussion, if the full rotational symmetry is recovered, we expect the states belonging to this family to be degenerate in parity and thus the lowest mass states, which are the ones that we can measure more precisely, to appear as a doublet. Also this prediction agrees with the data of .
The simplest lattice operators, constructed according to the character table, are shown in fig. 5.
### A.3 Exercise 3: Character Expansion for the $`SU(N)`$ group.
In this exercise we shall give some basic informations on the character expansion for $`SU(N)`$ groups, and shall then expand, as an example the Wilson action in the character basis (for further details see Ref. ). Notice that, even if we deal in particular with the SU$`(N)`$ groups, most of the results that we shall discuss can hold for any Lie group $`G`$ and with minor modifications also for discrete groups.
The irreducible characters $`\chi _r(U)`$ are the traces of the irreducible representations (labelled by $`r`$) of the group. They form a complete orthonormal basis for the class functions on the group. A function $`f(U)`$ on the group is called a “class function” if it satisfies the relation:
$$f(U)=f(VUV^{})V\mathrm{SU}(N).$$
(E3.1)
In particular, the characters themselves are class functions. The pure gauge action, eq. (13), is a class function.
The following orthogonality relations between characters hold:
$$𝑑U\chi _r(U)\chi _s^{}(U)=\delta _{r,s},$$
(E3.2)
$$\underset{r}{}d_r\chi _r(UV^{})=\delta (U,V),$$
(E3.3)
where $`dU`$ denotes the Haar measure (normalized to unity) on $`SU(N)`$ and $`d_r`$ denotes the dimension of the $`r^{\mathrm{th}}`$ representation.
Besides the above orthogonality relations there are two two other integration formulas of the characters which turn out to be very useful in the construction of SC expansions:
$$𝑑U\chi _r(V_1U)\chi _s(U^{}V_2)=\delta _{r,s}\frac{\chi _r(V_1V_2)}{d_r};$$
(E3.4)
$$𝑑U\chi _r(UV_1U^{}V_2)=\frac{1}{d_r}\chi _r(V_1)\chi _r(V_2).$$
(E3.5)
Any class function can be expanded in the basis of the characters:
$$f(U)=\underset{r}{}\chi _r(U)f_r,$$
(E3.6)
where the sum is over the set of all irreducible representations of the group, and the coefficients $`f_r`$ are given by
$$f_r𝑑U\chi _r^{}(U)f(U).$$
(E3.7)
Let us construct now the character expansion for the Wilson action.
The Boltzmann factor associated to each plaquette in the Wilson action is (see eq.(13)) :
$$\mathrm{e}^{\frac{\beta }{N}\mathrm{ReTr}U_{\mu \nu }(n)}=\underset{r}{}F_r(\beta )\chi _r(U_{\mu \nu }(n)),$$
(E3.8)
Notice that a factor 2 has been eliminated in the Boltzmann weight with respect to eq.(13) so as to avoid a double counting of the plaquettes. The coefficients $`F_r`$ are given by:
$$F_r(\beta )𝑑U\mathrm{e}^{\frac{\beta }{N}\mathrm{ReTr}U}\chi _r^{}(U)=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\mathrm{det}I_{r_jj+i+n}(\frac{\beta }{N}).$$
(E3.9)
The $`r_j`$’s are a set of integers labelling the representation $`r`$ and they are constrained by: $`r_1\mathrm{}r_N=0`$. The indices $`1i,jN`$ label the entries of the $`N\times N`$ matrix of which the determinant is taken and $`I_n(\beta )`$ denotes the modified Bessel function of order $`n`$.
As a consequence of the factor $`d_r`$ at the denominator in eq.s (E3.4, E3.5) the relevant coefficients in the character expansion (E3.8), namely the ones that will appear in the strong coupling expansions, are not the $`F_r`$ themselves, but the following normalized coefficients:
$$D_r(\beta )=\frac{F_r(\beta )}{d_rF_0(\beta )}.$$
(E3.10)
Let us see in more detail two examples: the case of $`SU(2)`$ and the large $`N`$ limit.
Character expansion for $`SU(2)`$ We can parameterize the most general matrix $`U`$ belonging to $`SU(2)`$ by using the Pauli $`\sigma `$ matrices.
$$U=cos(\theta /2)+i\stackrel{}{\sigma }\stackrel{}{n}sin(\theta /2)(0\theta <4\pi ),$$
(E3.11)
where $`\stackrel{}{n}`$ is a three dimensional normalized vector.
The normalized Haar measure is in this case
$$DU=\mathrm{sin}^2(\theta /2)\frac{d\theta }{2\pi }\frac{d^2\stackrel{}{n}}{4\pi }.$$
(E3.12)
The irreducible representations are labeled by the angular momentum $`j=0,\frac{1}{2},1,\mathrm{}`$ and have dimension $`d_j=2j+1`$.
The character of the $`j^{th}`$ irreducible representation is:
$$\chi _j(U)=\frac{\mathrm{sin}(j+\frac{1}{2})\theta }{\mathrm{sin}(\theta /2)}$$
(E3.13)
The Wilson action is in this case
$$exp\{\frac{1}{2}\beta \chi _{\frac{1}{2}}(U)\}exp\{\beta \mathrm{cos}(\theta /2)\}$$
(E3.14)
where $`U`$ is the plaquette variable.
If we insert eq.(E3.13) and (E3.14) in eq.(E3.9) we immediately recognize one of the integral representations of the modified Bessel functions.
Thus the expansion of the Wilson action in the character basis is
$$exp\{\frac{1}{2}\beta \chi _{\frac{1}{2}}(U)\}=\underset{j}{}2(2j+1)\frac{I_{2j+1}(\beta )}{\beta }\chi _j(U)$$
(E3.15)
From this one can immediately recover the expression for the normalized coefficients $`D_j(\beta )`$
$$D_j(\beta )=\frac{I_{2j+1}(\beta )}{I_1(\beta )}$$
(E3.16)
Character expansion in the large $`N`$ limit It is easy to see that in the large $`N`$ limit we find a finite value for the coefficients $`D_r(\beta )`$ only if we simultaneously take the $`\beta \mathrm{}`$ limit while keeping the $`\beta /N`$ ratio fixed (in agreement with the ’t Hooft prescription). In this limit the coefficients $`F_r(\beta )`$ turn out to have a very simple form. In particular in the region $`\beta /N<1`$ one finds:
$`F_0(\beta /N)`$ $``$ $`\mathrm{e}^{\left(\frac{\beta }{2}\right)^2},`$
$`F_f(\beta /N)`$ $``$ $`{\displaystyle \frac{\beta }{2}}\mathrm{e}^{\left(\frac{\beta }{2}\right)^2},`$
where the index $`f`$ denotes the fundamental representation (whose dimension is $`N`$). The above relations imply that in the large $`N`$ limit
$$D_f(\beta /N)=\frac{\beta }{2N}.$$
(E3.17)
Similar simplified relations hold also for higher representations.
### A.4 Exercise 4: Evaluate the first order of the strong coupling expansion of the Wilson loop in $`YM`$ theories.
Let us evaluate the first term in the strong coupling expansion the expectation value $`\chi _f(U_c)`$, where $`U_c`$ is the ordered product of gauge variables along a Wilson loop “C” of size $`R\times T`$ and $`\chi _f`$ denotes the character of the fundamental representation. Following eq.(28) the expectation value is defined as:
$$\chi _f(U_c)=\frac{_{n,\mu }dU_\mu (n)\chi _f(U_c)e^{S_W}}{Z}$$
(E4.1)
The first step is to insert in eq.(E4.1) the character expansion of the Wilson action (see eq.(E3.8)). The first non vanishing term in the expansion is the one in which we keep for all the plaquettes inside the Wilson loop (along the cristallographic plane, which ensures that we are keeping the minimum number of terms) and only for them, exactly the term in the expansion proportional to the fundamental representation. See fig. (6). In this way for all the links inside the Wilson loop and along the border we exactly find integrals of the type of eq.(E3.4).
This allow to perform all the group integrations in the expectation value, link after link. Each plaquette inside the loop gives a contribution $`F_f(\beta )`$, they are exactly $`RT`$. Each integration over the links gives a factor $`1/d_f`$. Again these are $`RT`$ (one must take into account the fact that for each integration that we perform some of the remaining links join together and thus at the end the total number of link integrals is not $`2RT`$ but only $`RT`$). Finally we must keep into account the $`Z`$ factor at the denominator of the expectation value. The simplest way to do this is to reorganize the strong coupling expansion so as to factorize also in front of the numerator the same factor $`Z`$. This simply amounts to normalize the coefficients of the expansion dividing them by $`F_0(\beta )`$. Collecting everything together we find
$$\chi _f(U_c)\left(\frac{F_f(\beta )}{d_fF_0(\beta )}\right)^{RT}D_f(\beta )^{RT}.$$
(E4.2)
This explains, by the way, why we introduced the normalized coefficients $`D_r(\beta )`$ in eq.(E3.10).
By using the definition of $`\sigma `$ (see eq.(43)) we immediately obtain from (E4.2)
$$\sigma =\mathrm{log}D_f(\beta )$$
(E4.3)
### A.5 Exercise 5: Evaluate the first order of the strong coupling expansion of the lowest glueball mass in $`YM`$ theories.
The solution of this exercise goes along the same lines of the one on the Wilson loop. The only non trivial point is that we must find the surface of minimal area bordered by the two plaquettes. If we are interested in the lowest glueball state (i.e. the $`0^{++}`$ state) we know (see sect. 3.1.2) that it is enough to study the connected correlator of two elementary spacelike plaquettes (in the fundamental representation) located at two values of the time coordinate $`t_1`$ and $`t_2`$ in the limit in which $`t|t_2t_1|\mathrm{}`$. In this limit the connected correlator decays exponentially, i.e.
$$\chi _f(U_{ij}(𝐱,t_1))\chi _f(U_{ij}(𝐱,t_2))e^{Mt}$$
(E5.1)
where $`(U_{ij}(𝐱,t_1))`$ is the plaquette (with spacelike indices $`(i,j)`$) located in the point $`(𝐱,t_1)`$ of the lattice, and $`M`$ is the mass of the lowest glueball. It is easy to see that with this geometry the minimal surface connecting the two plaquettes is a long tube made of $`4t`$ plaquettes. Hence at the first order in the strong coupling expansion we have
$$\chi _f(U_{ij}(𝐱,t_1))\chi _f(U_{ij}(𝐱,t_2))D_f(\beta )^{4t}$$
(E5.2)
from which we immediately see that
$$M(0^{++})=4\mathrm{log}(D_f(\beta ))$$
(E5.3)
### A.6 Exercise 6: The effective string contribution to a rectangular Wilson loop.
In this exercise we construct the effective string theory contribution for a Wilson loop in the infrared limit, assuming a simple Nambu-Goto action for the string. As discussed in sect. 3.5 the Nambu-Goto action reduces in this limit to the theory of $`d2`$ free massless scalar fields. In this exercise we shall compute the corresponding partition function $`Z_q(R,T)`$ which appears in eq. (52) following the discussion of ref. (and references therein).
The Nambu string action is given by the area of the world–sheet:
$$S=\sigma _0^T𝑑\tau _0^R𝑑\varsigma \sqrt{g},$$
(E6.1)
where $`g`$ is the determinant of the two–dimensional metric induced on the world–sheet by the embedding in $`R^d`$:
$`g=det(g_{\alpha \beta })`$ $`=`$ $`det_\alpha X^\mu _\beta X^\mu .`$
$`(\alpha ,\beta =\tau ,\varsigma ,\mu =1,\mathrm{},d)`$
and $`\sigma `$ is the string tension.
The reparametrization and Weyl invariances of the action (E6.1) require a gauge choice for quantization. We choose the ”physical gauge”
$`X^1`$ $`=`$ $`\tau `$
$`X^2`$ $`=`$ $`\varsigma `$ (E6.3)
so that $`g`$ is expressed as a function of the transverse degrees of freedom only:
$`g`$ $`=`$ $`1+_\tau X^i_\tau X^i+_\varsigma X^i_\varsigma X^i`$
$`+_\tau X^i_\tau X^i_\varsigma X^j_\varsigma X^j(_\tau X^i_\varsigma X^i)^2`$
$`(i=3,\mathrm{},d).`$
The fields $`X^i(\tau ,\varsigma )`$ satisfy Dirichlet boundary conditions on $`M`$:
$$X^i(0,\varsigma )=X^i(T,\varsigma )=X^i(\tau ,0)=X^i(\tau ,R)=0.$$
(E6.5)
Due to the Weyl anomaly this gauge choice can be performed at the quantum level only in the critical dimension $`d=26`$. However, the effect of the anomaly is known to disappear at large distances , which is the region we are interested in.
Expanding the square root in Eq. (E6.1) we obtain, discarding terms of order $`X^4`$ and higher
$`S`$ $`=`$ $`\sigma RT+{\displaystyle \frac{\sigma }{2}}{\displaystyle d^2\xi X^i(^2)X^i}`$ (E6.6)
$`^2`$ $`=`$ $`_\tau ^2+_\varsigma ^2.`$ (E6.7)
It is easy to see that this expansion of the action corresponds, for the partition function, to an expansion in powers of $`(\sigma RT)^1`$. Therefore the action (E6.6) describes the infrared limit of the model defined by Eq. (E6.1), and will be relevant to the physics of large Wilson loops. The contribution of the fluctuations of the flux–tube to the Wilson loop expectation value in the infrared limit will be the partition function of our CFT, given by
$$Z_q(R,T)\left[det(^2)\right]^{\frac{d2}{2}}.$$
(E6.8)
The determinant must be evaluated with Dirichlet boundary conditions.
The spectrum of $`^2`$ with Dirichlet boundary conditions is given by the eigenvalues
$$\lambda _{mn}=\pi ^2\left(\frac{m^2}{T^2}+\frac{n^2}{R^2}\right)$$
(E6.9)
corresponding to the normalized eigenfunctions
$$\psi _{mn}(\xi )=\frac{2}{\sqrt{RT}}\mathrm{sin}\frac{m\pi \tau }{T}\mathrm{sin}\frac{n\pi \varsigma }{R}.$$
(E6.10)
The determinant appearing in Eq. (E6.8) can be regularized with the $`\zeta `$-function technique: defining
$$\zeta _^2(s)\underset{mn=1}{\overset{\mathrm{}}{}}\lambda _{mn}^s$$
(E6.11)
the regularized determinant is defined through the analytic continuation of $`\zeta _^2^{}(s)`$ to $`s=0`$:
$$det(^2)=\mathrm{exp}\left[\zeta _^2^{}(0)\right].$$
(E6.12)
The series in Eq. (E6.11) can be transformed, using the Poisson summation formula, to read
$`\zeta _^2(s)={\displaystyle \frac{1}{2}}\left({\displaystyle \frac{R^2}{\pi ^2}}\right)^s\zeta _R(2s)+{\displaystyle \frac{\sqrt{\pi }Im\tau \mathrm{\Gamma }(s1/2)}{2\mathrm{\Gamma }(s)}}\left({\displaystyle \frac{R^2}{\pi ^2}}\right)^s\zeta _R(2s1)`$
$`+{\displaystyle \frac{2\sqrt{\pi }}{\mathrm{\Gamma }(s)}}\left({\displaystyle \frac{T^2}{\pi ^2}}\right)^s{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{p=1}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{\pi p}{nIm\tau }}\right)^{s1/2}K_{s1/2}(2\pi pnIm\tau )`$ (E6.13)
where $`\tau =iT/R`$, $`\zeta _R(s)`$ is the Riemann $`\zeta `$ function and $`K_\nu (x)`$ is a modified Bessel function. The derivative $`\zeta _^2^{}(s)`$ can be analytically continued to $`s=0`$ where it is given by
$$\zeta _^2^{}(0)=\mathrm{log}(\sqrt{2R})\frac{i\pi \tau }{12}\underset{n=1}{\overset{\mathrm{}}{}}\mathrm{log}(1q^n)$$
(E6.14)
where we have defined
$$qe^{2\pi i\tau }.$$
(E6.15)
Introducing the Dedekind $`\eta `$-function
$$\eta (\tau )=q^{1/24}\mathrm{\Pi }_{n=1}^{\mathrm{}}(1q^n)$$
(E6.16)
we obtain finally
$$det(^2)=\mathrm{exp}[\zeta _^2^{}(0)]=\frac{\eta (\tau )}{\sqrt{2R}}$$
(E6.17)
and
$$Z_q(R,T)\left[\frac{\eta (\tau )}{\sqrt{R}}\right]^{\frac{d2}{2}}.$$
(E6.18)
Substituting in Eq. (52) we obtain
$$<W(R,T)>=e^{\sigma RT+p(R+T)+k}\left[\frac{\eta (\tau )}{\sqrt{R}}\right]^{\frac{d2}{2}}.$$
(E6.19)
Notice, as a concluding remark, that it is clear from the above discussion that the Nambu-Goto action that we studied in this exercise is only an instance of a large class of bosonic effective string models which reduce to the CFT studied in this exercise in the infrared limit. This is one of the possible explanations for the “string universality” discussed in sect. 3.5.3.
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# Lithium-6 from Solar Flares
## 1 Introduction
The solar wind lithium isotopic ratio, (<sup>6</sup>Li/<sup>7</sup>Li)<sub>sw</sub> =0.032$`\pm `$0.004, has recently been determined from measurements in lunar soil (Chaussidon & Robert 1999). As these authors point out, this value greatly exceeds the expected photospheric ratio, based on the fact that <sup>7</sup>Li in the photosphere is depleted by over a factor of 100 relative to its protosolar value (i.e. the photospheric vs. the meteoritic abundance, Grevesse, Noels, & Sauval 1996), and that this depletion, due to burning at the bottom of the convection zone (Brun, Turck-Chieze, & Zahn 1999), should lead to a much more severe depletion of <sup>6</sup>Li, which burns at a lower temperature than <sup>7</sup>Li. In addition, there exist observational upper limits on the photospheric ratio, (<sup>6</sup>Li/<sup>7</sup>Li)<sub>ph</sub>$``$0.01 (Müller, Peytremann, & de la Reza 1975) and (<sup>6</sup>Li/<sup>7</sup>Li)<sub>ph</sub>$``$0.03 (Ritzenhoff, Schröter, & Schmidt 1997). Chaussidon & Robert (1999) thus suggest that the measured solar wind <sup>6</sup>Li must be solar flare produced. However, they only consider <sup>6</sup>Li production by spallation from C, N and O. The demonstration that solar flares can indeed account for the <sup>6</sup>Li in the solar wind has very important implications on many problems in solar physics.
Light element production by accelerated particle interactions was treated in detail (e.g. Ramaty et al. 1997). In non-solar settings, and for accelerated particles of predominantly low energy, the dominant reactions are <sup>4</sup>He($`\alpha `$,p)<sup>7</sup>Li, <sup>4</sup>He($`\alpha `$,n)<sup>7</sup>Be (with <sup>7</sup>Be decaying to <sup>7</sup>Li) and <sup>4</sup>He($`\alpha `$,x)<sup>6</sup>Li (where x stands for either a proton and a neutron, or a deuteron). In solar flares, however, the reaction <sup>4</sup>He(<sup>3</sup>He,p)<sup>6</sup>Li is also very important (Mandzhavidze, Ramaty, & Kozlovsky 1997a), both because of its very low threshold energy and because for solar energetic particles <sup>3</sup>He/<sup>4</sup>He can be as large as 1 or even larger (e.g. Reames 1998). Such <sup>3</sup>He/<sup>4</sup>He enhancements are one of the main characteristics of the acceleration mechanism responsible for impulsive solar energetic particle events, as distinguished from gradual events, based on the duration of the accompanying soft X-ray emission. The <sup>3</sup>He enrichment is thought to be due to stochastic acceleration through gyroresonant wave particle interactions which preferentially accelerate the <sup>3</sup>He (Temerin & Roth 1992; Miller & Viñas 1993). Concerning the particles which interact at the Sun, evidence for accelerated <sup>3</sup>He enrichment was obtained from the detection (Share & Murphy 1998) of a gamma-ray line at 0.937 MeV produced by the reaction <sup>16</sup>O(<sup>3</sup>He,p)<sup>18</sup>F (Mandzhavidze, Ramaty, & Kozlovsky 1997b; 1999). Using gamma-ray data from 20 flares, Mandzhavidze et al. (1999) showed that for essentially all of these flares <sup>3</sup>He/<sup>4</sup>He can be as large as 0.1, while for some of them values as high as 1 are possible. In addition, they showed that for the particles that interact and produce gamma rays, <sup>3</sup>He enrichments are present for both impulsive and gradual flares. Thus, we can expect <sup>3</sup>He/<sup>4</sup>He$`\stackrel{>}{}`$0.1 for most flares that produce gamma rays and isotopes at the Sun.
In the present Letter we carry out new calculations of Li production and re-calculate (see Ramaty & Simnett 1991) the average accelerated ion irradiation of the Sun, to show that flare accelerated particle interactions produce enough <sup>6</sup>Li which, combined with photospheric <sup>7</sup>Li, can account for the solar wind <sup>6</sup>Li/<sup>7</sup>Li measured in lunar soil.
## 2 Li Production
We employ the nuclear code described in detail in Ramaty et al. (1997) which includes, in addition to the $`\alpha `$$`\alpha `$ reactions mentioned above, also Li production from C, N and O. The cross section for the additional reaction, <sup>4</sup>He(<sup>3</sup>He,p)<sup>6</sup>Li, is shown in Figure 1, together with the cross sections for the $`\alpha `$$`\alpha `$ reactions producing <sup>6</sup>Li and <sup>7</sup>Li. For the <sup>3</sup>He induced reaction we obtained the cross section for <sup>6</sup>Li production in the ground state, from threshold (2.34 MeV/nucleon) to 8.2 MeV/nucleon, by detailed balance using the cross section for the inverse exothermic reaction <sup>6</sup>Li(p,<sup>3</sup>He)<sup>4</sup>He (Angulo et al. 1999). We added the contribution of the reaction for producing <sup>6</sup>Li in the 3.56 MeV excited state which decays to the ground state by photon emission, using data from Harrison (1967). The total cross section at 9.3 MeV/nucleon is from Koepke and Brown (1977), and at 18 and 20.4 MeV/nucleon from Halbert, van der Woude, & O’Fallon (1973). At higher energies we extrapolated the cross section as expected for reactions with 2 particles in the exit channel.
Gamma-ray production in solar flares results predominantly from thick target interactions, meaning that particles accelerated in the upper portions of coronal loops produce nuclear reactions as they slow down in the denser chromospheric region of the loops (e.g. Ramaty & Murphy 1987). We adopt the same model for Li production. The upper panel in Figure 2 shows the resultant thick target <sup>6</sup>Li yields, normalized to unit incident total number of protons of energy greater than 30 MeV, $`N_\mathrm{p}(>`$30)=1. The energy spectra of the accelerated particles are power laws in kinetic energy per nucleon, with spectral index $`s`$ (Ramaty, Mandzhavidze, & Kozlovsky 1996). The evidence for enhanced <sup>3</sup>He/<sup>4</sup>He was mentioned above. There is also evidence that $`\alpha `$/p could exceed the canonical 0.1, with possible value around 0.5 (Share & Murphy 1997; Mandzhavidze et al. 1999). Thus in Figure 2 we show results for $`\alpha `$/p = 0.1 and 0.5, and <sup>3</sup>He/<sup>4</sup>He=0, 0.1 and 1. We see in the upper panel that the <sup>3</sup>He enrichment very significantly increases the lithium production, especially for steep spectra. That the lithium production is mainly due to $`\alpha `$ particles and <sup>3</sup>He nuclei can be seen by comparing the six upper curves with the lowest one, for which we set the $`\alpha `$ particle and <sup>3</sup>He abundances to zero, so that all the <sup>6</sup>Li in this case is produced in C, N and O interactions. Considering the flare produced isotopic ratios in the lower panel, we see that while in the absence of <sup>3</sup>He, <sup>6</sup>Li/<sup>7</sup>Li is at most unity, much larger ratios are possible with enhanced <sup>3</sup>He/<sup>4</sup>He.
## 3 Average Solar Proton Irradiation
To calculate the average flare produced lithium, we estimate the average proton irradiation of the Sun, $`\dot{N}_\mathrm{p}`$($`>`$30MeV) measured in protons per second, where the average is taken over a solar cycle. We follow the method described by Ramaty & Simnett (1991). We start with the flare size distribution measured in 0.3 to 1 MeV bremsstrahlung because observations in this energy range give the most complete sample of solar flare gamma-ray emission (see Vestrand et al. 1999). To minimize the effects of anisotropic electrons (e.g. Miller & Ramaty 1989) we employ the distribution derived for flares near the solar limb (Dermer 1987). For flares at heliocentric longitudes 60 to 90, observed from March 1980 to February 1986 (approximately half a solar cycle) the size distribution, measured in number of flares per unit $`F_\mathrm{B}`$, can be approximated by $`dn/dF_\mathrm{B}8.5F_\mathrm{B}^{1.1}`$, where 10$`\stackrel{<}{}F_\mathrm{B}\stackrel{<}{}`$6500 photons cm<sup>-2</sup> is the observed 0.3 to 1 MeV bremsstrahlung fluence at Earth per flare . The total number of $`>`$0.3 MeV emitting flares per solar cycle is obtained by integrating the above expression multiplied by a factor of 12, where a factor of 6 takes into account the whole solar surface and a factor of 2 the other half of solar cycle. We thus obtain 375 flares, which compares well with the 175 flares listed by Vestrand et al. (1999) from which $`>`$0.3 MeV bremsstrahlung was observed with the Solar Maximum Mission (SMM) over almost a whole solar cycle. This latter number should be corrected for anisotropy effects, and must be multiplied by a factor of 2 since SMM only observes half the solar surface. The required average irradiation is then given by
$`\dot{N}_\mathrm{p}(>30)={\displaystyle \frac{12}{T}}{\displaystyle _{10}^{6500}}𝑑F_\mathrm{B}{\displaystyle \frac{dn}{dF_\mathrm{B}}}N_\mathrm{p}(F_\mathrm{B}),`$ (1)
where T is the number of seconds in 11 years and $`N_\mathrm{p}(F_\mathrm{B})`$ is the number of protons above 30 MeV expressed as a function of $`F_\mathrm{B}`$. To derive this relationship, we first employ the result of Murphy et al. (1990) that for flares near the limb $`F_\mathrm{B}/F_\mathrm{N}4.5`$, where $`F_\mathrm{N}`$ is the total nuclear deexcitation line emission fluence observed at Earth. Next we use the nuclear deexcitation code (e.g. Ramaty et al. 1996) to derive $`N_\mathrm{P}(>`$30)/$`F_\mathrm{N}`$. This ratio depends on the spectrum and composition of the accelerated particles, in particular $`\alpha `$/p. Ramaty et al. (1996) have derived the distribution of power law spectral indexes from gamma-ray data, showing that for a sample of 19 flares the mean $`s4`$. For this value of $`s`$ we find that $`N_\mathrm{p}(>`$30)/$`F_\mathrm{N}`$ = 1.7$`\times `$10<sup>29</sup> and 6.6$`\times `$10<sup>29</sup> protons/(nuclear deexcitation photons cm<sup>-2</sup>), for $`\alpha `$/p=0.5 and 0.1, respectively. By using $`F_\mathrm{B}/F_\mathrm{N}`$=4.5 and these $`N_\mathrm{p}(>`$30)/$`F_\mathrm{N}`$ to derive $`N_\mathrm{p}(F_\mathrm{B})`$, equation 1 yields $`\dot{N}_\mathrm{p}`$($`>`$30MeV)=3.5$`\times `$10<sup>25</sup> and 1.4$`\times `$10<sup>26</sup> protons s<sup>-1</sup>, for $`\alpha `$/p=0.5 and 0.1, respectively.
## 4 The Solar Wind <sup>6</sup>Li/<sup>7</sup>Li
Even though a detailed treatment of the time dependent evolution of Li in the solar atmosphere is beyond the scope of this paper, we now show that <sup>6</sup>Li production in solar flares could indeed account for the solar wind <sup>6</sup>Li/<sup>7</sup>Li. To demonstrate this we assume the following: (i) all the flare produced <sup>6</sup>Li is evacuated by the solar wind, (ii) the photospheric <sup>6</sup>Li that is the remnant of its protosolar abundance is negligible, and (iii) the solar wind (<sup>7</sup>Li/H)<sub>sw</sub> is equal to the photospheric value (<sup>7</sup>Li/H)<sub>ph</sub>=1.4$`\times `$10<sup>-11</sup> (Grevesse et al. 1996). The solar wind (<sup>6</sup>Li/<sup>7</sup>Li)<sub>sw</sub> is then given by
$`\left({\displaystyle \frac{{}_{}{}^{6}\mathrm{Li}}{{}_{}{}^{7}\mathrm{Li}}}\right)_{\mathrm{sw}}={\displaystyle \frac{\dot{N}_\mathrm{p}(>30)Q(^6\mathrm{Li})/\dot{F}_{\mathrm{sw}}}{(^7\mathrm{Li}/\mathrm{H})_{\mathrm{ph}}}},`$ (2)
where $`Q(^6\mathrm{Li})`$ is plotted in Figure 2 and $`\dot{F}_{\mathrm{sw}}`$$``$6$`\times `$10<sup>35</sup> s<sup>-1</sup> is the average solar wind proton flux (Dupree 1996). Taking $`s=4`$, 0.1$`<`$$`\alpha `$/p$`<`$0.5 and 0.1$`<`$<sup>3</sup>He/<sup>4</sup>He$`<`$1, we obtain 0.007$`<`$(<sup>6</sup>Li/<sup>7</sup>Li)<sub>sw</sub>$`<`$0.06. This range is consistent with the observed value of 0.032$`\pm `$0.004. Several effects could lead to lower or higher calculated (<sup>6</sup>Li/<sup>7</sup>Li)<sub>sw</sub>. Clearly there are uncertainties in our estimate of $`\dot{N}_\mathrm{p}(>`$30 MeV), in particular there could be a large number of smaller gamma-ray flares, which have not yet been observed, and if they had steep ion energy spectra and high <sup>3</sup>He/<sup>4</sup>He they would contribute significantly to <sup>6</sup>Li production. On the other hand, some of the flare-produced <sup>6</sup>Li could be mixed downward to the photosphere and lost from the solar wind. The calculated (<sup>6</sup>Li/<sup>7</sup>Li)<sub>sw</sub> would also be lower if (<sup>7</sup>Li/H)<sub>sw</sub> were higher than (<sup>7</sup>Li/H)<sub>ph</sub>, a possibility since Li has low first ionization potential, a factor that biases coronal abundances relative to those of the photosphere (e.g. Reames 1998). Nevertheless, the better than order of magnitude agreement between the calculated and measured (<sup>6</sup>Li/<sup>7</sup>Li)<sub>sw</sub> provides good support to the possibility that the measured <sup>6</sup>Li in lunar soil is indeed solar flare produced.
It is of some interest to compare the average <sup>6</sup>Li production, 6$`\times `$10<sup>22</sup>$`<`$\[$`\dot{N}_\mathrm{p}`$$`(>`$30)$`Q(^6`$Li)\] $`<`$5$`\times `$10<sup>23</sup> atoms s<sup>-1</sup>, with the contribution of the 19 large SMM flares from which gamma-ray line emission was observed. Using the method detailed in Mandzhavidze et al. (1999), for each flare we derive $`s`$ and $`N_\mathrm{p}`$($`>`$30). Then using $`Q`$(<sup>6</sup>Li) from Figure 2, taking into account that for 5 of the 19 flares $`\alpha `$/p$``$0.5 (Mandzhavidze et al. 1999), we obtain the flare-by-flare <sup>6</sup>Li productions which yield averages over the 9 year SMM observing period of 1$`\times `$10<sup>22</sup> and 7$`\times `$10<sup>22</sup> <sup>6</sup>Li atoms s<sup>-1</sup>, if for all flares <sup>3</sup>He/<sup>4</sup>He=0.1 and 1, respectively. Thus, about 15% of the <sup>6</sup>Li production that we derived using the flare size distribution could result from 19 of the largest flares. Concerning the contributions of individual flares, as much as a few times 10<sup>30</sup> Li atoms could be produced by a large flare and most of these would be <sup>6</sup>Li (Mandzhavidze et al. 1997a)
## 5 Discussion and Conclusions
We demonstrated that it is possible to produce enough <sup>6</sup>Li by flare accelerated particles to account for the measured <sup>6</sup>Li/<sup>7</sup>Li in lunar soil that is thought to originate from solar wind implantation. The presence of enriched accelerated particle <sup>3</sup>He is essential for the production of sufficient <sup>6</sup>Li. We note that the radioactive <sup>26</sup>Al in the early solar system is thought to be produced in <sup>3</sup>He induced reactions (Lee et al. 1998). This raises the possibility that some of the meteoritic <sup>6</sup>Li could also be of local early solar system origin.
Kotov et al. (1996) claimed that flare accelerated particle interactions could account for all the photospheric lithium. If this were true, since the solar wind acceleration is not expected to significantly alter the lithium isotopic ratio, the solar wind <sup>6</sup>Li/<sup>7</sup>Li should exceed 0.2 (Figure 2), contrary to the observed value of 0.03. This confirms the previous result of Mandzhavidze et al. (1997a) that production in flares does not make a significant contribution to the average photospheric lithium. But the fact that as much as 10<sup>30</sup> Li atoms are produced in large solar flares, suggests that flare produced lithium may be detected in a small area of the solar surface near the foot points of the flaring loops shortly after the time of the flare (see Livshits 1997). In this connection, it is interesting to point out that Ritzenhoff et al. (1997) don’t rule out the presence of <sup>6</sup>Li near a sunspot at a value close to their reported upper limit <sup>6</sup>Li/<sup>7</sup>Li$``$0.03, which in fact coincides with the measured solar wind value.
Further research in this area requires direct measurement of lithium and its isotopic ratio in the solar wind, spectroscopic measurements of <sup>6</sup>Li in the photosphere, and the detection of gamma rays from small flares that would lead to a more precise determination of the proton irradiation of the Sun. All of these should lead to new insights into the processes of transport and mixing in the solar atmosphere and of the acceleration of the solar wind.
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# Gröbner Basis Techniques for Computing Actions of K-Categories KEYWORDS: Gröbner basis, K-category, Action, Kan extension.
## 1 Introduction
The paper is motivated by a question which arises from two pieces of research. Firstly, the work of Brown and Heyworth which extends rewriting techniques to enable the computation of left Kan extensions over the category of sets. It is well known that left Kan extensions can be defined over categories other than $`\mathrm{𝖲𝖾𝗍𝗌}`$. Secondly, the ‘folklore’, made explicit in that rewriting theory is a special case of noncommutative Gröbner basis theory. It is therefore natural to ask whether Gröbner bases can provide a method for computing Kan extensions beyond the special case of rewriting.
To answer this question completely, fully exploiting the computational power of Gröbner basis techniques relating to Kan extensions is the ultimate aim. This paper provides a first step by showing how standard noncommutative Gröbner basis procedures can be used to calculate left Kan extensions of $`K`$-category actions. In the final section of the paper a number of interesting problems arising from the work are identified.
## 2 Background
This paper builds on work of Brown and Heyworth on extensions of rewriting methods.
The standard expression of rewriting is in terms of words $`w`$ in a free monoid $`\mathrm{\Delta }^{}`$ on a set $`\mathrm{\Delta }`$. This may be extended to terms $`x|w`$ where $`x`$ belongs to a set $`X`$ and the link between $`x`$ and $`w`$ is in terms of an action. More precisely, we suppose a monoid $`A`$ acts on the set $`X`$ on the right, and there is given a morphism of monoids $`F`$: $`AB`$ where $`B`$ is given by a presentation with generating set $`\mathrm{\Delta }`$. The result of the rewriting will then be normal forms for the induced action of $`B`$ on $`F_{}(X)`$. This gives an important extension of rewrite methods.
In fact monoids may be replaced by categories, and sets by directed graphs. This gives a formulation in terms of left Kan extensions, or induced actions of categories, which is explained in . Further, categories can be replaced by $`K`$-categories, as will be described later.
Let $`𝖠`$ be a category. A category action $`X`$ of $`𝖠`$ is a functor $`X:𝖠\mathrm{𝖲𝖾𝗍𝗌}`$. Let $`𝖡`$ be a second category and let $`F:𝖠𝖡`$ be a functor. Then an extension of the action $`X`$ along $`F`$ is a pair $`(E,\epsilon )`$ where $`E:𝖡\mathrm{𝖲𝖾𝗍𝗌}`$ is a functor and $`\epsilon :XEF`$ is a natural transformation. The left Kan extension of the action $`X`$ along $`F`$ is an extension of the action $`(E,\epsilon )`$ with the universal property that for any other extension of the action $`(E^{},\epsilon ^{})`$ there exists a unique natural transformation $`\alpha :EE^{}`$ such that $`\epsilon ^{}=\alpha \epsilon `$.
The problem that has been introduced is that of “computing a Kan extension”. In order to keep the analogy with computation and rewriting for presentations of monoids we propose a definition of a presentation of a left Kan extension. The papers were very influential on our choices.
A left Kan extension data $`(X^{},F^{})`$ consists of small categories $`𝖠`$, $`𝖡`$ and functors $`X^{}:𝖠\mathrm{𝖲𝖾𝗍𝗌}`$ and $`F^{}:𝖠𝖡`$. A left Kan extension presentation is a quintuple $`𝒫:=kan\mathrm{\Gamma }|\mathrm{\Delta }|RelB|X|F`$ where $`\mathrm{\Gamma }`$ and $`\mathrm{\Delta }`$ are (directed) graphs; $`X:\mathrm{\Gamma }\mathrm{𝖲𝖾𝗍𝗌}`$ and $`F:\mathrm{\Gamma }P\mathrm{\Delta }`$ are graph morphisms to the category of sets and the free category on $`\mathrm{\Delta }`$ respectively; and $`RelB`$ is a set of relations on the free category $`P\mathrm{\Delta }`$.
Formally, we say $`𝒫`$ presents the left Kan extension $`(E,\epsilon )`$ of the left Kan extension data $`(X^{},F^{})`$ where $`X^{}:𝖠\mathrm{𝖲𝖾𝗍𝗌}`$ and $`F^{}:𝖠𝖡`$ if $`\mathrm{\Gamma }`$ is a generating graph for $`𝖠`$ and $`X:\mathrm{\Gamma }\mathrm{𝖲𝖾𝗍𝗌}`$ is the restriction of $`X^{}:𝖠\mathrm{𝖲𝖾𝗍𝗌}`$; $`cat\mathrm{\Delta }|RelB`$ is a category presentation for $`𝖡`$ and $`F:\mathrm{\Gamma }P\mathrm{\Delta }`$ induces $`F^{}:𝖠𝖡`$.
We expect that a left Kan extension $`(E,\epsilon )`$ is given by a set $`EB`$ for each $`B\mathrm{Ob}\mathrm{\Delta }`$ and a function $`Eb:EB_1EB_2`$ for each $`b:B_1B_2𝖡`$ (defining the functor $`E`$) together with a function $`\epsilon _A:XAEFA`$ for each $`A\mathrm{Ob}𝖠`$ (the natural transformation).
The main result of defines rewriting procedures on $`T:=_{B\mathrm{Ob}\mathrm{\Delta }}_{A\mathrm{Ob}\mathrm{\Gamma }}XA\times P\mathrm{\Delta }(FA,B)`$ which is basically a set with a partial right action of the arrows of $`P\mathrm{\Delta }`$.
Two kinds of rewriting are involved here. The first is the familiar $`x|ulvx|urv`$ given by a relation $`(l,r)`$ – these rules are known as the ‘$`E`$-rules’. The second derives from a given action of certain words on elements, so allowing rewriting $`x|F(a)vxa|v`$ – these rules are known as the ‘$`\epsilon `$-rules’. Further, the elements $`x`$ and $`xa`$ may belong to different sets. When such rewriting procedures complete, the associated normal form gives in effect a computation of what we call the Kan extension defined by the presentation.
The ‘folklore’ of the relation of rewriting and Gröbner basis techniques, alluded to in and is made explicit in .
The polynomial ring $`K[X^{}]`$ consists of all polynomials having coefficients in the field $`K`$ and terms from $`X^{}`$ together with the usual operations of polynomial addition and (noncommutative) multiplication. Given a generating set for an ideal $`I`$ in this ring it is a problem to determine whether two given polynomials $`f`$ and $`g`$ are equivalent modulo the ideal i.e. whether they occur within the same congruence class. If a Gröbner basis $`G`$ can be constructed for $`I`$ from the original generating set then the congruence problem can be solved. The Gröbner basis calculation depends on a well-ordering of $`X^{}`$ and a definition of polynomial reduction, which is determined by comparing leading terms. In the noncommutative case it is not always successful.
The key observation is that the rewriting techniques used in calculating a monoid $`M`$ from a set of generators $`X`$ and a rewrite system $`R`$ compatible with an ordering $`>`$ corresponds step-by-step to the Gröbner basis techniques used in calculating the congruence classes of the polynomial ring $`K[X^{}]`$ with respect to the ideal generated by the difference binomials $`lr`$ for $`(l,r)`$ in $`R`$.
This provides the background to our problem of determining whether Gröbner bases can be used to calculate Kan extensions other than in the special case of rewriting systems. The first observation is that Gröbner bases involve polynomials, so we should examine how the addition operation is represented in categories.
## 3 $`K`$-Category Actions
We use the definitions of . Let $`K`$ be a field. A $`K`$-category is a category whose hom-sets (a hom-set is the set of all morphisms between a given pair of objects) are $`K`$-modules. A morphism of $`K`$-categories or $`K`$-functor $`F`$ preserves the $`K`$-module structure of the hom-sets so $`F(a+b)=F(a)+F(b),F(ka)=kF(a)`$ for all arrows $`a,b`$ such that $`a+b`$ is defined and scalars $`k`$ in $`K`$.
The free $`K`$-category on the graph $`\mathrm{\Delta }`$ is the category $`P_K\mathrm{\Delta }`$ whose objects are the objects of $`\mathrm{\Delta }`$ and whose arrows $`\mathrm{Arr}P_K\mathrm{\Delta }`$ are all polynomials of the form $`p=k_1m_1+\mathrm{}+k_nm_n`$ where $`k_1,\mathrm{},k_nK`$ and $`m_1,\mathrm{},m_nP\mathrm{\Delta }(B_1,B_2)`$ for some $`B_1,B_2\mathrm{Ob}\mathrm{\Delta }`$. We will refer to $`m_1,\mathrm{},m_n`$ as the terms which occur in $`f`$. Note that functions $`src`$ and $`tgt`$ are well-defined as $`src(f):=src(m_1)=\mathrm{}=src(m_n)`$ and $`tgt(f):=tgt(m_1)=\mathrm{}=tgt(m_n)`$.
The relations of a $`K`$-category can be of the form $`p=q`$ where both sides have the same source and target. Therefore $`R`$ will be assumed to be a set of polynomials $`pq`$ i.e. a subset of $`\mathrm{Arr}P_K\mathrm{\Delta }`$. If $`R=\{r_1,\mathrm{},r_n\}`$ is such a set of relations on $`P_K\mathrm{\Delta }`$ then the congruence generated by $`R`$ is defined as follows:
$$f=_Rh\text{ if and only if }f=h+k_1p_1r_1q_1+\mathrm{}+k_np_nr_nq_n$$
for some $`k_1,\mathrm{},k_nK`$ and $`p_1,\mathrm{},p_n,q_1,\mathrm{},q_n\mathrm{Arr}P_K\mathrm{\Delta }`$ where $`src(f)=src(h)=src(u_1)=\mathrm{}=src(u_n)`$ and $`tgt(f)=tgt(h)=tgt(v_1)=\mathrm{}=tgt(v_n)`$ and $`u_1r_1v_1,\mathrm{},u_nr_nv_n`$ are defined in $`\mathrm{Arr}P_K\mathrm{\Delta }`$. The $`K`$-category $`P_K\mathrm{\Delta }/=_R`$ whose elements are the congruence classes of $`\mathrm{Arr}P_K\mathrm{\Delta }`$ with respect to $`F`$ is known as the factor $`K`$-category.
###### Definition 3.1
Let $`K`$ be a field. A $`K`$-category presentation is a pair $`cat_K\mathrm{\Delta }|R`$ where $`\mathrm{\Delta }`$ is a graph and $`R\mathrm{Arr}P_K\mathrm{\Delta }`$. The $`K`$-category it presents is the factor category $`P_K\mathrm{\Delta }/=_R`$.
Our first result enables the use of Buchberger’s algorithm to compute Gröbner bases which enable the specification of the morphisms of a $`K`$-category presented in this way.
Let $`>`$ be an admissible well-ordering on $`\mathrm{Arr}P\mathrm{\Delta }`$ i.e. $`>`$ is Noetherian and compatible with the operation of path concatenation. Define the leading term of a polynomial $`f`$ to be the term occurring in $`f`$ which is the greatest path in $`\mathrm{\Delta }`$ with respect to $`>`$ and denote it $`\mathrm{𝙻𝚃}(f)`$. Define a reduction relation $`_R`$ on $`\mathrm{Arr}P_K\mathrm{\Delta }`$ by $`ffk_iu_ir_iv_i`$ when $`u_i(\mathrm{𝙻𝚃}(r_i))v_i`$ occurs in $`f`$ with coefficient $`k_iK`$ for $`u_i,v_i\mathrm{Arr}P\mathrm{\Delta }`$, $`r_iR`$. The reflexive, symmetric and transitive closure of $`_R`$ is denoted $`\underset{R}{\overset{}{}}`$. If the reduction relation $`_R`$ is complete (i.e. Noetherian and confluent) then we say that $`R`$ is a Gröbner basis.
###### Lemma 3.2
$$\frac{\mathrm{Arr}P_K\mathrm{\Delta }}{=_R}\frac{\mathrm{Arr}P_K\mathrm{\Delta }}{\underset{R}{\overset{}{}}}$$
Proof It is clear from the definitions that the equivalence relation $`\underset{R}{\overset{}{}}`$ is contained in $`=_R`$.
For the converse, suppose $`f=_Rh`$. Then there exist $`r_1,\mathrm{},r_nR`$ and $`p_1,\mathrm{},p_n,q_1,\mathrm{},q_nP_K\mathrm{\Delta }`$, such that $`f=h+p_1r_1q_1+\mathrm{}+p_nr_nq_n`$. By splitting $`p_i`$ and $`q_i`$ into their component terms for $`i=1,\mathrm{},n`$ we obtain $`f=h+k_1u_1r_1v_1+\mathrm{}+k_ju_jr_iv_j+\mathrm{}+k_tu_tr_nv_t`$ for some $`k_1,\mathrm{},k_tK`$, $`u_1,\mathrm{},u_t,v_1,\mathrm{},v_tP\mathrm{\Delta }`$. It follows immediately from this that $`f\underset{R}{\overset{}{}}h`$. $`\mathrm{}`$
###### Proposition 3.3
The relation $`_R`$ is Noetherian on $`\mathrm{Arr}P_K\mathrm{\Delta }`$.
The matches of $`R`$ are the pairs of polynomials $`(r_1,r_2)`$ whose leading terms overlap on some subword i.e. $`u\mathrm{𝙻𝚃}(r_1)v=\mathrm{𝙻𝚃}(r_2)`$ or $`\mathrm{𝙻𝚃}(r_2)=u\mathrm{𝙻𝚃}(r_2)v`$ or $`u\mathrm{𝙻𝚃}(r_1)=\mathrm{𝙻𝚃}(r_2)v`$ or $`\mathrm{𝙻𝚃}(r_1)v=u\mathrm{𝙻𝚃}(r_2)`$ for some $`u,v\mathrm{Arr}P\mathrm{\Delta }`$. If there is a match between $`r_1`$ and $`r_2`$ we may write $`u_1\mathrm{𝙻𝚃}(r_1)v_1=u_2\mathrm{𝙻𝚃}(r_2)v_2`$ for some $`u,v\mathrm{Arr}P\mathrm{\Delta }`$. The S-polynomial resulting from a match is then the difference $`u_1r_1v_1u_2r_2v_2`$. The set of S-polynomials of a finite set of polynomials is finite and can be computed.
###### Lemma 3.4
If all S-polynomials resulting from matches of $`R`$ reduce to zero by $`_R`$ then $`_R`$ is confluent on $`\mathrm{Arr}P_K\mathrm{\Delta }`$.
Outline Proof Observing that $`\mathrm{Arr}P_K\mathrm{\Delta }`$ is a subset of the free $`K`$-algebra on $`(\mathrm{Arr}P\mathrm{\Delta })^{}`$ we can deduce that the relation $`_R`$ is confluent on the free $`K`$-algebra. The fact that $`_R`$ preserves source and target enables us to deduce that $`_R`$ cannot reduce an element of $`\mathrm{Arr}P_K\mathrm{\Delta }`$ to anything not defined in $`\mathrm{Arr}P_K\mathrm{\Delta }`$. Thus $`_R`$ is confluent on $`\mathrm{Arr}P_K\mathrm{\Delta }`$. $`\mathrm{}`$
Buchberger’s algorithm calculates the S-polynomials of a system $`R`$ and attempts to reduce them to zero by $`_R`$. If an S-polynomial cannot be reduced it is added to the system. The S-polynomials of the modified system $`R^{}`$ are then computed – the process looping until a system is found whose S-polynomials can all be reduced to zero.
###### Theorem 3.5 (Buchberger’s Algorithm and $`K`$-category Presentations)
If it terminates, then Buchberger’s algorithm applied to $`(R,>)`$, will return a Gröbner basis for $`=_R`$ on $`\mathrm{Arr}P_K\mathrm{\Delta }`$.
Proof All that remains to be verified is that S-polynomials resulting from matches found in $`R`$ can be added to $`R`$ without altering $`\underset{R}{\overset{}{}}`$. We assume all polynomials in $`R`$ to be monic (possible since $`K`$ is a field). Now S-polynomials result from two types of overlap.
For the first case let $`r_1,r_2`$ be polynomials in $`R`$ such that $`u\mathrm{𝙻𝚃}(r_1)=\mathrm{𝙻𝚃}(r_2)v`$ for some $`u,v\mathrm{Arr}P\mathrm{\Delta }`$. Then the S-polynomial is $`s:=\mathrm{𝚛𝚎𝚖}(r_2)vu\mathrm{𝚛𝚎𝚖}(r_1)\mathrm{Arr}P_K\mathrm{\Delta }`$ where $`\mathrm{𝚛𝚎𝚖}(r_i):=r_i\mathrm{𝙻𝚃}(r_i)`$ for $`i=1,2`$. Now $`\mathrm{𝚛𝚎𝚖}(r_2)vu\mathrm{𝚛𝚎𝚖}(r_1)=ur_1r_2v`$ therefore $`s=\mathrm{𝚛𝚎𝚖}(r_2)vu\mathrm{𝚛𝚎𝚖}(r_1)=_R0`$, and hence the congruence generated by $`R^{}:=R\{s\}`$ coincides with $`=_R`$.
For the second case let $`r_1,r_2`$ be polynomials in $`R`$ such that $`u\mathrm{𝙻𝚃}(r_1)v=\mathrm{𝙻𝚃}(r_2)`$ for some $`u,v\mathrm{Arr}P\mathrm{\Delta }`$. Then the S-polynomial is $`s:=\mathrm{𝚛𝚎𝚖}(r_2)u\mathrm{𝚛𝚎𝚖}(r_1)v\mathrm{Arr}P_K\mathrm{\Delta }`$. Now $`\mathrm{𝚛𝚎𝚖}(r_2)u\mathrm{𝚛𝚎𝚖}(r_1)v=ur_1vr_2`$ therefore $`s=\mathrm{𝚛𝚎𝚖}(r_2)u(r_1)v=_R0`$, and hence the congruence generated by $`R^{}:=R\{s\}`$ coincides with $`=_R`$. $`\mathrm{}`$
An example of an application of the results proved above can be found in Section 5.
## 4 Left Kan Extensions
We obtain a further result by expressing the presentation of a noncommutative polynomial algebra as a problem of computing a left Kan extension over framed modules $`\mathrm{𝖪𝖬𝗈𝖽𝗌}`$ (modules over a fixed field).
###### Definition 4.1
A left Kan extension data for $`\mathrm{K}`$-categories $`(M^{},F^{})`$ consists of small categories $`𝖠`$, $`𝖡`$ and functors $`M^{}:𝖠\mathrm{𝖪𝖬𝗈𝖽𝗌}`$ and $`F^{}:𝖠𝖡`$. A left Kan extension presentation for $`\mathrm{K}`$-categories is a quintuple $`𝒫:=kan\mathrm{\Gamma }|\mathrm{\Delta }|RelB|M|F`$ where
1. $`\mathrm{\Gamma }`$ and $`\mathrm{\Delta }`$ are (directed) graphs;
2. $`M:\mathrm{\Gamma }\mathrm{𝖪𝖬𝗈𝖽𝗌}`$ and $`F:\mathrm{\Gamma }P_K\mathrm{\Delta }`$ are graph morphisms to the category of $`K`$-modules and the free $`K`$-category on $`\mathrm{\Delta }`$ respectively;
3. and $`RelB`$ is a set of relations on the free $`K`$-category $`P_K\mathrm{\Delta }`$.
Formally, we say $`𝒫`$ presents the left Kan extension $`(E,\epsilon )`$ of the left Kan extension data $`(M^{},F^{})`$ where $`M^{}:𝖠\mathrm{𝖪𝖬𝗈𝖽𝗌}`$ and $`F^{}:𝖠𝖡`$ if $`\mathrm{\Gamma }`$ is a generating graph for $`𝖠`$ and $`M:\mathrm{\Gamma }\mathrm{𝖪𝖬𝗈𝖽𝗌}`$ is the restriction of $`M^{}:𝖠\mathrm{𝖪𝖬𝗈𝖽𝗌}`$; $`cat_K\mathrm{\Delta }|RelB`$ is a $`K`$-category presentation for $`𝖡`$ and $`F:\mathrm{\Gamma }P_K\mathrm{\Delta }`$ induces $`F^{}:𝖠𝖡`$.
We expect that a left Kan extension $`(E,\epsilon )`$ is given by a set $`EB`$ for each $`B\mathrm{Ob}\mathrm{\Delta }`$ and a function $`Eb:EB_1EB_2`$ for each $`b:B_1B_2𝖡`$ (defining the $`K`$-functor $`E`$) together with a function $`\epsilon _A:XAEFA`$ for each $`A\mathrm{Ob}𝖠`$ (the natural transformation).
For the following theorem it is helpful to note that $`=_{FQ}^r`$ will denote the right congruence generated by $`FQ`$. Square brackets $`[]_{FQ}^r`$ denote the corresponding congruence classes.
###### Theorem 4.2 (Congruences on Algebras are Kan Extensions)
Let $`𝒫:=kan\mathrm{\Gamma }|\mathrm{\Delta }|RelB|M|F`$ be a presentation of a Kan extension for $`K`$-categories where:
1. $`\mathrm{\Gamma }`$ is the graph with one object $`A`$ and a collection of arrows $`Q`$,
2. $`\mathrm{\Delta }`$ is the graph with one object $`B`$ and a set of arrows $`X`$,
3. $`RelB`$ is a set of polynomial relations $`RK[X^{}]`$,
4. $`M:𝖠\mathrm{𝖪𝖬𝗈𝖽𝗌}`$ maps $`A`$ to $`K[1]`$ and the arrows of A to the identity morphism,
5. $`F:𝖠P_K\mathrm{\Delta }`$ maps the arrows of $`𝖠`$ to polynomials of $`K[X^{}]`$
Then the left Kan extension presented by $`𝒫`$ is $`(E,\epsilon )`$ where
1. $`E(B)`$ is isomorphic to $`(K[X^{}]/=_R)/=_{FQ}^r`$,
2. $`E(b)`$ is defined by $`E(b)[p]_R:=[pb]_R`$,
3. $`\epsilon :MEF`$ is given by $`\epsilon _AM(q):=[[q]_R]_{FQ}^r`$.
Outline Proof It is required to verify that $`E`$, as defined above, is a well-defined $`K`$-functor. This is quite routine and comes from the fact that the congruence preserves addition, scalar multiplication and right-multiplication. To verify that $`\epsilon `$ is a natural transformation of $`K`$-functors is straightforward, remembering that $`M(q)`$ is the identity morphism on $`K[1]`$. To check the universal property we suppose there is another such pair $`(E^{},\epsilon ^{})`$ and by drawing the commutative diagram we find that there is a unique natural transformation $`\alpha :EE^{}`$ defined by $`\alpha (b):=E^{}(b)(\epsilon ^{}(1_A))`$ for $`bX^{}`$. $`\mathrm{}`$
It is not claimed that this result is at all deep or difficult, given the results of but it allows the possibility of using Gröbner bases to compute different types of left Kan extensions.
###### Corollary 4.3
Gröbner bases can be used to compute left Kan extensions of the above type.
Outline Proof Let $`𝒫`$ be as above. Define the $`P_K\mathrm{\Delta }`$-set as
$$T:=MA\times \mathrm{Arr}P_K\mathrm{\Delta }$$
and write the terms $`A|p`$ where $`p\mathrm{Arr}P_K\mathrm{\Delta }`$. Define the system of polynomials $`𝒮:=(S_E,S_\epsilon )`$ where
$$S_E:=R\text{ and }S_\epsilon :=\{A|FqA|1:qQ\}$$
The results in describe Gröbner basis procedures for one-sided ideals in finitely presented noncommutative algebras over fields. The polynomials defining the $`K`$-algebra $`𝖡`$ as a quotient of the free $`K`$-algebra $`P_K\mathrm{\Delta }`$ are combined with the polynomials defining a right congruence $`=_{FQ}^r`$ of $`𝖡`$, by using a tagging notation. Standard noncommutative Gröbner basis techniques can then be applied to the mixed set of polynomials, thus calculating $`𝖡/=_{FQ}^r`$ whilst working in a free structure, avoiding the complication of computing in $`𝖡`$.
Suppose $`𝒢`$ is a Gröbner basis for $`𝒮`$. Then the Kan extension is given in the following way:
1. $`E(B):=\mathrm{𝙸𝚁𝚁}`$,
2. $`E(b):A|p\mathrm{𝚒𝚛𝚛}(A|pb)`$, for $`A|p`$ in $`E(B)`$, $`b`$ in $`K[X^{}]`$
3. $`\epsilon _A(1):=A|1`$
where $`\mathrm{𝚒𝚛𝚛}(A|pb)`$ is the irreducible result of repeated reduction of $`A|pb`$ by $`_𝒢`$ and $`\mathrm{𝙸𝚁𝚁}`$ is the set of all irreducible terms of $`T`$. $`\mathrm{}`$
###### Remark 4.4
It is worth remarking that, as with the rewriting methods developed in , the Gröbner basis methods developed in which are referred to above do not require changes in the existing programs. The use of tags enables the combination of polynomials giving the conditions for the action of the Kan extension together with the polynomials giving the conditions for the natural transformation.
## 5 Examples
The first example illustrates the previous section, showing that the standard Gröbner basis computation is the computation of a Kan extension and extending the example to make clear the type of calculation used for right congruences of algebras. The second example demonstrates the use of Gröbner bases to calculate the morphisms of a $`K`$-category given by a presentation. In each case we consider the left Kan extension given by a presentation $`𝒫:=\mathrm{\Gamma }|\mathrm{\Delta }|RelB|M|F`$.
###### Example 5.1
*Let $`\mathrm{\Gamma }`$ be the trivial graph with one object $`A`$. Let $`\mathrm{\Delta }`$ be the graph with one object $`B`$ and arrows $`X:=[e_1,e_2,e_3]`$. Let $`RelB`$ be the set of polynomials*
$$R:=\{e_1e_1e_1,e_2e_2e_2,e_3e_3e_3,e_3e_1e_1e_3,e_2e_1e_2e_1e_2e_1+\frac{2}{9}e_2\frac{2}{9}e_1,e_3e_2e_3e_2e_3e_2+\frac{2}{9}e_3\frac{2}{9}e_2\}.$$
*$`F:\mathrm{\Gamma }P_K\mathrm{\Delta }`$ be inclusion and define $`M(A):=K[1]`$. The system $`𝒮`$ consists only of untagged polynomials $`R`$ because there are no non-trivial arrows in $`𝖠`$. We use the length-lexicographic ordering with $`e_3>e_2>e_1`$ to obtain Gröbner basis for the congruence generated by $`𝒮`$ in $`K[X^{}]`$ by adding*
$$e_3e_2e_1e_3e_2e_3e_2e_1+\frac{2}{9}e_2e_1\frac{2}{9}e_1e_3$$
*to $`R`$. The irreducible terms $`\mathrm{𝙸𝚁𝚁}`$ in this case are sums of $`K`$-multiples of the following terms*
$`\{A|1`$, $`A|e_1,A|e_2,A|e_3`$, $`A|e_1e_2,A|e_1e_3,A|e_2e_1,A|e_2e_3,A|e_3e_2`$, $`A|e_1e_2e_1,A|e_1e_2e_3,A|e_1e_3e_2,A|e_2e_1e_3,A|e_2e_3e_2,A|e_3e_2e_1`$, $`A|e_1e_2e_1e_3,A|e_1e_2e_3e_2,A|e_1e_3e_2e_1,A|e_2e_1e_3e_2,A|e_2e_3e_2e_1\}.`$
*In this example the tag “$`A|`$” is redundant: the $`K`$-module $`EB`$ is a $`K`$-algebra, in fact it is the Hecke algebra $`H_4`$. Suppose now that $`\mathrm{\Gamma }`$ has one arrow $`q`$ whose image under $`F`$ is $`e_2e_1`$. The system of polynomials $`𝒮`$ for the Kan extension now has an $`\epsilon `$-polynomial namely $`A|e_2e_1A|1`$. Applying Buchberger’s Algorithm with the length-lexicographic ordering $`e_3>e_2>e_1>A|`$ results in a Gröbner basis of mixed polynomials:*
$`\{e_1e_1e_1,e_2e_2e_2,e_3e_3e_3,e_3e_1e_1e_3,e_2e_1e_2e_1e_2e_1+\frac{2}{9}e_2\frac{2}{9}e_1,e_3e_2e_3e_2e_3e_2+\frac{2}{9}e_3\frac{2}{9}e_2,`$
$`e_3e_2e_1e_3e_2e_3e_2e_1+\frac{2}{9}e_2e_1\frac{2}{9}e_1e_3,A|e_2e_1A|1,A|e_1A|1\}`$.
*The right congruence classes of $`e_2e_1`$ on $`H_4`$ are represented by sums of $`K`$-multiples of the following irreducible terms i.e. $`\mathrm{𝙸𝚁𝚁}`$ consists of:*
$`\{A|1`$, $`A|e_2,A|e_3`$, $`A|e_2e_3,A|e_3e_2`$, $`A|e_2e_3e_2,A|e_3e_2e_1`$, $`A|e_2e_3e_2e_1\}.`$
*Here the tag “$`A|`$” is necessary in the computation of the Gröbner basis. The final results may be written as right congruence classes $`[e_3e_2]^r`$, say, instead of tagged terms $`A|e_3e_2`$ but the representation as tagged terms allows us to determine whether $`e_1e_2e_3`$ and $`e_2e_3`$ occur in the same class: reducing $`A|e_1e_2e_3`$ and $`A|e_2e_3`$ has the same result, so they are congruent.*
###### Example 5.2
*Let $`𝖡`$ be the $``$-category generated by the graph $`\mathrm{\Delta }`$:*
*The arrows of the free category $`P_K\mathrm{\Delta }`$ are sums of $``$-multiples of terms occuring in the same column of the following table (the hom-sets consisting solely of identities are omitted):*
| $`B_1B_2`$ | $`B_1B_3`$ | $`B_1B_4`$ | $`B_1B_5`$ | $`B_2B_2`$ | $`B_2B_3`$ | $`B_4B_3`$ | $`B_5B_4`$ | $`B_5B_3`$ |
| --- | --- | --- | --- | --- | --- | --- | --- | --- |
| $`a`$ | $`ac`$ | $`h`$ | $`e`$ | $`1_{B_2}`$ | $`c`$ | $`g`$ | $`f`$ | $`j`$ |
| $`ab`$ | $`abc`$ | $`ef`$ | | $`b`$ | $`bc`$ | | | $`fg`$ |
| $`ab^2`$ | $`ab^2c`$ | | | $`b^2`$ | $`b^2c`$ | | | |
| | | | | | | | | |
| $`ab^n`$ | $`ab^nc`$ | | | $`b^n`$ | $`b^nc`$ | | | |
*Let $`R`$ be the relations defining $`𝖡`$*
$$R:=\{ab^3ab^2ab+a,b^3cb^2cbc+c,abc+defg,ac+dhg,fgj\}$$
*Applying the length-lexicographic ordering with $`a<b<c<d<e<f<g<h<j`$ it can be checked that $`R`$ is a Gröbner basis. It can therefore be immediately deduced that the arrows of $`𝖡`$ are uniquely represented by $``$-multiples of terms occurring in the same column of the following table:*
| $`B_1B_2`$ | $`B_1B_3`$ | $`B_1B_4`$ | $`B_1B_5`$ | $`B_2B_2`$ | $`B_2B_3`$ | $`B_4B_3`$ | $`B_5B_4`$ | $`B_5B_3`$ |
| --- | --- | --- | --- | --- | --- | --- | --- | --- |
| $`a`$ | $`ac`$ | $`h`$ | $`e`$ | $`1_{B_2}`$ | $`c`$ | $`g`$ | $`f`$ | $`j`$ |
| $`ab`$ | $`abc`$ | $`ef`$ | | $`b`$ | $`bc`$ | | | |
| $`ab^2`$ | $`ab^2c`$ | | | $`b^2`$ | $`b^2c`$ | | | |
| | | | | | | | | |
| | | | | $`b^n`$ | | | | |
## 6 Further Questions
### 6.1 Induced Modules
It would be useful to phrase the results of Section 4 in terms of induced modules, relating it to the commutative case in .
### 6.2 Extensions of Gröbner basis techniques
To apply rewriting to Kan extensions we had to generalise it. We have not yet discovered how precisely to generalise Gröbner bases to apply to *any* Kan extension of $`K`$-categories over $`\mathrm{𝖪𝖬𝗈𝖽}`$.
### 6.3 Rings with Many Objects
Mitchell’s classic work, generalises noncommutative homological ring theory to (pre)additive category theory . His work motivates the investigation of Gröbner basis techniques for $`K`$-categorical Kan extensions by the potential for Gröbner bases to provide more powerful methods of computation (of homology or cohomology) in this setting.
### 6.4 Term rewriting and Monads
Term rewriting systems, widely used throughout computer science, are similar to algebraic theories (algebraic theories declare term constructors, term rewriting systems declare term constructors and rewrite constructors). Algebraic theories can be modelled by finitary monads over $`\mathrm{𝖲𝖾𝗍𝗌}`$. Term rewriting systems can be modelled by finitary monads over the category of preorders $`\mathrm{𝖯𝗋𝖾}`$. This has been useful in providing categorical proofs of rewriting theories. The particularly interesting point is that term rewriting systems can be modelled as monads over a more complex base category. So $`𝖢`$-algebraic theories can be modelled by finitary monoids on $`𝖢`$. There is a relation between monads, adjoint functors and Kan extensions. We need to investigate the relation between string rewriting for Kan extensions and the monads modelling algebraic theories and term rewriting systems.
### 6.5 Petri nets
Gröbner basis procedures can be usefully applied in Petri net analysis. To every Petri net there is an associated category – a Petri category . How does the structure for Petri categories relate to Kan extensions? Are the Gröbner basis techniques usefully extended by relating these two areas or are they in fact the means by which the areas can be related?
### 6.6 Automatic Structure
For groups, monoids and coset systems there is a well-known concept of an automatic structure. These systems are special cases of Kan extensions so it is natural to ask what would be the definition of an automatic structure for a left Kan extension in general.
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# REFERENCES
A New Hypothesis for Layers of High Reflectivity Seen in MST Radar Observations
G. C. Asnani and M. K. Rama Varma Raja$`{}_{}{}^{}`$
Indian Institute of Tropical Meteorology, Dr. Homi Bhabha Road , Pune 411 008(India). E-mail: asnani@giaspn01.vsnl.net.in or Prof.asnani@vsnl.com
$`{}_{}{}^{}`$ Department of Physics, University of Pune, Pune 411 007(India). E-mail: rama@physics.unipune.ernet.in
We have worked with MST Radar located at Gadanki, Tirupati, India(13.47 degree N, 79.18 degree E). Altitude of the location is approximately 100 meters. This is the first and only MST Radar operating in India. It is in Tropical Monsoon region. Monsoon moves northward across Tirupati around 1st June and withdraws southward across Tirupati in the first week of December. Our interest has been to examine the characteristics of vertical velocity and other wind characteristics associated with monsoon. We examined the characteristics of range-compensated Signal-to-Noise ratio ($`r^2`$ SNR). We call this as Reflectivity. We examined the MST Radar data set for 14 months (September 1995 to November 1996 except May 1996). The height range considered for the analysis was between 3.6 km and 21 km above ground. Our findings have been on the following lines:
1) Except on Thunderstorm occasions, high reflectivity-regions are in the form of layers 1-2 km thick. One such high-reflectivity layer is always present near tropopause level 17 km. Below the tropopause, there is a bunch of such high-reflectivity layers, generally 3-4 in number, between 4 and 11 km. The atmospheric layer between 11 and 15 km is generally free from high reflectivity layer.
2) On a thunderstorm occasion, there is a deep high-reflectivity layer extending from 4-11 km. After the passage of the thunderstorm, this deep layer of high reflectivity breaks up into layers.
3) Layers of high reflectivity occur throughout the year even outside monsoon season, when ITCZ is far away from the region and they cannot be attributed as originating from deep convective clouds. When visual observation and MST observations are taken simultaneously at Gadanki, visual clouds are estimated to be at the same levels as the high reflectivity layers seen through MST observations. Indian satellite pictures over Gadanki also suggest similar heights of clouds as given by MST Radar for high reflectivity layers.
4) Satellite pictures in infra-red range show much more extensive areas of cloudiness than the pictures in the satellite visible range. In other words, there are extensive sub-visual clouds in the atmosphere.
5) Every month early morning and late evening, in twilight hours we see beautiful cirrus clouds in the form of cloud streets, streaks or sheets on several occasions which may or may not be clearly visible at other hours.
6) On looking at the sky frequently, one gets the impression that the sky is not clear even though we do not see clouds or haze layers. In these haze layers, appear clouds in the form of cloud streets, streaks or cloud sheets at some times. The appearance of these clouds gives the clear impression that there are waves in the atmosphere which give visible clouds in the region of upward motion associated with these waves. Again in each bigger wave cloud, there are smaller and smaller wave clouds and clearances; there are waves within waves. When the clouds disappear, they leave a sort of haze. Hence, the clouds form out of haze and leave some haze after dissipating. Even when there are no visible clouds during the day or night, the structure of the atmosphere is patchy in appearance. One gets a clear impression that there are waves and waves, bigger and smaller waves in the atmosphere which are creating patches of haze and sometimes visible clouds in the atmosphere. These are due to gravity-type waves with a very wide spectrum of horizontal and vertical wavelengths. As we know from theoretical and observational evidence, horizontal wavelengths are 1-2 orders of magnitude larger than vertical wavelengths. If and when they occur in the atmosphere, they have a tendency to take horizontally spread layered structure with vertical depth 1-2 orders of magnitude smaller than horizontal extent.
7) In the atmosphere, we visualize three classes of waves: a.Inertial waves or Rossby-type waves: Their horizontal extent is of the order of a few thousand kilometers and vertical wavelengths of the order of 10 km. In the mechanism of their formation, we have to consider the rotation of the earth and the resulting coriolis force. Their period is of the order of a few days. b. Gravity waves: These arise mainly from local horizontal pressure gradients arising out of gravitational weight of the overlying air column, air accelerating from higher pressure towards lower pressure. Horizontal accelerations and displacements are accompanied by appropriate vertical accelerations and displacements to conform to the requirements of law of conservation of mass i.e., equation of continuity.
Vertical displacements and accelerations of air parcels can also arise from buoyancy forces. Heavier parcel tends to sink down while lighter air parcels tend to rise up in an environment of horizontal density gradients. These are called Brunt-Vaisala oscillations. These horizontal gradients of density arise out of differences in temperature, humidity and hydro-meteor loading.
This hydrometeor-loading needs a little elaboration:
(i) Every parcel of cloud air contains at least one hydrometeor in liquid water or solid form. Inside each hydrometeor is an aerosol which acts as a nucleus which has induced condensation and/or freezing.
(ii) Invariably, hydrometeor has higher density than the surrounding air. As such, it tends to fall down due to gravitational force.
(iii) As it descends down, it exchanges sensible heat and moisture with its environment. Hence, its volume, mass, and density undergo a change during its descent.
(iv) As the hydrometeor descends down with gravitational acceleration through the air parcel, there is frictional resistance/viscous resistance to its vertical motion. soon, the hydrometeor loses its acceleration and descends down with what is called ”Terminal velocity”. Where has its weight gone? Its weight is taken up by the air parcel which is offering resistance to its vertical movement and acceleration. In other words, the air parcel becomes heavier to the extent it has taken over the gravitational acceleration of the falling hydrometeor. If the hydrometeor falls with its terminal velocity, it means that its entire weight is taken over by the air parcel. As such the air parcel has become heavier by the total weight of the hydrometeor; the density of the air parcel has effectively increased.
If only half of the gravitational acceleration of the descending hydrometeor has been taken over by the air parcel, then the weight of the air parcel has increased only by half the weight of the hydrometeor.
The hydrometeor gradually loses its gravitational acceleration and gives its weight to the surrounding air parcel gradually and not instantaneously. As such, the air parcel takes the load of the hydrometeor gradually during the downward trajectory of the descending hydrometeor.
In many calculations, for the sake of simplicity, it is assumed that as soon as condensation or freezing takes place in the atmosphere, the load of the hydrometeors is immediately taken over by the surrounding air. However, we have to recognize that the hydrometeor-loading occurs gradually through a finite interval of time.
We should also remember that during the fall, the hydrometeor is simultaneously undergoing a change in its volume, mass and density. Hence, a realistic, quantitative estimate of hydrometeor loading effect on the air parcel needs careful calculation. However, nature takes care of the process and creates varying density effects on the cloud air parcel as the hydrometeor descends.
(v) In addition to buoyancy fall of the hydrometeor, there are upward and downward motions of air parcels inside the cloud. These upward and downward motions of air parcels inside the cloud create further complications in the calculation of hydrometeor-loading effect on air parcel.
(vi) In addition to pure buoyancy forces operating in a class of waves called gravity waves, there also occur what are known as Kelvin-Helmholtz waves due to presence of vertical shear of horizontal winds which is almost always present.
(vii) This class of gravity waves of Kelvin-Helmholtz type have very small wavelengths of the order of centimeters and meters and correspondingly small periods of the order of a few minutes. Earth’s rotation or Coriolis force does not perceptibly come into the calculations for these waves.
c. Inertio-Gravity waves: Between the two extremes of large inertial waves and small gravity waves, there is an intermediate class of waves which may be termed as inertio-gravity waves in which Coriolis force plays some role along with gravity force. There is literature on the subject of inertio-gravity waves ( , , , etc.), but more work needs to be done on this class of waves. Orographic influence also come into play.
8) When we fly in an aircraft, large-scale weather and clouds are influenced by Rossby-type waves. When we look at the sky through cockpit or through the window near the window-seat, during day time, we immediately get the view of air clouds at different levels and also gravity waves of various dimensions near the flight level. We see fairly large waves with estimated wavelengths of the order of tens of kilometers, along with embedded smaller and smaller waves, thick clouds, thin clouds, thinner clouds, space filled with haze and space clear of visible clouds.
In our view, MST Radar reflectivity pattern gives us spot view of these numerous beautiful waves.
9) We have analyzed the field of MST Radar reflectivity as seen at Gadanki along with MST Radar measured wind fields. We have interpreted the reflectivity fields within a conceptual model given below:
(i) Inertio-gravity waves in the atmosphere generate layers of upward/downward motion, high/low humidity and high/low temperature lapse rates. The layers of upward and downward motion are regularly seen in the vertical wind field given by MST Radar. The vertical wavelength of these inertio-gravity waves has wide spectrum depending on orography and diabatic heating; vertical wavelength of about 5 km is a more frequently observed wavelength. The corresponding horizontal wavelength is of the order of 200 km, the more dominant wavelength visible in satellite picture is 500 km in the direction of wind and 1000 km across the wind.
The layers of high relative humidity created by inertio-gravity waves are favorable for the formation of layered clouds, which we call ”Mother Cloud Layers”; these clouds may be visible or sub-visible.
(ii) Hydrometeors inside a ”Mother Cloud Layer” tend to fall down attaining their respective terminal velocities. During their stay inside the clouds, the hydrometeors exchange heat, moisture, mass and momentum with the environmental air on small micro-scales. These exchanges between the hydrometeors and the ”Mother Clouds’s air” create strong gradients of temperature, humidity, density and momentum. Density variations are also created through hydrometeor-loading.
Electrical charges are also generated during the processes of condensation, evaporation, freezing, melting and sublimation.
These micro-physical gradients in density along with prevailing wind field in the vertical generate internal gravity waves in the form of Brunt-Vaisala oscillations, Kelvin-Helmholtz waves and other waves of different horizontal and vertical wavelengths. These wavelengths range from a few millimeters to tenths of meters in the vertical and from a few meters to about thousand meters in the horizontal. Known laws of physics and dynamics suggest that the wavelengths may be still smaller, equivalent to the distances between adjacent parcels of air exchanging heat moisture, mass and momentum, with the hydrometeor embedded in the parcel.
(iii) The strong gradients of temperature, humidity and density created by these micro-physical and micro-dynamical processes in the air surrounding the hydrometeor or an ensemble of hydrometeors cause strong variations in the refractive index of air parcel, in respect of electro-magnetic lidar and radar beams impinging on the air parcels. In turn, this causes high reflectivity/scatter of the impinging lidar/radar beam. In respect of VHF MST radar Bragg-type reflection/scatter is a dominant type of reflection/ scatter.
We examined the size of air parcels giving the highest values of reflectivity, at Gadanki. We came to the conclusion that their horizontal extent is of the order of 1 km while their vertical extent is of the order of 100 m.
Indian MST radar beam oriented in vertical has a half-wavelength of about 3m. As such, Indian MST radar is capable of detecting reflectivity patterns of vertical wavelengths of the order of about 6m. These small-scale variations in reflectivity appeared in the form of very delicate embroidery inside the large-scale reflectivity of the ”Mother Cloud Layer”.
(iv) These variations in the reflectivity pattern may look like turbulent fluctuation in the atmosphere. If we do not associate these fluctuations with clouds and hydrometeors inside the clouds, the clouds might appear as clear-air turbulence as has been prevalent in the current explanation appearing in literature connected with MST radars. This prevalent explanation has faced paradoxes, the main paradox being of thin horizontal sheets of turbulence.
If we free our thinking from the concept of clear-air turbulence and turn our thinking in the direction of visible or sub-visible clouds containing hydrometeors, the apparent paradoxes of clear-air turbulence causing high- MST radar reflectivity get immediately resolved.
The existence of sub-visible clouds occupying much larger area than the visible clouds has now been established beyond question, through latest observations by satellites in the infra-red range, by aircraft flying through cirrus cloud air and by lidars operating in high altitude aircraft sending their beams through visible and sub-visible cloud layers.
(v) Using over 2,50,000 observational data points for MST radar reflectivity and vertical wind shear (deduced from corresponding MST radar wind observations) spread over 14 months (from 1995 September to 1996 November), we plotted scatter diagrams of MST radar reflectivity versus vertical wind shear. We are pleasantly surprised to find that reflectivity decreases almost exponentially as vertical wind shear increases. If mechanical turbulence was the main cause of high reflectivity, we should see reflectivity increasing with vertical wind shear, and not decreasing almost exponentially. Scatter diagrams for each of the 14-months are presented in . Also, Scatter diagrams for four representative months (January, April, July and October), for Indian monsoon region, are presented in ,. This shows that mechanical turbulence is not the principal cause of high-MST radar reflectivity. had hypothesized that turbulence may not be the primary cause of high-reflectivity seen in MST radars.
(vi) Knowing the importance of cirrus clouds, the world scientific community launched the programme known as FIRE (First ISCCP Regional Experiment). This FIRE programme concentrated on the study of layer clouds (Cirrus Clouds in the upper troposphere and low level stratus clouds in the lower troposphere). FIRE I programme was executed during the period 1985-1990 while FIRE II programme was executed during the period 1990-1995. FIRE III programme is proposed to be executed in the beginning of this new millennium, with particular emphasis on the tropics. Results of FIRE I have been summarized in a special issue of Monthly Weather Review (November, 1990); results of FIRE II have been summarized in a special issue of Journal of Atmospheric Sciences (December, 1995). This topic is an important component of CLIVAR programme in the tropics.
The results of FIRE I and FIRE II have broadly confirmed that there is fine micro structure inside the layered clouds which can be interpreted easily, atleast qualitatively, in terms of micro-physical and micro-dynamical processes presented above.
Also, the latest Numerical Modeling Work on Cirrus Clouds (for example , ,) show that there are fine-scale structures of various in-cloud parameters including temperature, humidity and ice-concentrations.
In fact in our view these fine scale structures created inside the visible or sub-visible ”Mother Cloud Layers” can cause steep refractive index gradients sufficient to cause high-MST radar reflectivity.
10) We expect that the conceptual model of MST radar reflectivity presented above will give a new orientation to the thinking and interpretation of MST radar reflectivity. It will provide a satisfactory, physically and dynamically acceptable, interpretation of the reflectivity patterns seen in the MST radar observation.
11) A few corollaries follow from this conceptual model:
(i) Since, high-MST radar reflectivity layer, 1-2 km thick, is always observed near the tropopause, it may or may not be directly connected to inertio-gravity waves. The mechanism for the formation and sustenance of high-reflectivity layer is realized as follows:
(a) The temperature near the tropopause level are very cold ( 200 K); temperature lapse rate is stable, 2-4 degree C/km; as such vertical mixing of air is inhibited and is weak. Relative humidity below the tropopause is high . Water substance and aerosols injected into the upper troposphere by deep convection near ITCZ remains below the tropopause while the relative humidity is very low above the tropopause. Water substance and aerosols form visible or sub-visible cirrus cloud layer or haze layer below the tropopause. Through the micro-physical and micro-dynamical processes mentioned earlier, the air layer develops strong vertical gradients or discontinuities in the refractive index with respect to MST radar beam impinging on the air there. This gives high-reflectivity echoes near the tropopause level on all days of the year.
When the cloud gets dissolved, it leaves large number of aerosols suspended there. By themselves, aerosols are not likely to be detected by MST radar. But the layer of high content of aerosols can be detected and has been detecting near the tropopause (,, etc.).
In the tropics, ITCZ injects aerosols and moisture which remain trapped in the layer immediately below the tropopause. In extra-tropics, this injection is done by extra-tropical cyclone waves and polar front. On the same reasoning, as given for tropics, a layer with high content of aerosols and thin cirrus ice crystals will also get formed below the tropopause in extra-tropics. Thus a layer having high content of aerosols and thin cirrus cloud is expected to envelope the whole earth’s atmosphere near tropopause. This has been substantiated by observations of (Nee et al., 1995, 1998; FIRE I and FIRE II observations and their results published in the Special issues of Monthly Weather in November 1990 and Journal of Atmospheric Sciences in December 1995 respectively). Asnani et al. had hypothesized the existence of an aerosol layer near tropical tropopause, based on the same mechanism mentioned above.
(ii) As we have stated above, observations confirm existence of alternate layers of upward and downward motion with vertical depth of 2-3 km. We should expect accumulation and depletion of both water substance and aerosol substance near the levels of vertical convergence. Further vertical upward motion tends to create higher relative humidity and higher temperature lapse rate; vertical downward motion tends to do the opposite, creating stable lapse rate and drier air. Hence we should expect to find layers of high and low relative humidity in the vertical. Sensitive instrumentation is required to detect this type of structure in the atmosphere; this structure is likely to be missed by the ordinary radiosonde instruments. Such layers have been detected by special effort .
(iii) As stated earlier, these alternate layers of upward and downward motion are associated with inertio-gravity waves, which are always present in the atmosphere. These vertical wave motion will also tend to split large convective clouds, particularly in their decaying stage, into layered clouds. While convective instability in the atmosphere tends to generate deep convective clouds in the tropical atmosphere, these inertio-gravity waves inhibit the formation and sustenance of these deep convective clouds in the tropical atmosphere.
(iv) Cirrus clouds have a capacity to retain their existence for a long time, even away from the source of their formation. Hence, visible or sub-visible cirrus clouds are likely to be seen at many places, with or without deep convective clouds. The upward vertical motion associated with inertio-gravity waves tends to generate cirrus clouds, visible or sub-visible at many places in the atmosphere. If nothing else, these clouds influence the radiative heat budget of the earth’s atmosphere.
Acknowledgment: The authors thank Prof. D. Narayana Rao and his group at Department of Physics, Sri Venkateswara University, Tirupati (India) for considerable help in Indian MST Radar data collection and analysis. Mr. M. K. Rama Varma Raja thanks Dr. (Mrs.) P. S. Salvekar and the Director, Indian Institute of Tropical Meteorology for providing necessary support and facilities during the course of research work.
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# Haldane-gap excitations in the low-𝐻_𝑐 1-dimensional quantum antiferromagnet Ni(C5H14N2)2N3(PF6) (NDMAP).
## I Introduction
The unique properties of the one-dimensional (1D) integer-spin Heisenberg antiferromagnet (HAF) have captivated the minds of condensed matter physicists for the last 20 years. In total defiance of the quasi-classical picture of magnetism, the ground state in this system is a spin-singlet- a “quantum spin liquid” with only short-range (exponentially decaying) spin correlations. The excitation spectrum is also rather unique, and, even in the absense of any magnetic anisotropy, features a so-called Haldane energy gap. The wealth of theoretical and experimental results accumulated to date provide a fairly complete understanding of the physics involved, and few mysteries remain, as far as the behavior of the idealized model is concerned. Particularly revealing was the neutron scattering work done on real quasi-1D $`S=1`$ HAF compounds such as CsNiCl<sub>3</sub>, NENP (Ni(C<sub>2</sub>H<sub>8</sub>N<sub>2</sub>)<sub>2</sub>NO<sub>2</sub>(ClO<sub>4</sub>)), and Y<sub>2</sub>BaNiO<sub>5</sub> (Ref. ).
One of the few remaining unresolved issues is the behavior of Haldane-gap antiferromagnets in high magnetic fields. An external magnetic field splits the triplet of Haldane excitations, driving one of the modes to zero energy at some critical field $`H_c`$. At this soft-mode quantum phase transition, static long-range antiferromagnetic correlations appear. In essence, the external field suppresses zero-point quantum spin fluctuations that are responsible for the destruction of long-range order in the spin-liquid state, and restores long-range spin coherency. Despite the success of high-field bulk studies of several material, it has been highly frustrating, that for most existing systems the actual values of the critical fields are unobtainable in neutron scattering experiments. As a result, the properties of the high-field phase are still largely unknown. It is even unclear, whether or not the high-field phase is commensurate, and virtually nothing is known about the excitation spectrum.
Traditionally, NENP was the workhorse of high-field studies. For this compound large single crystal samples can be fabricated, and the critical field $`H_c=9`$ T (when applied parallel to the chain-axis) is, in principle, accessible in a neutron-scattering experiment. In practice, it is however easier to perform measurements with a magnetic field applied perpendicular to the spin chains. The corresponding critical field for NENP is, unfortunately, much larger: $`H_c1113`$ T. In any case, measurements deep within the high-field phase are not possible. Moreover, due to certain structural features, the phase transition in NENP is smeared out. The g-tensors of the $`S=1`$ Ni<sup>2+</sup> ions are staggered in this compound and lead to an effective staggered field when an external uniform field is applied. As a result, a static staggered magnetization appears at arbitrary weak applied fields, and, strictly speaking, no additional spontaneous symmetry breaking occurs at $`H_c`$.
The recent discovery of NDMAP (Ni(C<sub>5</sub>H<sub>14</sub>N<sub>2</sub>)<sub>2</sub>N<sub>3</sub>(PF<sub>6</sub>)), a relatively easily crystallized Haldane-gap compound with a critical field of only around 4 T, and no staggering of g-tensors within spin chains, promises to make the high-field phase readily accessible in neutron scattering studies. The crystal structure of this material is similar to that of NENP, and is visualized in Fig. 1. The AF spin chains run along the $`c`$ axis of the tetragonal structure (space group $`P_{nmn}`$, $`a=18.046`$ Å, $`b=8.705`$ Å, and $`c=6.139`$ Å), and are composed of octahedrally-coordinated $`S=1`$ Ni<sup>2+</sup> ions bridged by triplets of nitrogen atoms. These long nitrogen bridges account for a realtively small in-chain AF exchange constant $`J=2.6`$ meV, as estimated from bulk $`\chi (T)`$ measurements. ESR and specific heat studies revealed a transition to the high-field phase at $`H_c^{}=3.4`$ T, and $`H_c^{}=5.8`$ T, extrapolated to $`T0`$, for a magnetic field applied parallel and perpendicular to the chain-axis, respectively. This anisotropy of critical field is attributed to single-ion easy-plane magnetic anisotropy of type $`DS_z^2`$. The anisotropy constant was obtained from bulk susceptibility data: $`D/J0.3`$. Using the well-known numerical result $`\mathrm{\Delta }_z0.41J+2pD`$, $`\mathrm{\Delta }_x0.41JpD`$, $`p2/3`$, one can thus estimate the Haldane gap energies: $`\mathrm{\Delta }_{}2.1`$ meV and $`\mathrm{\Delta }_{}0.54`$ meV.
While NDMAP appears to be an ideal model system for neutron scattering experiments in the high-field phase, additional characterization, particularly a measurement of inter-chain interactions and in-plane anisotropy, is required before such a study can be carried out. Finding the 3D AF zone-center is especially important, since it is at this wave vector that static long-range correlations are expected to appear in the high-field phase. Obviously, inelastic neutron scattering is the most direct method to obtain this information. In the present paper we report the results of a zero-field inelastic cold-neutron scattering study of deuterated NDMAP single-crystal samples, aimed at extracting this information.
## II Experimental procedures and results
Fully deuterated NDMAP single crystal samples were grown from solution as described in Ref. . It was observed that the crystals tend to shatter when cooled to low temperature, and even more so when warmed back up. When wrapped in aluminum foil the crystals do not fall apart, but the width of the mosaic distribution increases dramatically with thermal cycling. Most inelastic neutron scattering data were collected on a 140 mg single crystal sample that was taken trough the cooling cycle only twice. The mosaic of the as-grown sample was roughly $`0.3^{}`$ and increased to $`1.5^{}`$ and $`3^{}`$ after the first and second cooling, respectively. Inelastic neutron scattering measurements were performed on the SPINS 3-axis spectrometer installed on the cold neutron source at the National Institute of Standards and Technology Center for Neutron Research. Pyrolitic graphite crystals set for their (002) reflection were used as monochromator and analyzer. Beam divergences were approximately $`40^{}80^{}80^{}240^{}`$ through the instrument, with a cooled Be filter between the monochromator and sample. The measurements were done with a fixed final neutron energy $`E_f=2.8`$ meV. The crystal was mounted with either the $`a`$ or $`b`$ crystallographic axis vertical, making $`(0,k,l)`$ and $`(h,0,l)`$ reciprocal-space planes accessible to measurements. The sample was cooled to $`T=1.4`$ K in an ”ILL-Orange” cryostat with a 70 mm diameter sample well.
The primary goal of the experiment was to determine the $`a`$\- and $`b`$-axis dispersion relations and the polarizations of the two lower-energy Haldane-gap excitations, relevant for the soft-mode transition at $`H_c`$. Most of the data were collected in constant-$`Q`$ scans at the 1D AF zone-center $`l=0.5`$ in the energy transfer range 0–1 meV. In these dispersion measurements $`Q`$-resolution is particularly important, so a flat analyzer was used. Typical scans for different momentum transfers perpendicular to the $`c`$-axis are shown in Fig. 2. At all wave vectors a well-defined peak corresponding to the lower-energy Haldane excitation doublet is clearly seen around 0.5 meV energy transfer.
To observe the higher-energy Haldane-gap excitation and obtain an accurate measurement of the anisotropy constant we performed a constant-$`Q`$ scan in the range 0–2.4 meV (Fig. 3). To maximize intensity we used a horizontally-focusing analyzer pointing $`𝐜^{}`$ towards the analyzer to maintain wave vector resolution along the chain. The scattering angle was varied so the projection of wave vector transfer along the chain was $`0.5𝐜^{}`$ throughout the scan. In addition to the peak seen in the low-energy scans, a weaker feature is observed at $`\mathrm{}\omega 2`$ meV that can be attributed to the $`c`$-axis-polarized Haldane gap mode.
## III Data analysis and discussion
The measured constant-$`Q`$ scans were analyzed using a parameterized model cross section, numerically convoluted with the Cooper-Nathans 3-axis spectrometer resolution function. Near the 1D AF zone-center the single mode approximation (SMA) for the dynamic structure factor of isolated Haldane spin chain is known to work extremely well. For each channel of spin polarization the dynamic structure factor $`S^{(\alpha \alpha )}(𝒒,\omega )`$ can be written as:
$`S^{(\alpha \alpha )}(𝒒,\omega )|f(q)|^2(1{\displaystyle \frac{𝒒𝒆_\alpha }{q^2}})\times `$ (1)
$`\times {\displaystyle \frac{1\mathrm{cos}(𝒒𝒄)}{2}}{\displaystyle \frac{Zv}{\omega _{\alpha ,𝒒}}}\delta (\mathrm{}\omega \mathrm{}\omega _{\alpha ,𝒒}),`$ (2)
$`(\mathrm{}\omega _{\alpha ,𝒒})^2=\mathrm{\Delta }_\alpha ^2+v^2\mathrm{sin}^2(𝒒𝒄).`$ (3)
Here $`v`$ is the spin wave velocity, given by $`v2.49J`$. The dimension-less constant $`Z`$ defines the static staggered susceptibility of a Haldane spin chain: $`\chi _\pi =Zv/\mathrm{\Delta }`$.<sup>*</sup><sup>*</sup>*Here we have adopted the notation used in Ref. In the above expression we have included the magnetic form factor $`f(q)`$ for Ni<sup>2+</sup> and the polarization factor $`\left(1\frac{𝒒𝒆_\alpha }{q^2}\right)`$. The structure factor for weakly-coupled chains can be calculated in the Random Phase Approximation (RPA). The expression for $`S^{(\alpha \alpha )}(𝒒,\omega )`$ does not change explicitly, but the excitations acquire dispersion perpendicular to the chain axis:
$$(\mathrm{}\omega _{\alpha ,𝒒})^2=\mathrm{\Delta }_\alpha ^2+v^2\mathrm{sin}^2(𝒒𝒄)+ZvJ^{}(𝒒).$$
(4)
In this formula $`J^{}(𝒒)`$ is the Fourier transform of inter-chain exchange interactions, that we assume to be isotropic (of Heisenberg type). According to numerical calculations $`Z1.26`$. The form of $`J^{}`$ can be guessed by looking at the crystal structure. The smallest inter-chain Ni-Ni distance (8.705Å) is along the $`b`$ crystallographic axis. We shall denote the corresponding exchange constant as $`J_y`$. The next-smallest inter-chain Ni-Ni distance (10.478Å) is along the (0.5,0.5,0.5) direction. This interaction, however, is frustrated by in-chain AF interactions (any site in one chain couples to two consecutive sites in another chain), and is thus irrelevant, within the RPA, at momentum transfers $`𝒒𝒄\pi `$. Finally, the third-smallest inter-chain distance (18.046Å) is along $`a`$. The corresponding coupling constant $`J_x`$ is expected to be very small, due to the large “bond” length. The Fourier transform of inter-chain interactions can thus be written as:
$$J^{}(𝒒)=2J_x\mathrm{cos}(𝒒𝒂)+2J_y\mathrm{cos}(𝒒𝒃).$$
(5)
The chain-axis exchange constant $`J`$ is not very important for our measurements: for scans collected at the 1D AF zone-center it only slightly influences the peak shapes, due to a non-zero $`Q`$-resolution along the chain-axis. We can therefore safely use the value previously obtained from susceptibility measurements. The relevant parameters of the model are thus the three gap energies $`\mathrm{\Delta }_x`$, $`\mathrm{\Delta }_y`$ and $`\mathrm{\Delta }_z`$, two inter-chain exchange constants $`J_x`$ and $`J_y`$ and an overall intensity prefactor. Three additional parameters were used to describe the background: an energy-independent component, and the intensity and width for a Lorenzian profile centered at $`\mathrm{}\omega =0`$ to account for incompletely resolved elastic scattering.
The difference between $`\mathrm{\Delta }_x`$ and $`\mathrm{\Delta }_y`$, usually a result of in-plane magnetic anisotropy of type $`E(S_x^2S_y^2)`$, is too small to observe two and separate corresponding peaks in any single constant-$`Q`$ scan. Distinct polarization factors for the two branches however allow us to extract both parameters in a global fit to the data measured at different momentum transfers perpendicular to the chain-axis. As a first step, we simultaneously analyzed the data collected in $`(0,k,l)`$ reciprocal-space plane (11 scans with a total of 382 data points). A very good global fit using Eqs. 2 and 4 is obtained with a residual $`\chi ^2=1.6`$. The solid lines in Fig. 2a–c were calculated from the globally optimized parameter set. The shaded areas represent the contribution of each mode, and the dashed lines show the background level. Figure 4 (right) shows the obtained $`b`$-axis dispersion relation (solid line), that has a minimum at $`k=0.5`$. Symbols in this figure represent the excitation energies obtained in fits to individual scans. A similar global fit (4 scans, 149 data points, $`\chi ^2=1.7`$) was applied to all constant-$`Q`$ scans measured in the $`(h,0,l)`$ scattering plane (solid lines in Fig. 2d and e). In this case the excitation energies at the zone-boundary $`(0,0,0.5)`$ were fixed to the values obtained from the $`(0,k,l)`$-plane global fit. Only the $`a`$-axis exchange constant $`J_x`$ was refined. Dispersion of magnetic excitations along this direction is barely detectable. The dispersion relation obtained from our fits is nontheless plotted in a solid line in Fig. 4 (left). The fitting analysis suggests a shallow minimum at $`h=0.5`$. From Eq. 5 we can thus guess that the global minimum of the 3D dispersion (3D AF zone-center) is located at $`(0.5,0.5,0.5)`$. The energy gap for $`c`$-axis polarized excitations was determined by fitting the model cross section to the wide-range constant-$`Q`$ scan measured with the horizontally-focusing configuration (Fig. 3, solid line). In this procedure the values of $`\mathrm{\Delta }_x`$ and $`\mathrm{\Delta }_y`$ were fixed. The parameters determined by the analysis described above are as follows: $`\mathrm{\Delta }_x=0.42(0.03)`$ meV, $`\mathrm{\Delta }_y=0.52(0.06)`$ meV, $`\mathrm{\Delta }_z=1.86(0.1)`$ meV, $`J_y=1.8(0.4)10^3`$ meV, and $`J_x=3.5(3.0)10^4`$ meV.
The in-chain exchange constant $`J`$ and anisotropy parameter $`D`$ are readily obtained from the measured gap energies: $`J0.81(\mathrm{\Delta }_x+\mathrm{\Delta }_y+\mathrm{\Delta }_z)=2.28`$ meV, $`D\frac{1}{4}(2\mathrm{\Delta }_z\mathrm{\Delta }_x\mathrm{\Delta }_y)=0.70`$ meV, and $`D/J=0.30`$, which is consistent with the bulk susceptibility result of Ref. The relative strength of inter-chain interactions are $`J_y/J710^4`$ and $`J_x/J1.310^4`$. These ratios are very similar to those found in NENP.
## IV Summary
The small value of critical fields in NDMAP, compared to those in NENP, are a result of smaller in-chain exchange interactions and somewhat larger easy-plane magnetic anisotropy. The lowest-energy excitation in NDMAP is polarized along the crystallograhic $`a`$–axis. The 3D magnetic zone-center appears to be $`(0.5,0.5,0.5)`$. Future experimental studies of the high-field phase should thus concentrate on this region of reciprocal space.
###### Acknowledgements.
We would like to thank S. M. Shapiro for illuminating discussions, S.-H. Lee for his assistance with the measurements at NIST, and R. Rothe for technical support at BNL. Work at Brookhaven National Laboratory was carried out under Contract No. DE-AC02-98CH10886, Division of Material Science, U.S. Department of Energy. Work at JHU was supported by the NSF through DMR-9801742. This work used instrumentation supported by NIST and the NSF through DMR-9423101. Work at RIKEN was supported in part by a Grant-in-Aid for Scientific Research from the Japanese Ministry of Education, Science, Sports and Culture.
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# 1. Introduction.
## 1. Introduction.
By a classical result in graph theory, the number of labeled trees<sup>1</sup><sup>1</sup>1Precise definitions for this and the following terms can be found in Section 2. on $`n1`$ vertices is $`n^{n2}`$. We endow the set $`𝒯_n`$ of labeled trees on $`n1`$ vertices with the uniform probability, giving weight $`n^{2n}`$ to each tree.
Each tree in $`𝒯_n`$ comes with its incidence matrix, the $`n\times n`$ matrix with entry $`ij`$ equal to $`1`$ if there is an edge between vertices $`i`$ and $`j`$ and to $`0`$ else. Each such (symmetric) matrix has $`n`$ (real) eigenvalues, which by definition form the spectrum of the corresponding tree. This leads in turn to $`nn^{n2}=n^{n1}`$ eigenvalues counted with multiplicities for $`𝒯_n`$ as a whole. In the sequel, we wish concentrate on the multiplicity of the eigenvalue $`0`$. Let $`Z(T)`$ be the multiplicity of the eigenvalue $`0`$ in the spectrum of the incidence matrix of the tree $`T`$, i.e. the dimension of the kernel. For each $`n1`$, the restriction $`Z_n`$ of $`Z`$ to $`𝒯_n`$ is a random variable. We set $`z_n=_{T𝒯_n}Z_n(T)`$. The expectation of $`Z_n(T)`$ is $`𝔼(Z_n)=z_n/n^{n2}`$.
To illustrate these definitions, we give a direct counting of $`z_1,\mathrm{},z_4`$ in appendix A.
Our aim is to prove :
###### Theorem 1.
Let $`z_n`$ be the total multiplicity of the eigenvalue $`0`$ in the spectra of the $`n^{n2}`$ labeled trees on $`n`$ vertices. Then :
i) Closed form :
$`z_n`$ $`=`$ $`n^{n1}2{\displaystyle \underset{2mn}{}}(1)^mn^{nm}m^{m2}\left({\displaystyle \genfrac{}{}{0pt}{}{n1}{m1}}\right)`$
$`{\displaystyle \frac{z_n}{n^{n2}}}`$ $``$ $`𝔼(Z_n)=n\left(12{\displaystyle \underset{2mn}{}}{\displaystyle \frac{(1)^m}{m}}\left({\displaystyle \frac{m}{n}}\right)^m\left({\displaystyle \genfrac{}{}{0pt}{}{n}{m}}\right)\right).`$
ii) Formal power series identity :
$$x^2+2xxe^x=\underset{n1}{}\frac{z_n}{n!}\left(xe^xe^{xe^x}\right)^n.$$
and
###### Corollary 2.
For large $`n`$, $`𝔼(Z_n)`$ has an asymptotic expansion in powers of $`1/n`$, whose first two terms are
$$𝔼(Z_n)=(2x_{}1)n+\frac{x_{}^2(x_{}+2)}{(x_{}+1)^3}+O(1/n)$$
where $`x_{}=0.5671432904097838729999\mathrm{}`$ is the unique real root of $`x=e^x`$. In particular, the average fraction of the spectrum occupied by the eigenvalue $`0`$ in a large random tree is asymptotic to $`2x_{}1=0.1342865808195677459999\mathrm{}`$.
Remark 3. We do not try to justify here that fluctuations in random trees become small when the number of vertices is large. However, it is expected that $`𝔼(Z_n^2)𝔼(Z_n)^2`$ grows only linearly with the number of vertices, so that in an appropriate sense the fraction of the spectrum occupied by the eigenvalue $`0`$ in an infinite random tree is $`2x_{}1`$ with probability 1.
Remark 4. With the explicit formula above, it is easy to list the first terms in the sequence $`(z_n)_{n1}`$, which are
$$1,0,3,8,135,1164,21035,322832,7040943,153153620,4048737099,\mathrm{}$$
To prove part i) of Theorem 1 we establish a few preparatory lemmas of independent interest. Then we prove ii) using Lagrange inversion and get Corollary 2 with the steepest descent method.
But first, we need to fix conventions and notations.
## 2. Definitions.
Even if we are interested ultimately only in trees, we shall need more general graphs (for instance, forests) in the proofs, so we give for the sake of completeness a collection definitions. Most of them are standard, and the reader is encouraged to skip this this section and come back to it only when needed. The fundamental definition is
Definition 5. A simple graph $`G`$ is a pair $`(V,E)`$ where $`V`$ is a finite set called the set of vertices and $`E`$ is a subset of $`V^{(2)}\{\{x,y\},xV,yV,xy\}`$ called the set of edges.
Remark 6. The adjective simple refers to the fact that there is at most one edge between two vertices and that edges are pairs of distinct vertices. As we have no use of more general graphs in the sequel, we shall from now on use graph for simple graph.
Definition 7. If $`V`$ is empty, then we say that the graph $`G`$ is empty. If $`\{x,y\}`$ belongs to $`E`$, we say that there is an edge between $`x`$ and $`y`$ and that $`x`$ and $`y`$ are adjacent vertices in $`G`$. The vertices adjacent to a given vertex $`x`$ are called the neighbors of $`x`$. The number of neighbors of a vertex $`x`$ is called the degree of $`x`$. A leaf of $`G`$ is a vertex of degree $`1`$. Two edges of $`G`$ with a common vertex are called adjacent edges
Definition 8. A labeled graph on $`n1`$ vertices is a graph with vertex set $`[n]=\{1,\mathrm{},n\}`$. The incidence matrix of a labeled graph on $`n`$ vertices is the $`n\times n`$ matrix with entry $`ij`$ equal to $`1`$ if there is an edge between vertices $`i`$ and $`j`$ and to $`0`$ else.
Remark 9. If the graph $`G`$ has $`|V|=n1`$ vertices<sup>2</sup><sup>2</sup>2For any finite set $`S`$ , $`|S|`$ is the number of elements in $`S`$., any bijection between $`V`$ and $`[n]`$ defines a labeled graph. The incidence matrices for different bijections differ only by a permutation of the lines and columns. In particular the eigenvalues are independent of the bijection. They are real because, by construction, incidence matrices are symmetric.
Definition 10. The spectrum of a graph is the set of eigenvalues (counted with multiplicities) of any of the associated incidence matrices. By convention, the spectrum of the empty graph is empty.
Definition 11. A subgraph of a graph $`G=(V,E)`$ is a graph $`(W,F)`$ such that $`WV`$ and $`FE`$. An induced subgraph of $`G`$ is a graph $`(W,F)`$ such that $`WV`$ and $`F=EW^{(2)}`$.
Definition 12. We say that two vertices $`x`$ and $`x^{}`$ $`V`$ are in the same component of $`G`$ if there is a sequence $`x=x_1,\mathrm{},x_n=x^{}`$ in $`V`$ such that adjacent terms in the sequence are adjacent in $`G`$ (taking $`n=1`$ shows that luckily $`x`$ and $`x`$ are in the same component). This gives a partition of $`V`$. Each component defines an induced subgraph of $`G`$ which is called a connected component of $`G`$. Then $`G`$ can be thought of as the disjoint union of its connected components. We say that $`G`$ is connected if it has only one connected component.
Definition 13. A polygon in a graph $`G`$ is a sequence $`x_0,x_1,\mathrm{},x_n`$, $`n3`$ of vertices such that adjacent terms in the sequence are adjacent in $`G`$, $`x_0=x_n`$ and $`x_1,\mathrm{},x_n`$ are distinct.
Definition 14. A forest is a graph without polygons. A tree is a non-empty connected forest.
Remark 15. Clearly a subgraph of a forest is a forest. The connected components of a non empty forest are trees. One shows easily that that a tree with $`n2`$ vertices has at least two leaves. Then a simple induction shows that a tree is exactly a connected graph for which the number of vertices is $`1`$ plus the number of edges. A classical theorem of Cayley states that there are $`n^{n2}`$ labeled trees on $`n`$ vertices(see for instance proposition 5.3.2 in ).
## 3. Two preparatory lemmas.
The first lemma is a characterization of the dimension of the kernel of incidence matrices viewed as a function on forests.
###### Lemma 16.
The function $`Z`$ which associates to any forest the multiplicity of the eigenvalue $`0`$ in its spectrum is characterized by the following properties :
i) The function $`Z`$ takes the value $`0`$ on $`\mathrm{}`$, the empty forest.
ii) The function $`Z`$ takes the value $`1`$ on , the forest with one vertex.
iii) The function $`Z`$ is additive on disjoint components, i.e. if the forest $`F`$ is the union of two disjoint forests $`F_1`$ and $`F_2`$ then $`Z(F)=Z(F_1)+Z(F_2)`$
iv) The function $`Z`$ is invariant under “leaf removal”, i.e. if $`x`$ is a leaf of $`F`$, $`y`$ its (unique) $`neighbor`$, $`V^{}=V\backslash \{x,y\}`$, and $`F^{}`$ is the subforest of $`F`$ induced by $`V^{}`$ then $`Z(F)=Z(F^{})`$.
Remark 17. That the function $`Z`$ fulfills properties i)–iv) was no doubt known decades ago (see for instance section 8.1, Hückels theory, in ). We give a proof, because in the sequel we want to emphasize and use the simple fact that these properties characterize the function $`Z`$.
Proof of Lemma 16. First, we show that the function $`Z`$ has properties i)–iv). In fact, this is true for general graphs (not only forests). Properties i) and ii) follow from the definition of $`Z`$, property iii) follows from the fact that the incidence matrix can be put in block diagonal form, each block corresponding to a connected component. Property iv) is only slightly more complicated. With an appropriate labeling of the vertices, the incidence matrix $`𝐌`$ of $`F`$ can be decomposed as
$$𝐌=\left(\begin{array}{ccc}0& 1& \mathrm{𝟎}\\ 1& 0& 𝐍\\ \mathrm{𝟎}& {}_{}{}^{t}𝐍& 𝐌^{}\end{array}\right)$$
where the first line and column are indexed by the leaf $`x`$, the second line and column is indexed by its neighbor $`y`$, $`𝐍`$ describes the edges between this neighbor and $`V^{}`$, and $`𝐌^{}`$ is the incidence matrix for $`V^{}`$. Then $`𝐯=^t(v_1,v_2,𝐯^{})`$ is in the kernel of $`𝐌`$ if and only if
$`v_2`$ $`=`$ $`0`$
$`v_1`$ $`=`$ $`\mathrm{𝐍𝐯}^{}`$
$`𝐌^{}𝐯^{}`$ $`=`$ $`^t𝐍v_2.`$
So $`v_2=0`$ which reported in the third equation gives $`𝐌^{}𝐯^{}=\mathrm{𝟎}`$ implying that $`𝐯^{}`$ is in the kernel of $`𝐌^{}`$, and then the second equation just tunes $`v_1`$ the appropriate value. So the kernels of $`𝐌`$ and $`𝐌^{}`$ have the same dimension. This proves iv).
Now, any tree with more than $`1`$ vertex has leaves, so leaf removal as defined in iv) allows to reduce the forest $`F`$ to a (possibly empty) family of isolated vertices (all connected components have only one vertex). Hence, there is at most one function, namely $`Z`$, that can satisfy properties i)–iv).
Remark 18. Leaf removal and additivity give an efficient algorithm to compute the multiplicity of the eigenvalue 0 for a given forest, especially when this forest is given as a drawing.
The next lemma gives a practically awful but theoretically useful formula for the function $`Z`$.
###### Lemma 19.
Let $`L`$ be the function on forests defined by
i’) The function $`L`$ takes value $`0`$ on $`\mathrm{}`$, the empty forest.
ii’) The function $`L`$ takes value $`1`$ on , the forest with one vertex.
iii’) The function $`L`$ takes value $`0`$ on disconnected forests.
iv’) The function $`L`$ takes value $`2(1)^{n1}`$ on trees with $`n2`$ vertices.
Then, for any forest $`F`$
$$Z(F)=\underset{F^{}F}{}L(F^{})=\underset{T^{}F}{}L(T^{})$$
where the first sum is over induced subforests of $`F`$, and the second over induced subtrees of $`F`$.
Remark 20. For a given forest, there is a much nicer formula, directly connected to the geometry of the forest (again, see for instance section 8.1, Hückels theory, in ). In fact, let $`Q(F)`$ be the maximum among the cardinals of sets of pairwise non-adjacent edges in $`F`$, and $`N(F)`$ be the number of vertices in $`F`$. Then $`Z(F)=N(F)2Q(F)`$. It is easy to show that $`N(F)2Q(F)`$ satisfies properties i)–iv) of Lemma 16. In particular, a possible way to maximize the number of non-adjacent edges in $`F`$ in the situation iv) is to do so on $`F^{}`$ and add the edge $`\{x,y\}`$. Anyway, this explicit formula allows us to restate our theorems in terms of the random variable $`Q_n`$, the restriction of $`Q`$ to $`𝒯_n`$. For instance, in a large random tree on $`n`$ vertices, one can find about $`(1x_{})n`$ pairwise non-adjacent edges. Note that $`1x_{}=0.4328567095902161270000\mathrm{}`$ is not much smaller than $`0.5`$ (the upper bound for $`Q(T)/N(T)`$ for a given tree because $`Z(T)=N(T)2Q(T)`$ is always nonnegative).
Proof of Lemma 19. Our strategy is to use the characterization of $`Z`$ in Lemma 16. First, we observe that the second equality is a trivial consequence of i’) and iii’). We define a new function $`Z^{}`$ on the set of forests by
$$Z^{}(F)\underset{T^{}F}{}L(T^{})$$
(where the sum is over induced subtrees of $`F`$) and show that $`Z^{}`$ satisfies properties i)–iv) of Lemma 16.
As the empty forest has no non-empty induced subtree i’) implies i).
In the same vein, the forest with one vertex as only one non-empty induced subtree, namely itself, so ii’) implies ii).
If the forest $`F`$ is the union of two disjoint forests $`F_1`$ and $`F_2`$, an induced subtree of $`F`$ is either an induced subtree of $`F_1`$ or an induced subtree of $`F_2`$, and the sum defining $`Z^{}(F)`$ splits as $`Z^{}(F_1)+Z^{}(F_2)`$, showing that $`Z^{}`$ satisfies property iii).
Now, if $`x`$ is a leaf of $`F`$ and $`y`$ its neighbor, we define $`V^{}=V\backslash \{x,y\}`$, $`V^{\prime \prime }=\{x,y\}`$ and consider $`F^{}`$ and $`F^{\prime \prime }`$, the subforests of $`F`$ induced by $`V^{}`$ and $`V^{\prime \prime }`$ respectively. We split the sum defining $`Z^{}(F)`$ in three pieces. The first is over the induced subtrees of $`F^{}`$. This is just the sum defining $`Z^{}(F^{})`$. The second is over the induced subtrees of $`F^{\prime \prime }`$, which is a tree on two vertices. Its subtrees are itself, with weight $`L(F^{\prime \prime })=2(1)^{21}=2`$, and two trees with one vertex, each with weight $`L(\text{})=1`$, so this second sum gives $`0`$. The third sum is over induced subtrees that have vertices in both $`V^{}`$ and $`V^{\prime \prime }`$. If this sum is not empty, every tree that appears in it has $`y`$ as a vertex (by connectivity) and has at least two vertices (because the tree consisting of $`y`$ alone has already been counted). Then we can group these trees in pairs, a tree containing $`x`$ being paired with the same tree but with $`x`$ and the edge $`\{x,y\}`$ deleted. The function $`L`$ takes opposite values on the two members of a pair, so the third sum contributes $`0`$. Hence $`Z^{}`$ satisfies property iV). So $`Z^{}(F)=Z^{}(F^{})`$.
Remark 21. These two lemmas have an obvious extension to bicolored forests. If we use black and white as colors, and count the zero eigenvectors having value zero on white vertices, we only need to replace ii) in Lemma 16 by
ii) The function $`Z`$ takes value $`1`$ on $``$, the forest with one vertex colored in black and $`0`$ on $``$, the forest with one vertex colored in white.
and ii’) and iii’) in Lemma 19 by
ii’) The function $`L`$ takes value $`1`$ on $``$, the forest with one vertex colored in black and $`0`$ on $``$, the forest with one vertex colored in white.
iii’) The function $`l`$ takes value $`(1)^{n1}`$ on trees with $`n2`$ vertices.
The proofs remain the same.
Remark 22. The formula
$$Z(F)=\underset{F^{}F}{}L(F^{})$$
can be inverted using inclusion-exclusion to give
$$L(F)=\underset{F^{}F}{}(1)^{|V(F)||V(F^{})|}Z(F^{}).$$
This identity has an application in random graph theory , which is why we got interested in Lemma 19 in the first place.
## 4. Main proofs.
We have now the necessary tools to prove theorem 1.
Proof of Theorem 1. By Lemma 16
$$z_n\underset{T𝒯_n}{}Z_n(T)=\underset{m=1}{\overset{n}{}}\underset{T𝒯_n}{}\underset{T^{}𝒯_m}{\overset{T^{}T}{}}L(T^{}).$$
As the function $`L`$ depends only on the number of vertices, for fixed $`m`$ the double sum $`_{T𝒯_n}_{T^{}𝒯_m}^{T^{}T}`$ is simply a multiplicity. We count this multiplicity as follows : we remove from $`T`$ the edges of $`T^{}`$, so we are left with m trees, each with a special vertex, the one belonging to $`T^{}`$. This is by definition what is called a planted forest (or rooted forest) with $`n`$ vertices and $`m`$ trees. The number of such objects is $`m\left(\genfrac{}{}{0pt}{}{n}{m}\right)n^{nm1}`$ (see for instance proposition 5.3.2 in ). Conversely, starting from such a planted forest with $`m`$ trees (each with a special vertex) and $`n`$ vertices, we can build a tree on the special vertices in $`m^{m2}`$ ways. So
$$\underset{T𝒯_n}{}\underset{T^{}𝒯_m}{\overset{T^{}T}{}}1=m^{m1}\left(\genfrac{}{}{0pt}{}{n}{m}\right)n^{nm1}.$$
Hence summation over $`m`$ gives
$$z_n=n^{n1}2\underset{2mn}{}(1)^mn^{nm1}m^{m1}\left(\genfrac{}{}{0pt}{}{n}{m}\right).$$
Simple rearrangements lead to the two equivalent formulæ in i), the first one making clear that $`z_n`$ is an integer.
To obtain the generating function in ii), we need a mild extension of the Lagrange inversion formula (see for instance section 5.4 in ), which states that if $`f(x)`$ is a formal power series in $`x`$ starting as $`f(x)=x+O(x^2)`$ and $`g(x)`$ is an arbitrary formal power series in $`x`$,
$$\left(gf^1\right)(t)=g(0)+\underset{n1}{}\frac{1}{n}\left[\frac{x^ng^{}(x)}{f(x)^n}\right]_{n1}t^n.$$
where $`[h(v)]_k`$ is by definition the $`k^{th}`$ coefficient of the formal power series $`h(v)`$.
As an immediate application, we see that if $`t=xe^x`$ then
$$x=\underset{m1}{}(m)^{m1}\frac{t^m}{m!}$$
and
$$xx^2/2=\underset{m1}{}(m)^{m2}\frac{t^m}{m!}.$$
Now we introduce $`y=te^t`$ and define a sequence $`z_n^{},n1`$ by
$$x^2+2xxe^x=\underset{n1}{}z_n^{}\frac{y^n}{n!},$$
but instead of applying directly the Lagrange inversion formula to $`y=xe^xe^{xe^x}`$, we first substitute the $`t`$-expansion (already obtained by Lagrange inversion) on the left-hand side which yields
$$2\underset{m1}{}(m)^{m2}\frac{t^m}{m!}t,$$
and then apply Lagrange inversion with $`y=te^t`$. The result is
$$\frac{z_n^{}}{n!}=\frac{1}{n}\left[e^{nt}\left(12\underset{m2}{}\frac{(m)^{m2}}{(m1)!}t^{m1}\right)\right]_{n1}.$$
Straightforward expansion of this formula shows that $`z_n^{}=z_n`$, and this proves the generating function representation in ii).
Remark 23. The derivation of ii) is quite artificial. It turns out that random graph theory gives a natural proof using the formula mentioned in remark 22.
Proof of Corollary 2. This time we use Lagrange inversion with $`y=xe^xe^{xe^x}`$, in a contour integral representation<sup>3</sup><sup>3</sup>3We include the factor $`\frac{1}{2i\pi }`$ in the symbol $``$.. So
$$\frac{z_n}{n!}=\frac{1}{n}\frac{dx}{(xe^xe^{xe^x})^n}(1+x)(2e^x),$$
where the contour is a small anticlockwise-oriented circle around the origin. For large $`n`$ we use the steepest descent method to obtain the asymptotic expansion of $`z_n`$. As $`\frac{d}{dx}xe^xe^{xe^x}=(1+x)(1xe^x)e^xe^{xe^x}`$, the saddle points of $`xe^xe^{xe^x}`$ are $`x=1`$ and the solutions to $`x=e^x`$. This equation has a unique real root, $`x_{}`$, which is positive. Numerically, $`x_{}=0.5671432904097838729999\mathrm{}.`$ On the other hand, $`x=e^x`$ has an infinite number of complex solutions, coming in complex conjugate pairs. Asymptotically, the imaginary parts of these zeroes are evenly spaced by about $`2\pi `$, while their real parts are negative and grow logarithmically in absolute value. Consideration of the landscape produced by the modulus of the function $`xe^xe^{xe^x}`$ shows that the small circle around the origin can be deformed to give the union of two steepest descent curves, one passing through $`x=1`$ and the other through $`x=x_{}`$. These two curves are asymptotic to the two lines $`y=\pm \pi `$ at $`x+\mathrm{}`$. Hence, despite the fact that the value of $`xe^xe^{xe^x}`$ is the same, namely $`1/e`$, at all the complex saddle points and at $`x_{}`$, the complex saddle points do not contribute to the asymptotic expansion of $`z_n`$ at large $`n`$. Moreover, the point $`x=1`$ only gives subdominant contributions because $`e^1e^{e^1}`$ is larger than $`1/e`$ in absolute value. So we concentrate on the asymptotic expansion around $`x_{}`$. As
$$\mathrm{log}xe^xe^{xe^x}=1\frac{(x_{}+1)}{2x_{}}(xx_{})^2+O((xx_{})^3)$$
we infer that
$$e^n\sqrt{2\pi n}\frac{dx}{(xe^xe^{xe^x})^n}(1+x)(2e^x)$$
has an asymptotic expansion in powers of $`1/n`$. Hence, by use of Stirling’s formula for $`n!`$, we conclude that $`𝔼(Z_n)=z_n/n^{n2}`$ has an asymptotic expansion in powers of $`1/n`$. The first two terms are obtained by brute force.
## Appendix A Examples of direct multiplicity counting.
This appendix gives the counting of trees and multiplicities of $`0`$ in the spectrum for trees on $`n=1,2,3`$ or $`4`$ vertices.
Example 24. For $`n=1`$ there is only one tree, , and one way to label it, giving a total of $`1=1^{12}`$ tree on one vertex. The incidence matrix is $`(0)`$, so the eigenvalue $`0`$ occurs with multiplicity $`z_1=1`$.
Example 25. For $`n=2`$ there is only one tree, , and one way to label it, giving a total of $`1=2^{22}`$ tree on two vertices. The incidence matrix is
$$\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),$$
so the eigenvalue $`0`$ occurs with multiplicity $`z_2=0`$.
Example 26. For $`n=3`$ there is only one tree, , and three ways to label it, giving a total of $`3=3^{32}`$ trees on three vertices. Up to permutation of rows and columns, the incidence matrix for each of these three labeled trees is
$$\left(\begin{array}{ccc}0& 1& 0\\ 1& 0& 1\\ 0& 1& 0\end{array}\right),$$
which has zero as an eigenvalue with multiplicity $`1`$ (a corresponding eigenvector is $`{}_{}{}^{t}(1,0,1)`$), so there is a total of $`3\times 1`$ zero eigenvalues, and $`z_3=3`$
Example 27. For $`n=4`$ there are two trees, ($`12`$ ways to label it), and ($`4`$ ways to label it), giving a total of $`12+4=16=4^{42}`$ trees on three vertices. Up to permutation of rows and columns, the two incidence matrices are
$$\left(\begin{array}{cccc}0& 1& 0& 0\\ 1& 0& 1& 0\\ 0& 1& 0& 1\\ 0& 0& 1& 0\end{array}\right)\text{and}\left(\begin{array}{cccc}0& 1& 0& 0\\ 1& 0& 1& 1\\ 0& 1& 0& 0\\ 0& 1& 0& 0\end{array}\right).$$
The first does not have $`0`$ as an eigenvalue, whereas the second has zero as an eigenvalue with multiplicity $`2`$ (corresponding eigenvectors are for instance $`{}_{}{}^{t}(1,0,1,0)`$ and $`{}_{}{}^{t}(1,0,0,1)`$), so there is a total of $`12\times 0+4\times 2`$ zero eigenvalues, and $`z_4=8`$.
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# Capture of 𝛼 Particles by Isospin-Symmetric Nuclei
## I Introduction
The astrophysical importance of $`\alpha `$ capture on target nuclei with $`N=Z`$ is manifold. In the Ne- and O-burning phase of massive stars, alpha capture reaction sequences are initiated at <sup>24</sup>Mg and <sup>28</sup>Si, respectively, and determine the abundance distribution prior to the Si-burning phase . Nucleosynthesis in explosive Ne and explosive O burning in type II supernovae depend on reaction rates for $`\alpha `$ capture on <sup>20</sup>Ne to <sup>36</sup>Ar . While many of these processes are in quasistatistical equilibrium , the reaction rates itself are important for a reliable description of nucleosynthesis in the subsequent cooling phase. An $`\alpha `$ capture chain on such self-conjugate nuclei actually determines the production of <sup>44</sup>Ti , which contributes to the supernova light curve by the energy release from its $`\beta `$ decay to <sup>44</sup>Ca via <sup>44</sup>Sc. Charged-particle reaction networks have to consider $`\alpha `$-capture rates in the conditions of an $`\alpha `$-rich freeze-out , or for an extended rp-process for proper treatment of the $`\gamma `$ induced $`\alpha `$ break-up of selfconjugate nuclei . Such calculations involve also highly unstable nuclei, thus calling for a reliable prediction of the respective reaction rates.
Due to isospin selection rules, $`E1`$ $`\gamma `$ transitions with isospin $`T=0T=0`$ are forbidden. Likewise, $`M1`$ transitions will be strongly suppressed. Because of isospin conservation, only states with isospin $`T=0`$ can be populated by $`\alpha `$ capture on $`N=Z`$ ($`T=0`$) targets. This leads to a strong suppression of the $`\gamma `$ transitions in the compound nucleus and thus of the ($`\alpha `$,$`\gamma `$) cross section of self-conjugate nuclei .
Previous theoretical work used in astrophysical calculations either neglected this isospin effect or accounted for it only in a phenomenological way with arbitrary suppression factors . In this work we want to improve on the prediction of reaction rates on isospin-symmetric targets and compare the theoretical values to newly compiled available experimental information. In section II the statistical model of nuclear reactions is introduced and the method for obtaining isospin suppression factors as well as the resulting theoretical cross sections and reaction rates are presented. Section III reviews the available experimental information and presents reaction rates newly derived from all available experimental data covering different reaction channels. A summary and conclusion will be given in Section IV.
## II The Statistical Model
### A Introduction
The majority of nuclear reactions in astrophysics can be described in the framework of the statistical model (compound nucleus mechanism, Hauser-Feshbach approach, HF; see e.g. ), provided that the level density of the compound nucleus is sufficiently large in the contributing energy window . This description assumes that the reaction proceeds via a compound nucleus which finally decays into the reaction products. With a sufficiently high level density, average cross sections
$$\sigma ^{\mathrm{HF}}=\sigma _{\mathrm{form}}b_{\mathrm{dec}}=\sigma _{\mathrm{form}}\frac{\mathrm{\Gamma }_{\mathrm{final}}}{\mathrm{\Gamma }_{\mathrm{tot}}}$$
(1)
can be calculated which can be factorized into a cross section $`\sigma _{\mathrm{form}}`$ for the formation of the compound nucleus and a branching ratio $`b_{\mathrm{dec}}`$. This branching ratio describes the probability of the decay into the channel of interest relative to the total decay probability into all possible exit channels. The partial widths $`\mathrm{\Gamma }`$ as well as $`\sigma _{\mathrm{form}}`$ are related to (averaged) transmission coefficients, which comprise the central quantities in any HF calculation.
Many nuclear properties enter the computation of the transmission coefficients: mass differences (separation energies), optical potentials, GDR widths, level densities. The transmission coefficients can be modified due to pre-equilibrium effects which are included in width fluctuation corrections (see also , and references therein) and by isospin effects. It is in the description of the nuclear properties where the various HF models differ.
In astrophysical applications usually different aspects are emphasized than in pure nuclear physics investigations. Many of the latter in this long and well established field were focused on specific reactions. All or most ”ingredients”, like optical potentials for particle transmission coefficients, level densities, resonance energies and widths of giant resonances for predicting E1 and M1 $`\gamma `$-transitions, were deduced from experiments. As long as the statistical model prerequisites are met, this will produce highly accurate cross sections. For the majority of nuclei in astrophysical applications such information is not available. The real challenge is thus not the well-established statistical model, but rather to provide all these necessary ingredients in as reliable a way as possible, also for nuclei where none of such information is available.
### B The NON-SMOKER Code
For the calculations presented in this work we utilized the recently developed statistical model code NON-SMOKER . The current status of the code is outlined in the following.
For neutrons and protons the optical potential of Jeukenne, Lejeune & Mahaux is used, which is based on microscopic infinite nuclear matter calculations for a given density, applied with a local density approximation. It includes corrections of the imaginary part . The potential of McFadden & Satchler is used for $`\alpha `$ particles, which is based on extensive data. Deformed nuclei are treated by an effective spherical potential of equal volume (see e.g. ).
The level density treatment has been recently improved . Additionally, experimental level information (excitation energies, spins, parities) are included , as well as experimental nuclear masses .
The $`\gamma `$-transmission coefficients have to include the dominant E1 and M1 $`\gamma `$ transitions. The smaller, less important M1 transitions have usually been treated with the simple single particle approach $`TE^3`$ (see e.g. ). The E1 transitions are usually calculated on the basis of the Lorentzian representation of the Giant Dipole Resonance (GDR). Many microscopic and macroscopic models have been devoted to the calculation of GDR energies and widths. An excellent fit to the GDR energies is obtained with the hydrodynamic droplet model . An improved microscopic-macroscopic approach is used, based on dissipation and the coupling to quadrupole surface vibrations ; see also . Most recently it was shown that the inclusion of “soft mode” or “pygmy” resonances might have important consequences on the E1 transitions in neutron-rich nuclei far off stability. The pygmy resonances could be caused by a neutron skin which generates soft vibrational modes . It is still under discussion whether such modes exist. As the effect is negligible for nuclei close to stability, we did not include it in the calculations presented here.
### C Inclusion of Isospin Effects
The original Hauser-Feshbach equation implicitly assumes complete isospin mixing but can be generalized to explicitly treat the contributions of the dense background states with isospin $`T^<=T^{\mathrm{g}.\mathrm{s}.}`$ and the isobaric analog states with $`T^>=T^<+1`$ . In reality, compound nucleus states do not have unique isospin and for that reason an isospin mixing parameter $`\mu `$ was introduced , which is the fraction of the width of $`T^>`$ states leading to $`T^<`$ transitions; for complete isospin mixing $`\mu =1`$, for pure $`T^<`$ states $`\mu =0`$. In the case of overlapping resonances for each involved isospin, $`\mu `$ is directly related to the level densities $`\rho ^<`$ and $`\rho ^>`$, respectively. Isolated resonances can also be included via their internal spreading width $`\mathrm{\Gamma }^{}`$ and a bridging formula was derived to cover both regimes .
In order to determine the mixing parameter $`\mu =\mu (E)`$, experimental information for excitation energies of $`T^>`$ levels is used where available in the code NON-SMOKER. Experimental values for spreading widths are also tabulated . Similarly to the standard treatment for the $`T^<`$ states, a level density description is invoked above the last experimentally known $`T^>`$ level. Since the $`T^>`$ states in a nucleus ($`Z`$,$`N`$) are part of a multiplet, they can be approximated by the levels (and level density) of the nucleus ($`Z`$$``$1,$`N`$+1), only shifted by a certain energy $`E_\mathrm{d}`$. This displacement energy $`E_\mathrm{d}`$ can be calculated and it is dominated by the Coulomb displacement energy: $`E_\mathrm{d}=E_\mathrm{d}^{\mathrm{Coul}}+ϵ`$. In the absence of experimental level information, we use the formula from Woosley & Fowler as given by for the determination of the excitation energy of the first isobaric analog state.
The inclusion of the explicit treatment of isospin has two major effects on statistical cross section calculations in astrophysics: the suppression of $`\gamma `$ widths for reactions involving self-conjugate nuclei and the suppression of the neutron emission in proton-induced reactions. This paper focuses on the suppression of the $`\gamma `$ width in $`\alpha `$ capture reactions. Non-statistical effects, i.e. the appearance of isobaric analog resonances, will not be further discussed here.
The isospin selection rule for $`E1`$ transitions is $`\mathrm{\Delta }T=0,1`$ with transitions $`00`$ being forbidden. In the case of ($`\alpha `$,$`\gamma `$) reactions on targets with $`N=Z`$, the cross sections will be heavily suppressed because $`T=1`$ states cannot be populated in the compound nucleus due to isospin conservation. A suppression will also be found for capture reactions leading into self-conjugate nuclei, although somewhat less pronounced because $`T=1`$ states can be populated according to the isospin coupling coefficients.
In previous reaction rate calculations the suppression of the $`\gamma `$ widths was treated completely phenomenologically by dividing the total $`\gamma `$ widths (and thus the cross section) by quite uncertain factors of 5 and 2, for ($`\alpha `$,$`\gamma `$) reactions on self-conjugate nuclei and nucleon capture reactions going into self-conjugate nuclei, respectively, regardless of the target’s nature. These empirical factors were estimated from the scarce experimental data available at that time.
We are replacing these factors by including more isospin information into the calculation of the $`E1`$ and $`M1`$ suppression. It can be shown (see e.g. ) that $`00`$ $`E1`$ transitions are forbidden. An approximate suppression rule for $`\mathrm{\Delta }T=0`$ transitions in self-conjugate nuclei can also be derived for $`M1`$ transitions and leads to a suppression factor of about 1/150. The total suppression of the $`E1`$ transitions would be exact if isospin were an exact quantum number, at least in the simplifying limit that the wavelength of the transition involved is large compared with the nuclear size. However, the Coulomb force mixes states of different isospin to a small extent, so the “clean” selection rule becomes a suppression factor, too. The theoretical estimate for the factor as given by Ref. is 0.01.
Since only $`T^<=0`$ states can be populated by $`\alpha `$ capture, we need to know the mixing of these $`T^<`$ into $`T^>`$ states. Below the first isobaric analog state at $`E_\mathrm{d}`$, states should be unmixed. We describe that phenomenologically by suppressing both $`E1`$ and $`M1`$ transitions. The suppression factor $`f_{\mathrm{iso}}`$ is set equal for $`E1`$ and $`M1`$.
It is clear that the total suppression factor of the resulting $`\gamma `$ width must be related to the number of resonances in the compound nucleus. Therefore, in a straightforward way we set the suppression for $`\gamma `$ transitions from the excitation energy $`E`$ of the compound nucleus proportional to the ratio of the density of $`T=1`$ levels and the density of $`T=0`$ levels,
$$f_{\mathrm{iso}}(E)\frac{\rho ^>(E)}{\rho ^<(E)}.$$
(2)
We find a weak energy dependence of $`f_{\mathrm{iso}}`$.
### D Results
The cross sections and reaction rates for $`\alpha `$ capture were calculated for the $`N=Z`$ isotopes from Mg to Mo ($`12Z42`$). The NON-SMOKER results including the above description of isospin suppression of the $`\gamma `$ width are given in Tables VII to XII. In the next section we compare these results to experimental cross sections and rates, either directly measured or calculated from resonance parameters.
## III Comparison to experimental information
### A Experimental data
Experimental data on $`\alpha `$ capture reactions on self-conjugate nuclei 20$``$A$``$40 are rather sparse. To determine reliably the reaction rate at various stellar temperatures T<sub>9</sub> (in GK) detailed information on the number of contributing resonances $`n`$, the resonance strengths $`\omega \gamma _i`$ (in units eV), and the resonance energies E<sub>i</sub> (in units MeV) are necessary:
$$N_A<\sigma v>=1.5410^5A^{3/2}T_9^{3/2}\underset{i}{\overset{n}{}}\omega \gamma _ie^{\frac{11.605E_i}{T_9}},$$
(3)
with A as the reduced mass of the system. The resonance strength depends on the spin of the resonance $`J`$ and the partial widths of the entrance $`\mathrm{\Gamma }_\alpha `$ and exit channel $`\mathrm{\Gamma }_\gamma `$, and the total width $`\mathrm{\Gamma }`$
$$\omega \gamma =(2J+1)\frac{\mathrm{\Gamma }_\alpha \mathrm{\Gamma }_\gamma }{\mathrm{\Gamma }}.$$
(4)
Extensive resonance studies in the astrophysically relevant low energy range E$`{}_{\alpha }{}^{}`$1.5 MeV have only been made for <sup>20</sup>Ne($`\alpha `$,$`\gamma `$)<sup>24</sup>Mg , <sup>24</sup>Mg($`\alpha `$,$`\gamma `$)<sup>28</sup>Si , and also for <sup>28</sup>Si($`\alpha `$,$`\gamma `$)<sup>32</sup>S . For the reaction <sup>32</sup>S($`\alpha `$,$`\gamma `$)<sup>36</sup>Ar only insufficient data are available on resonances in the energy range E$`{}_{\alpha }{}^{}`$2.2 MeV . The experimentally observed level density in the compound nucleus <sup>36</sup>Ar indicates a significantly higher number of possible resonances . Even less information is available on resonances in <sup>36</sup>Ar($`\alpha `$,$`\gamma `$)<sup>40</sup>Ca. Resonance measurements have been performed in the energy range of the giant dipole resonance E<sub>α</sub>=6-17 MeV. Several strong resonances have been observed and the resonance strengths were determined . Measurements at lower energies, however, E<sub>α</sub>=3-6 MeV, did not yield any significant information on possible resonance states in this range . More experimental data are again available for the reaction <sup>40</sup>Ca($`\alpha `$,$`\gamma `$)<sup>44</sup>Ti. Several resonances were successfully measured in the energy range E<sub>α</sub>=2.75-4 MeV , additional measurements were performed in the energy range E<sub>α</sub>=4-6 MeV . Again, no information is available on lower energy resonances.
Due to the lack of low energy resonance information a direct comparison of the predicted HF rates and the rates obtained by the experimental data is only of limited use. However, additional and complementary information can be gained from resonance $`\alpha `$ elastic scattering measurements and lifetime measurements which yield information about the total width $`\mathrm{\Gamma }`$ of the resonance state. More important, however, are (p,$`\alpha )`$ studies with a self-conjugate compound nucleus. These measurements yield extensive information about the existence and characteristics of possible resonances in the $`\alpha `$ capture channel. In particular reactions like <sup>23</sup>Na(p,$`\alpha `$)<sup>20</sup>Ne , <sup>27</sup>Al(p,$`\alpha `$)<sup>24</sup>Mg , <sup>31</sup>P(p,$`\alpha `$)<sup>28</sup>Si <sup>35</sup>Cl(p,$`\alpha `$)<sup>32</sup>, and <sup>39</sup>K(p,$`\alpha `$)<sup>36</sup>Ar populate $`\alpha `$-unbound natural parity states in the self-conjugate compound nuclei and therefore complement the direct experimental information on resonances in <sup>20</sup>Ne($`\alpha ,\gamma `$)<sup>24</sup>Mg, <sup>24</sup>Mg($`\alpha ,\gamma `$)<sup>28</sup>Si, <sup>28</sup>Si($`\alpha ,\gamma `$)<sup>32</sup>S, <sup>32</sup>S($`\alpha ,\gamma `$)<sup>36</sup>Ar, and <sup>36</sup>Ar($`\alpha ,\gamma `$)<sup>40</sup>Ca, respectively. Particularly useful is information on the resonance strength,
$$\omega \gamma _{(p,\alpha )}=\frac{2J+1}{(2j_p+1)(2j_t+1)}\frac{\mathrm{\Gamma }_p\mathrm{\Gamma }_\alpha }{\mathrm{\Gamma }}$$
(5)
with $`j_p`$ and $`j_t`$ as projectile and target spin. Combined with information on the total width
$$\mathrm{\Gamma }\mathrm{\Gamma }_p+\mathrm{\Gamma }_\alpha $$
(6)
and on the resonance strengths of the correlated natural parity (p,$`\gamma `$) resonances
$$\omega \gamma _{(p,\gamma )}=\frac{2J+1}{(2j_p+1)(2j_t+1)}\frac{\mathrm{\Gamma }_p\mathrm{\Gamma }_\gamma }{\mathrm{\Gamma }}.$$
(7)
Such information allows to fit the partial widths $`\mathrm{\Gamma }_p`$, $`\mathrm{\Gamma }_\gamma `$, $`\mathrm{\Gamma }_\alpha `$ to match the total width and the observed (p,$`\gamma `$) and (p,$`\alpha `$) strengths. The ($`\alpha ,\gamma `$) resonance strength is calculated from these values. If only partial data about (p,$`\gamma `$) or (p,$`\alpha `$) resonance strengths and the total width is available, we adopted an alpha spectroscopic factor from empirical alpha-strength studies in self-conjugate nuclei in this mass range and calculated the $`\alpha `$ width of the level in terms of a simple potential model. If only either the total width, or one of the (p,$`\gamma `$) or (p,$`\alpha `$) strengths is available, both the alpha as well as the proton partial width need to be calculated. For the latter case the single particle spectroscopic factor was adopted from the average of the single particle strength distribution in the excitation range of this nucleus. For natural parity levels with no spectroscopic information available, the gamma width was calculated additionally. We used an average Weisskopf strength which has been matched to the known gamma strength distribution of the neighboring states to account empirically for the isospin $`\gamma `$ strength suppression in self-conjugate nuclei. Since there are no data available for proton capture on the short-lived <sup>43</sup>Sc, the strengths of low energy resonances in <sup>40</sup>Ca($`\alpha ,\gamma `$)<sup>44</sup>Ti are based purely on such estimates. The results are listed and compared with the available experimental data in Table I to VI. It can clearly be seen that there is systematically good agreement between the calculated and experimental resonance parameters.
### B Reaction Rates
The resonance parameters derived and discussed in the previous chapter allow us to calculate the reaction rate N$`{}_{A}{}^{}<\sigma v>`$ as a function of temperature using Equation 2. These rates are directly compared with the reaction rates N$`{}_{A}{}^{}<\sigma v>_{\mathrm{HF}}`$ based on the HF calculations with the code NON-SMOKER in Tables VII to XII. Shown is the ’experimental’ reaction rate N$`{}_{A}{}^{}<\sigma v>_{\mathrm{exp}}`$ which is based on the few directly observed ($`\alpha ,\gamma `$) resonances only. The ’empirical’ rate N$`{}_{A}{}^{}<\sigma v>_{\mathrm{emp}}`$ is calculated using Equation 2 on the basis of the extended resonance set and the associated level parameter analysis discussed in the previous section. The ratios of the ’experimental’ to the ’empirical’ rates is shown in Figs. 1$``$6. For the reactions <sup>20</sup>Ne($`\alpha ,\gamma `$)<sup>24</sup>Mg, <sup>24</sup>Mg($`\alpha ,\gamma `$)<sup>28</sup>Si, and <sup>28</sup>Si($`\alpha ,\gamma `$)<sup>21</sup>S a considerable amount of experimental alpha capture data is available for higher energies , therefore at temperatures above T$``$0.3 GK the experimental rates are in good agreement with the empirical rates, only at lower temperatures the experimental rates are substantially lower due to the lack of experimental data. Only very limited data on alpha capture resonances are available for <sup>32</sup>S($`\alpha ,\gamma `$)<sup>36</sup>Ar and <sup>36</sup>Ar($`\alpha ,\gamma `$)<sup>40</sup>Ca, this explains the substantial deviation between experimental and empirical rate over the entire temperature range.
The ratios of the statistical model rates discussed in section II and empirical rates discussed in section III are shown in Fig. 7 for the astrophysically relevant temperature region. Except for <sup>32</sup>Si($`\alpha ,\gamma `$)<sup>36</sup>Ar and <sup>36</sup>Ar($`\alpha ,\gamma `$)<sup>40</sup>Ca at temperature T$`<2`$ GK the deviations remain in the range of a factor of $``$3. This is to be expected from a global statistical model calculation. Towards higher temperature ($`T3.5`$ GK) the calculation is significantly improved in respect to the empirical rate and reaches a deviation of about 30%.
Towards lower temperature two different behaviors can be identified. The ratios of the theoretical and empirical rates for <sup>20</sup>Ne($`\alpha ,\gamma `$)<sup>24</sup>Mg and <sup>40</sup>Ca($`\alpha ,\gamma `$)<sup>44</sup>Ti increase towards lower temperatures, whereas the other ratios first decrease and finally strongly increase at the lowest temperatures $`T0.2`$ GK (not shown in Fig. 7 as the corresponding rates are already very low and outside the temperature region of interest). In the latter very low temperature range the reaction rates are dominated by low energy resonances, E$`{}_{R}{}^{}`$1 MeV. For these resonances the $`\alpha `$ width $`\mathrm{\Gamma }_\alpha `$ is typically smaller than the $`\gamma `$ width and determines the resonance strength $`\omega \gamma _{(\alpha ,\gamma )}`$. The reduction of $`\gamma `$-strength has no influence for these states. With only a few dominating resonances the assumptions of the statistical model are not valid anymore.
For the decrease of the ratios in the temperature range $`T<3`$ GK the proper energy dependence of the optical $`\alpha `$+nucleus potential is an important factor. The slope of the ratio plotted in Fig. 7 for each reaction is sensitive to the choice of the optical potential. The use of an equivalent square well potential which neglects absorption in the Coulomb barrier leads to a slightly less steep slope. Apparently, all available $`\alpha `$+nucleus potentials in statistical model calculations cannot account for the proper energy dependence at energies close to the Coulomb barrier. This is a well-known problem in such calculations.
The cases for <sup>20</sup>Ne($`\alpha ,\gamma `$)<sup>24</sup>Mg and <sup>40</sup>Ca($`\alpha ,\gamma `$)<sup>44</sup>Ti seem different, there the HF predictions are larger than the empirical rate and the ratios increase monotonically towards lower temperatures. However, the HF calculation is in agreement within a factor of $``$2 for the temperature range T$``$0.4 GK with the experimental value directly derived from $`\alpha `$ capture resonances in the reaction <sup>20</sup>Ne($`\alpha ,\gamma `$). For <sup>40</sup>Ca($`\alpha ,\gamma `$), the empirical rate is nearly identical with the experimental rate for temperatures T=1-5 GK as can be seen from figure 6. Most of the available experimental information about natural parity states in <sup>44</sup>Ti is based on $`\alpha `$ capture studies . Only two $`\alpha `$ unbound states at lower energies, E$`{}_{x}{}^{}`$7.5 MeV, are known from <sup>46</sup>Ti(p,t)<sup>44</sup>Ti two-particle transfer measurements . Also the results of <sup>40</sup>Ca(<sup>6</sup>Li,d)<sup>44</sup>Ti and <sup>40</sup>Ca(<sup>7</sup>Li,t)<sup>44</sup>Ti $`\alpha `$-transfer measurements indicate a fairly low level density in the excitation range of interest. Shell model calculations support the experimental results that the level density in the pf-shell nucleus <sup>44</sup>Ti is low . The low level density in <sup>44</sup>Ti may limit the applicability of the HF approach for calculating the reaction rate of <sup>40</sup>Ca($`\alpha ,\gamma `$)<sup>44</sup>Ti. This is supported by the fact that a variation in the level density description employed in the statistical model calculation has a significant impact on the reaction rate of <sup>40</sup>Ca($`\alpha ,\gamma `$)<sup>44</sup>Ti but leaves the remaining rates nearly unchanged. This high sensitivity is due to the low level density in the compound nucleus.
## IV Summary and Conclusion
In this paper we attempted to improve the statistical model description of $`\alpha `$ capture rates on self-conjugate nuclei, which are important for many nucleosynthesis processes in stellar and explosive He-burning.
The NON-SMOKER predictions presented here are typically lower than the results of previous Hauser Feshbach calculations , except for <sup>28</sup>Si($`\alpha `$,$`\gamma `$)<sup>32</sup>S for which the present theoretical rate is higher by a factor of 1.96 (see also Table 6 in ). Those previous calculations approximated the isospin effect by simply dividing the total $`\gamma `$ width by a factor of 5 for isospin conjugated nuclei and employed equivalent square well potentials for the calculation of the $`\alpha `$ transmission coefficients. As expected, differences are larger in comparison to calculations neglecting the isospin suppression.
The statistical model rates are compared with reaction rates derived directly from experiment or calculated from resonance parameters which have been measured through different reaction channels. These rates agree reasonably well with the HF predictions in the astrophysically interesting temperature range $`1<T_9<5`$ (see Fig. 7). At lower temperatures, the differences tend to be larger, due to the increasing level spacing and the importance of single resonances.
As the isobaric energy shifts are quite different for the considered compound nuclei but the general trend in the isospin suppression of the $`\gamma `$ widths is nevertheless well reproduced in the theoretical framework, we conclude that the approach presented here is valid in the energy range of astrophysical importance. The remaining differences have to be attributed to the description of the energy dependence of other nuclear properties, such as the optical $`\alpha `$+nucleus potential or the level density.
## ACKNOWLEDGMENTS
This work was partly supported by the Swiss NSF (grant 2000-053798.98) and by the National Science Foundation (grant PHY98-03757). T. R. is a PROFIL fellow of the Swiss NSF (grant 2124-055832.98).
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# Thermodynamics of spin 𝑺=𝟏/𝟐 antiferromagnetic uniform and alternating-exchange Heisenberg chains
## I Introduction
An antiferromagnetic alternating-exchange Heisenberg chain is one in which nearest-neighbor spins in the chain interact via a Heisenberg interaction, but with two antiferromagnetic (AF) exchange constants $`J_2J_1,J_1,J_20`$ which alternate from bond to bond along the chain; the alternation parameter is $`\alpha J_2/J_1`$. Here we will be concerned with the magnetic susceptibility $`\chi `$ and specific heat $`C`$ versus temperature $`T`$ of alternating-exchange chains consisting of spins $`S=1/2`$. The uniform AF Heisenberg chain is one limit of the alternating chain in which the two exchange constants are equal ($`\alpha =1,J_1=J_2J`$). At the other limit is the isolated dimer in which one of the exchange constants is zero ($`\alpha =0`$). The present work is a combined theoretical and experimental study of $`\chi (T)`$ and $`C(T)`$ of the $`S=1/2`$ AF alternating-exchange chain over the entire range $`0\alpha 1`$ of the alternation parameter, with the emphasis on the regime $`\alpha 1`$ at and close to the uniform chain limit. This latter regime is relevant for compounds showing second order spin dimerization transitions with decreasing $`T`$. The present work was originally motivated by our desire to accurately extract the temperature dependent energy gap $`\mathrm{\Delta }(T)`$ for magnetic excitations, the “spin gap”, from experimental $`\chi (T)`$ data for the $`S=1/2`$ chain/two-leg ladder compound $`\mathrm{NaV}_2\mathrm{O}_5`$ below its spin dimerization temperature $`T_\mathrm{c}34`$ K. We found that existing theory for the alternating-exchange chain was insufficient to accomplish this goal. In the present work we critically examine the predictions of previous theory, perform the required additional theoretical calculations, and then apply the results to extract $`\mathrm{\Delta }(T)`$ at $`TT_\mathrm{c}`$ from our $`\chi (T)`$ data for $`\mathrm{NaV}_2\mathrm{O}_5`$ single crystals. We have extended the original goal so that we also include theoretical and experimental studies of $`C(T)`$ and how this quantity relates to $`\chi (T)`$. In the remainder of this introduction we briefly review the prior theoretical results pertaining to $`\chi (T)`$ and $`C(T)`$ of the uniform and alternating-exchange chain to place our work in the proper context. We then review the experimental and theoretical background on $`\mathrm{NaV}_2\mathrm{O}_5`$ and describe the plan for the rest of the paper.
### A Theory
The $`\chi (T)`$ and $`C(T)`$ of both limits of the $`S=1/2`$ AF alternating-exchange Heisenberg chain are known exactly. For the dimer, the $`\chi (T)`$ is given by the exact Eq. (14) below and the exact $`C(T)`$ is also easily calculated. The $`\chi (T)`$ and $`C(T)`$ of the uniform chain for $`T0.4J/k_\mathrm{B}`$ ($`k_\mathrm{B}`$ is Boltzmann’s constant) were estimated from calculations for chains with $`11`$ spins by Bonner and Fisher in 1964; they extended their results by extrapolating to $`T=0`$, and in the case of $`\chi (T)`$ to the exact $`T=0`$ value. The exact solution for $`\chi (T)`$ of the uniform chain was obtained using the Bethe ansatz in 1994 by Eggert, Affleck and Takahashi, and compared with their low-$`T`$ results from conformal field theory. They found, remarkably, that $`\chi (T0)`$ has infinite slope. Their numerical $`\chi (T)`$ values are up to $`10`$% larger than the Bonner-Fisher extrapolation for $`T0.25J/k_\mathrm{B}`$ (for a comparison of the two predictions, see Fig. 8.1 in Ref. ). Their conformal field theory calculations showed that the leading order correction to the zero temperature limit is of the form $`\chi (T)=\chi (0)\{1+1/[2\mathrm{ln}(T_0/T)]\}`$, where the value of the scaling temperature $`T_0`$ is not predicted by the field theory. Such log terms are called “logarithmic corrections” in the literature. One of us recently presented numerical Bethe ansatz calculations of $`\chi (T)`$ with a higher absolute accuracy for $`\chi (T)`$ estimated to be $`1\times 10^7`$, and showed that the data are consistent with the above field theory prediction, with an additional higher order logarithmic correction, over the temperature range $`5\times 10^{25}k_\mathrm{B}T/J10^3`$. Corresponding $`C(T)`$ calculations were also carried out, and logarithmic corrections were studied for this quantity as well. Lukyanov has recently presented an exact theory for $`\chi (T)`$ and $`C(T)`$ at low $`T`$, including the exact value of $`T_0`$. In the present work, we compare the very recent numerical Bethe ansatz results of Klümper and Johnston with the predictions of Lukyanov’s theory and find agreement for $`\chi (T)`$ to high accuracy ($`1\times 10^6`$) over a temperature range spanning 18 orders of magnitude, $`5\times 10^{25}k_\mathrm{B}T/J5\times 10^7`$; the agreement in the lower part of this temperatures range is much better, $`𝒪(10^7)`$. For $`C(T)`$, the logarithmic correction in Lukyanov’s theory is insufficient to describe the Bethe ansatz data sufficiently accurately even at very low temperatures, so we derive the next two logarithmic corrections from the Bethe ansatz $`C(T)`$ data. For various applications, it would be desirable to have fits to the $`\chi (T)`$ and $`C(T)`$ Bethe ansatz data which extend to higher temperatures. We describe the formulation and implementation of fit functions, incorporating the influence of the logarithmic correction terms, which yield extremely precise fits to the data for both quantities over the entire 25 decades in temperature of the calculations, $`5\times 10^{25}k_\mathrm{B}T/J5`$.
The $`\chi (T)`$ in the intermediate regime $`0<\alpha <1`$ has been investigated analytically in the Hartree-Fock approximation and using numerical techniques. Of particular interest here is the regime $`\alpha 1`$, close to the uniform limit, which is the regime relevant to materials exhibiting a dimerization transition with decreasing $`T`$ such as occurs in materials exhibiting a spin-Peierls transition. There are no accurate theoretical predictions available for $`\chi (T)`$ of the alternating-exchange Heisenberg chain in this regime, which is the property usually used to initially characterize the occurrence of such a transition experimentally. To address this deficiency and to also cover a more extended $`\alpha `$ range, we carried out extensive quantum Monte Carlo (QMC) simulations and transfer-matrix density-matrix renormalization group (TMRG) calculations of $`\chi (T)`$ for $`0.05\alpha 1`$ over the temperature range $`0.002k_\mathrm{B}T/J_110`$.
An interesting issue is how the spin gap $`\mathrm{\Delta }`$ evolves with alternation parameter $`\alpha `$ as the uniform limit is approached, $`\alpha 1`$. Because the uniform chain is a gapless quantum-critical system, the introduction of alternating exchange along the chain has been theoretically predicted to yield a nonanalytic $`\mathrm{\Delta }(\alpha )`$ behavior for $`\alpha 1`$. We derive $`\mathrm{\Delta }(\alpha )`$ by fitting our low-$`t`$ TMRG $`\chi (T)`$ data by an expression which we formulated. The $`\mathrm{\Delta }(\alpha )`$ results are compared with those of previous numerical calculations and with the theoretical prediction. We infer from our data that the asymptotic critical regime is only entered for $`\alpha 0.99`$.
In order to be optimally useful for accurately modeling experimental $`\chi (T)`$ data for alternating-exchange chain compounds, our QMC and TMRG $`\chi (\alpha ,T)`$ results must first be accurately fitted by a continuous function of both $`\alpha `$ and $`T`$. We will introduce a general fit function which eventually proves capable of fitting these combined data for the alternating-exchange Heisenberg chain very accurately. We first fit the $`\chi (T)`$ of the uniform chain and isolated dimer using this function and then use the obtained fitting parameters as end-point parameters in the fit to our combined QMC and TMRG data for intermediate values of $`\alpha `$. The final fit function is a single two-dimensional function of $`\alpha `$ and $`T`$ for $`0\alpha 1`$ which can be used to extract the (possibly temperature-dependent) alternation parameter, exchange constants and spin gap from experimental $`\chi (T)`$ data for compounds for which the $`S=1/2`$ AF alternating-exchange Heisenberg chain Hamiltonian is appropriate. Our fit function will also be useful as a reference for $`\chi (T)`$ calculated from other related $`S=1/2`$ Hamiltonians such as that incorporating the spin-phonon interaction for spin-Peierls systems.
### B NaV<sub>2</sub>O<sub>5</sub>
Vanadium oxides show a remarkable variety of electronic behaviors. For example, the metallic fcc normal-spinel structure compound LiV<sub>2</sub>O<sub>4</sub> shows local momentlike behaviors above $`50`$ K, crossing over to heavy fermion behaviors below $`10`$ K. On the other hand, the $`d^1`$ compound CaV<sub>2</sub>O<sub>5</sub> has a two-leg trellis-ladder-layer structure in which all of the V atoms are crystallographically equivalent and is a Mott-Hubbard insulator. The $`\chi (T)`$ shows a spin-gap $`\mathrm{\Delta }/k_\mathrm{B}660`$ K arising from strong coupling of the V $`S=1/2`$ spins across a rung. Modeling of $`\chi (T)`$ by QMC simulations confirmed that this compound consists magnetically of V<sub>2</sub> dimers, with an intradimer AF exchange constant $`J/k_\mathrm{B}665`$ K and with very weak interdimer interactions.
The compound NaV<sub>2</sub>O<sub>5</sub> can also be formed. The crystal structure was initially reported in 1975 to consist of two-leg ladders as in CaV<sub>2</sub>O<sub>5</sub>, but in a non-centrosymmetric (acentric) structure (space group P2<sub>1</sub>mn) in which charge segregation occurs such that one leg of each ladder consists of V<sup>+4</sup> and the other of crystallographically inequivalent V<sup>+5</sup> ions. However, recently five different crystal structure investigations showed that the structure is actually centrosymmetric (space group Pmmn), with all V atoms crystallographically equivalent at room temperature, so that (static) charge segregation between the V atoms does not, in fact, occur. This result is consistent with <sup>51</sup>V NMR investigations which showed the presence of only one type of V atom at room temperature. This compound is thus formally a mixed-valent $`d^{0.5}`$ system, which has been considered in a one-electon-band picture to be a quarter-filled ladder compound. We note that from modeling optical excitations in the energy range 4 meV–4 eV, Damascelli and coworkers initially concluded that the room-temperature structure of NaV<sub>2</sub>O<sub>5</sub> is acentric; their analysis was consistent with the V atoms on a rung of a ladder having oxidation states of 4.1 and 4.9, respectively. However, this group subsequently explained that length- and/or time-scale-of-measurement issues may be involved in their interpretation, such that charge disproportionation between V atoms may only occur locally and possibly dynamically, which could then be consistent with the (average long-range) crystal structure refinements and NMR measurements. Theoretical support for this scenario was provided by Nishimoto and Ohta. Factor group analyses of the possible IR- and Raman-active phonon modes and comparisons with experimental observations at room temperature are consistent with the centrosymmetric space group for the compound. A first-principles electronic structure study based on the density functional method within the generalized gradient approximation showed that the total energy of the centric structure is about 1.0 eV/(formula unit) lower than that of the acentric structure, consistent with the recent structural studies.
One might expect that the hole-doping which occurs upon replacing Ca in CaV<sub>2</sub>O<sub>5</sub> by Na would result in metallic properties for NaV<sub>2</sub>O<sub>5</sub>, because of the nonintegral oxidation state of the V cations and of the crystallographic equivalence of these atoms. However, NaV<sub>2</sub>O<sub>5</sub> is a semiconductor. This has been explained by the formation of $`d^1`$ V-O-V molecular clusters on the rungs of the two-leg ladders, again resulting in a Mott-Hubbard insulator due to the on-site Coulomb repulsion, where in this case a “site” is a V-O-V molecular cluster. Thus a nonintegral oxidation state and crystallographic equivalence of transition metal atoms in a compound are not sufficient to guarantee metallic character simply by symmetry; all nearest-neighbor pairs, triplets, …, of transition metal atoms must also be crystallographically equivalent, which is not the case in NaV<sub>2</sub>O<sub>5</sub>, since a V<sub>2</sub> pair on a rung is not crystallographically equivalent to one on a leg in the two-leg ladders. In contrast, all V atoms and pairs of V atoms in mixed-valent fcc LiV<sub>2</sub>O<sub>4</sub> are respectively crystallographically equivalent, resulting in metallic character as demanded by symmetry.
The V-O-V rung molecular clusters which are coupled along the ladder direction in NaV<sub>2</sub>O<sub>5</sub> may be considered to form an effective $`S=1/2`$ one-dimensional (1D) chain. Experimental support for this picture, often quoted in the literature, is that the magnetic susceptibility (above $`T_\mathrm{c}`$, see below) is in agreement with the Bonner-Fisher prediction for the $`S=1/2`$ Heisenberg chain, as reported by Isobe and Ueda. Angle-resolved photoemission spectroscopy (ARPES) measurements on NaV<sub>2</sub>O<sub>5</sub> by Kobayashi et al. showed that the electronic structure is essentially 1D, despite the ostensibly 2D nature of the trellis layer, with dispersion in the oxygen and copper bands (below the Fermi energy) occurring only in the ladder direction ($`b`$-axis). Interestingly, the dispersion in the lowest binding energy part of the occupied Cu Hubbard band showed a lattice periodicity of 2$`b`$, which may reflect dynamical short-range AF and/or crystallographic ordering in the ladder direction. Temperature-dependent ARPES measurements on Na<sub>0.96</sub>V<sub>2</sub>O<sub>5</sub> by the same group from 120 to 300 K showed evidence for the predicted spin-charge separation in 1D magnetic systems.
A phase transition occurs in NaV<sub>2</sub>O<sub>5</sub> at a critical temperature $`T_\mathrm{c}33`$–36 K, below which the spin susceptibility $`\chi ^{\mathrm{spin}}0`$ as $`T0`$ and a lattice distortion occurs. The lattice distortion results in a supercell with lattice parameters $`2a\times 2b\times 4c`$. Therefore the transition was initially characterized as a possible spin-Peierls transition, which by definition is driven by magnetoelastic (spin-phonon) coupling, and in which there is no change in the charge/spin distribution within the rungs/V-O-V molecular clusters. The superstructure in the $`a`$ and $`c`$ directions, perpendicular to the V chains which run in the $`b`$ direction, would be a result of the phasing of the distortions in adjacent chains/ladders. In this interpretation, and within the adiabatic approximation (discussed later), one would expect that the magnetic properties above $`T_\mathrm{c}`$ should be close to those of the $`S=1/2`$ Heisenberg uniform chain, and of an $`S=1/2`$ alternating-exchange Heisenberg chain below $`T_\mathrm{c}`$.
It has become clear, however, that the phase transition occurring at $`T_\mathrm{c}`$ in NaV<sub>2</sub>O<sub>5</sub> is accompanied by charge ordering, in contrast to a classic spin-Peierls transition. Therefore, the magnetoelastic coupling may only play a secondary role, and the spin gap may be a secondary order parameter. In particular, <sup>51</sup>V NMR experiments showed the presence of (inequivalent) V<sup>+4</sup> and V<sup>+5</sup> below $`T_\mathrm{c}`$, whereas only one V species was present above $`T_\mathrm{c}`$. This result is consistent with the solution of the superstructure below $`T_\mathrm{c}`$ by Lüdecke and co-workers using synchrotron x-ray diffraction. Lüdecke et al. found that there are modulated and unmodulated chains of V atoms below $`T_\mathrm{c}`$, tentatively assigned to magnetic and nonmagnetic chains. One interpretation of the results is that the $`d^1`$ V<sup>+4</sup> cations segregate into alternate two-leg ladders which are isolated from each other within the $`\mathrm{V}_2\mathrm{O}_3`$ trellis layer by intervening two-leg ladders containing only nonmagnetic V<sup>+5</sup>. The anomalous strong increase in the thermal conductivity below $`T_\mathrm{c}`$ may also be due to charge ordering. From ultrasonic measurements of shear and longitudinal elastic constants, Schwenk and co-workers have suggested that the charge ordering is of the zig-zag type within each ladder. In each of these scenarios for charge ordering, static charge disproportionation occurs such that 1/2 of the V atoms have oxidation state +4 and the other half +5, consistent with the average formal oxidation state of +4.5 in the compound.
Köppen et al. have concluded from thermal expansion measurements that the phase transition at $`T_\mathrm{c}`$ actually consists intrinsically of two closely-spaced phase transitions separated by $`1`$ K, where the upper transition is thermodynamically of second order whereas the lower one is first order. However, a double transition was not found in their specific heat measurements on the same crystal, which they attributed to the 50 mK temperature oscillations required by their ac measurement technique which were thought to broaden the two transitions and render them indistinguishable.
The nature of the possible charge ordering pattern has been studied theoretically by several groups. Seo and Fukuyama predicted (at $`T=0`$) a static zig-zag chain of V<sup>+4</sup> ions on each two-leg ladder, with an interpenetrating zig-zag chain of V<sup>+5</sup> ions. They proposed that pairs of V<sup>+4</sup> spins, one each on adjacent ladders, would form spin singlets, resulting in the observed spin gap. A similar zig-zag charge configuration in each ladder was inferred by Mostovoy and Khomskii, with subsequent experimental support by Smirnov et al., and by Gros and Valenti. Motivated in part by the above thermal expansion measurement results of Köppen et al., Thalmeier and Fulde proposed that the charge ordering transition would result in one linear chain of V<sup>+4</sup> and one linear chain of V<sup>+5</sup> on each two-leg ladder, thereby then allowing a conventional spin-Peierls transition to occur at a slightly lower temperature, resulting in a double transition as reported by Köppen et al. A similar picture was put forward by Nishimoto and Ohta. Thalmeier and Yaresko have extensively discussed the linear-chain and zig-zag scenarios for charge ordering, and in addition have considered the alternating two-leg ladder charge ordering pattern of the type suggested by Lüdecke et al. They point out that in both the linear and zig-zag patterns, a secondary spin-Peierls dimerization or spin exchange anisotropy (in spin space) may be necessary to give a spin gap, whereas the two-leg ladder ordering has a spin gap even with no lattice distortion. Thalmeier and Yaresko describe the characteristic signatures of each of the charge-ordered models to be compared with experimental inelastic neutron scattering measurements. Riera and Poilblanc have discussed the influence of electron-phonon coupling on the derived charge- and spin-order phase diagrams.
We have carried out $`\chi (T)`$ measurements from 2 to 750 K on single crystals of NaV<sub>2</sub>O<sub>5</sub> along the ladder ($`b`$ axis) direction to further characterize and clarify the nature of the magnetic interactions and ordering below and above $`T_\mathrm{c}`$. We find that the magnetic properties above $`T_\mathrm{c}`$ are not accurately described by the $`S=1/2`$ Heisenberg uniform chain prediction with a $`T`$-independent $`J`$, although a mean-field ferromagnetic interchain coupling can explain these data. Using our theoretical $`\chi (\alpha ,T)`$ fit function for the AF alternating-exchange chain below $`T_\mathrm{c}`$, we find that $`\delta (0)(1\alpha )/(1+\alpha )=0.034(6)`$ and that the zero-temperature spin-gap of NaV<sub>2</sub>O<sub>5</sub> is $`\mathrm{\Delta }(0)/k_\mathrm{B}=103(2)`$ K. The $`\delta (T)`$ and $`\mathrm{\Delta }(T)`$ below $`T_\mathrm{c}`$ are extracted. A spin pseudogap is found to occur above $`T_\mathrm{c}`$ with a rather large magnitude. From our specific heat measurements on two crystals, we find that the magnetic specific heat at low temperatures $`T15`$ K is too small to be resolved experimentally, and that the spin entropy at $`T_\mathrm{c}`$ is too small to account for the entropy of the transition. A quantitative analysis shows that at least 77 % of the entropy change at $`T_\mathrm{c}`$ due to the transition(s) and associated order parameter fluctuations must arise from the lattice and/or charge degrees of freedom and less than 23 % from the spin degrees of freedom.
### C Plan of the paper
The rest of the paper is organized as follows. Our notation for the Heisenberg spin Hamiltonian and for the reduced susceptibility, temperature and spin gap are given immediately in Sec. II. Some general features of the high-temperature series expansion (HTSE) for $`\chi (T)`$ and $`C(T)`$ of $`S=1/2`$ Heisenberg spin lattices and the low-temperature limits of these quantities for one-dimensional (1D) systems with a spin gap are then given. We then specialize to the $`S=1/2`$ AF alternating-exchange Heisenberg chain in Sec. II C, where we discuss the HTSEs, the spin gap and the one-magnon dispersion relations $`E(\mathrm{\Delta },k)`$. In the latter subsection, we derive a one-parameter approximation for $`E(\mathrm{\Delta },k)`$ which correctly extrapolates to the $`\alpha 0`$ limit and which we will need in order to later fit the TMRG $`\chi (T)`$ data to extract $`\mathrm{\Delta }(\alpha )`$. We also show that the expressions for the low-$`T`$ limits of both $`\chi (T)`$ and $`C(T)`$ depend only on the spin gap (in addition to $`T`$). In Sec. III, we discuss overall features of the $`\chi (T)`$ and $`C(T)`$ for the uniform chain and then focus on the low-$`T`$ behavior. The explicit forms of the logarithmic corrections previously found for $`\chi (T)`$ are first discussed. We show that a low-$`T`$ expansion of the theory of Lukyanov gives the same first two corrections, and in addition gives the next higher-order term. We then compare the Bethe ansatz $`\chi (T)`$ results directly with the theory with no adjustable parameters or approximations. Logarithmic corrections are also found to be important to accurately describe the Bethe ansatz data for $`C(T)`$. We show that the lowest order correction is not sufficient to fit the data, and we derive the next two higher-order corrections by fitting the data at very low temperatures.
General features of our scheme to fit numerical $`\chi (T)`$ data are described in Sec. IV A, followed by a fit to the exact $`\chi (T)`$ for the antiferromagnetic Heisenberg dimer and two fits to the numerical $`\chi (T)`$ data for the uniform chain. Due to the special requirements of, and constraints on, the two-dimensional fit function necessary to accurately fit $`\chi (\alpha ,T)`$ data for the alternating-exchange chain over large ranges of both $`\alpha `$ and $`T`$, a separate section, Sec. IV E, is devoted to formulating and discussing this fit function. Using a fit function similar to that used to fit numerical $`\chi (T)`$ data, in the next section an extremely accurate and precise fit is obtained over 25 decades in temperature to the Bethe ansatz $`C(T)`$ data. Our QMC and TMRG $`\chi (T)`$ data for the alternating-exchange chain are presented and fitted in Sec. V, using as end-point parameters those determined for the uniform chain and the dimer, respectively. The spin gap $`\mathrm{\Delta }(\alpha )`$ is extracted for $`0.8\alpha 0.995`$ by fitting the TMRG $`\chi (\alpha ,T)`$ data at low temperatures in Sec. VI. Section VII contains a comparison of our numerical results with previous work. The $`\mathrm{\Delta }(\alpha )`$ values are compared with previous numerical results and with the theoretical prediction for the asymptotic critical behavior in Sec. VII A. Our $`\chi (T)`$ calculations are shown in Sec. VII B to be in good agreement with the previous numerical results of Barnes and Riera for $`0.2\alpha 0.8`$. The numerical calculations of Bulaevskii have been extensively used in the past by experimentalists to fit the $`\chi (T)`$ of spin-Peierls compounds, but up to now a detailed analysis of the predictions of this theory has not been given. We present such an analysis in Sec. VII C and compare our results with these predictions.
We begin the experimental part of the paper by studying the anisotropic magnetic susceptibility of a high quality NaV<sub>2</sub>O<sub>5</sub> single crystal in Sec. VIII A, where literature data on the anisotropy of the $`g`$ factor and Van Vleck susceptibility are compared with our results. In the following sections we illustrate the utility and application of many of the theoretical results derived and presented previously in the paper. In Sec. VIII B we present experimental $`\chi (T)`$ data for single crystals of NaV<sub>2</sub>O<sub>5</sub> and model these data in detail in Sec. VIII C using our QMC and TMRG $`\chi (T)`$ data fit function for the AF alternating-exchange Heisenberg chain. We show that qualitatively and quantitatively new information about the temperature dependences of the spin dimerization parameter and spin gap below $`T_\mathrm{c}`$ can be obtained from our modeling. This analysis also shows that spin dimerization fluctuations and a spin pseudogap are present above $`T_\mathrm{c}`$, and we quantitatively determine their magnitudes. Our specific heat measurements of NaV<sub>2</sub>O<sub>5</sub> single crystals and our extensive modeling of these data are presented in Sec. VIII D, where we obtain quantitative limits on the relative contributions of the lattice, spin and charge degrees of freedom to the change in the entropy due to the transition at $`T_\mathrm{c}`$ and to associated order parameter fluctuations. A summary and concluding discussion are given in Sec. IX.
## II Theory
In this paper we will only be concerned with the spin $`S=1/2`$ antiferromagnetic (AF) Heisenberg Hamiltonian
$$=\underset{<ij>}{}J_{ij}𝑺_i𝑺_j,$$
(1)
where $`J_{ij}`$ is the Heisenberg exchange interaction between spins $`𝑺_i`$ and $`𝑺_j`$ and the sum is over unique exchange bonds. A $`J_{ij}>0`$ corresponds to AF coupling, whereas $`J_{ij}<0`$ refers to ferromagnetic coupling. Note that magnetic nearest neighbors $`𝑺_j`$ of a given spin $`𝑺_i`$ in Eq. (1) need not be crystallographic nearest neighbors. A magnetic nearest neighbor of a given spin is any other spin with which the given spin has an exchange interaction.
For notational convenience, we define the reduced spin susceptibilities $`\chi ^{}`$ and $`\overline{\chi ^{}}`$, reduced temperatures $`t`$ and $`\overline{t}`$ and reduced spin gaps $`\mathrm{\Delta }^{}`$ and $`\overline{\mathrm{\Delta }^{}}`$ as
$$\chi ^{}\frac{\chi J^{\mathrm{max}}}{Ng^2\mu _\mathrm{B}^2},\overline{\chi ^{}}\frac{\chi J}{Ng^2\mu _\mathrm{B}^2},$$
(2)
$$t\frac{k_\mathrm{B}T}{J^{\mathrm{max}}},\overline{t}\frac{k_\mathrm{B}T}{J},$$
(3)
$$\mathrm{\Delta }^{}\frac{\mathrm{\Delta }}{J^{\mathrm{max}}},\overline{\mathrm{\Delta }^{}}\frac{\mathrm{\Delta }}{J},$$
(4)
where $`J^{\mathrm{max}}`$ and $`J`$ are, respectively, the largest and average exchange constants in the system, $`N`$ is the number of spins, $`g`$ is the spectroscopic splitting factor appropriate to the direction of the applied magnetic field relative to the crystallographic axes, and $`\mu _\mathrm{B}`$ is the Bohr magneton.
### A High-temperature series expansions for the spin susceptibility and magnetic specific heat
For any Heisenberg spin lattice (in any dimension) in which the spins are magnetically equivalent, i.e. where each spin has identical magnetic coordination spheres, the first three to four terms of the exact quantum mechanical high temperature series expansion of $`\chi ^{}(t)`$ have the same form, with a particularly simple form if the series is inverted. For $`S=1/2`$, one obtains
$$\frac{1}{4\chi ^{}t}=\underset{n=0}{\overset{\mathrm{}}{}}\frac{d_n}{t^n},$$
(6)
$$d_0=1,d_1=\frac{1}{4J^{\mathrm{max}}}\underset{j}{}J_{ij},d_2=\frac{1}{8J_{}^{\mathrm{max}}{}_{}{}^{2}}\underset{j}{}J_{ij}^2,$$
(7)
$$d_3=\frac{1}{24J_{}^{\mathrm{max}}{}_{}{}^{3}}\underset{j}{}J_{ij}^3.$$
(8)
Equation (7) is universal, but Eq. (8) holds only for spin lattices which are not geometrically frustrated for AF ordering and in which the magnetic and crystallographic nearest neighbors of a given spin are the same. Geometrically frustrated lattices typically contain closed triangular exchange paths within the spin lattice structure, such as in the 2D triangular lattice or in the 3D $`B`$ sublattices of the fcc $`AB_2\mathrm{O}_4`$ oxide normal-spinel and $`A_2B_2\mathrm{O}_7`$ oxide pyrochlore structures. The uniform and alternating-exchange chains considered in this paper are not geometrically frustrated, and the magnetic and crystallographic nearest neighbors of a given spin are the same. It has been found that the terms to $`𝒪(1/t^3)`$ on the right-hand-side of Eq. (6) are sufficient to quite accurately describe the susceptibilities of a variety of nonfrustrated zero-, one-, and two-dimensional $`S=1/2`$ AF Heisenberg spin lattices to surprisingly low temperatures $`t1`$. Higher order $`d_n/t^n`$ terms with $`n4`$ are dependent on the structure and dimensionality of the spin lattice. The Weiss temperature $`\theta `$ in the Curie-Weiss law $`\chi (T)=C/(T\theta )`$ is given by the universal expression $`\theta =d_1J^{\mathrm{max}}/k_\mathrm{B}`$.
Because the spin susceptibility and the magnetic contribution $`C(T)`$ to the specific heat can both be expressed, via the fluctuation-dissipation theorem and the Heisenberg Hamiltonian, respectively, in terms of the spin-spin correlation functions, there is a close relationship between these two quantities. In particular, just as there is a universal expression for the first three to four HTSE terms for $`\chi (T)`$ of a Heisenberg spin lattice as discussed above, a universal expression for the first one to two HTSE terms for $`C(T)`$ of such a spin lattice exists and is given for $`S=1/2`$ by
$$\frac{C(t)}{Nk_\mathrm{B}}=\frac{3}{32}\left[\frac{_jJ_{ij}^2}{t^2J_{}^{\mathrm{max}}{}_{}{}^{2}}+\frac{_jJ_{ij}^3}{2t^3J_{}^{\mathrm{max}}{}_{}{}^{3}}+𝒪\left(\frac{1}{t^4}\right)\right].$$
(9)
The sums are over all magnetic nearest-neighbor bonds of any given spin $`𝑺_i`$. The first term is universal but the second term holds only for geometrically nonfrustrated spin lattices in which the crystallographic and magnetic nearest-neighbors of any given spin are the same. Higher order terms all depend on the structure and dimensionality of the spin lattice.
A common misconception is that $`C=0`$ if the magnetic susceptibility of a local-moment system obeys the Curie-Weiss law. This is only true classically. For Heisenberg spin lattices, one can easily show that the Weiss temperature $`\theta `$ in the Curie-Weiss law arises from the first HTSE term \[$`𝒪(1/t)`$\] of the magnetic nearest-neighbor spin-spin correlation function, which is the same quantity that the first HTSE term of $`C(t)`$ arises from. Thus, e.g., for $`S=1/2`$ Heisenberg spin lattices at temperatures $`t1`$ at which the Curie-Weiss law holds, the magnetic specific heat is given by the universal first term of Eq. (9).
### B Low-temperature limit of the spin susceptibility and specific heat of 1D systems with a spin gap
##### Magnetic susceptibility.
For one-dimensional (1D) $`S=1/2`$ Heisenberg spin systems with a spin gap such as the $`S=1/2`$ two-leg ladder (and the alternating-exchange chain), Troyer, Tsunetsugu, and Würtz derived a general expression for $`\chi ^{}(t)`$ which approximately takes into account kinematic magnon interactions, given by
$$\chi ^{}(t)=\frac{1}{t}\frac{z(t)}{1+3z(t)},$$
(11)
$$z(t)=\frac{1}{\pi }_0^\pi \mathrm{e}^{\epsilon _k/t}d(ka),$$
(12)
where $`\epsilon _kE(k)/J^{\mathrm{max}}`$, $`E(k)`$ is the nondegenerate one-magnon (triplet) dispersion relation (the Zeeman degeneracy is already accounted for) and $`a`$ is the (average) distance between spins. This expression is exact in both the low- and high-temperature limits. For the isolated dimer, for which $`\epsilon _k=\mathrm{\Delta }^{}=1`$, Eq. (11) is exact at all temperatures. Inserting $`z(t)=\mathrm{e}^{1/t}`$ for the dimer into Eq. (11) yields the correct result
$$\chi ^{}(t)=\frac{\mathrm{e}^{1/t}}{t}\frac{1}{1+3\mathrm{e}^{1/t}},(\mathrm{dimer})$$
(14)
$$\chi ^{}(t0)=\frac{\mathrm{e}^{1/t}}{t}.$$
(15)
The $`\chi ^{}(t)`$ in Eq. (14) for the antiferromagnetic Heisenberg dimer is plotted in Fig. 1; the fit shown in the figure will be presented and discussed later in Sec. IV B.
At low temperatures $`t\mathrm{\Delta }^{}`$ and $`t`$ one-magnon bandwidth/$`J^{\mathrm{max}}`$, and for a dispersion relation with a parabolic dependence on wave vector $`k`$ near the band minimum
$$\epsilon _k\frac{E(k)}{J^{\mathrm{max}}}=\mathrm{\Delta }^{}+c^{}(ka)^2,$$
(16)
one can replace $`\epsilon _k`$ in Eq. (12) by the approximation (16) and replace the upper limit of the integral in Eq. (12) by $`\mathrm{}`$, yielding $`z(t)=\mathrm{e}^{\mathrm{\Delta }^{}/t}\sqrt{t}/(2\sqrt{\pi c^{}})`$. Substituting this result into Eq. (11) gives the low-$`t`$ limit
$$\chi ^{}(t0)=\frac{A}{t^\gamma }\mathrm{e}^{\mathrm{\Delta }^{}/t},$$
(18)
$$A=\frac{1}{2\sqrt{\pi c^{}}},\gamma =\frac{1}{2}.$$
(19)
This result is correct for any 1D $`S=1/2`$ Heisenberg spin system with a spin gap and with a nondegenerate (excluding Zeeman degeneracy) lowest-lying excited triplet magnon band which is parabolic at the band minimum.
On the other hand, the low-temperature limit of $`\chi ^{}(t)`$ for the isolated dimer in Eq. (15) is of the same form as Eq. (18), but with $`\gamma =1`$. Thus, for 1D systems consisting essentially of dimers which are weakly coupled to each other, a crossover from $`\gamma =1`$ to $`\gamma =1/2`$ is expected with decreasing $`t`$.
The parameters $`A`$ and $`\gamma `$ can be determined if very accurate $`\chi ^{}(t)`$ and $`\mathrm{\Delta }^{}`$ data are available. Taking the logarithm of Eq. (18) yields the low-$`t`$ prediction
$$\mathrm{ln}[\chi ^{}(t)]+\frac{\mathrm{\Delta }^{}}{t}=\mathrm{ln}A\gamma \mathrm{ln}t,$$
(21)
so plotting the left-hand-side vs $`\mathrm{ln}t`$ allows these two parameters to be determined. Alternatively, assuming $`\gamma =1/2`$, one can obtain estimates of $`A`$ and $`\mathrm{\Delta }^{}`$ using Eq. (18), according to
$$\mathrm{ln}(\chi ^{}\sqrt{t})=\mathrm{ln}A+\frac{\mathrm{\Delta }^{}}{t}$$
(22)
and/or
$$\frac{\mathrm{ln}(\chi ^{}\sqrt{t})}{(1/t)}=\mathrm{\Delta }^{}.$$
(23)
##### Specific heat.
The low-$`t`$ limit of the magnetic contribution $`C(T)`$ to the specific heat for the same model is calculated to be
$`{\displaystyle \frac{C(t0)}{Nk_\mathrm{B}}}={\displaystyle \frac{3}{2}}`$ $`\left({\displaystyle \frac{\mathrm{\Delta }^{}}{\pi c^{}}}\right)^{1/2}\left({\displaystyle \frac{\mathrm{\Delta }^{}}{t}}\right)^{3/2}`$ (25)
$`\times \left[1+{\displaystyle \frac{t}{\mathrm{\Delta }^{}}}+{\displaystyle \frac{3}{4}}\left({\displaystyle \frac{t}{\mathrm{\Delta }^{}}}\right)^2\right]\mathrm{e}^{\mathrm{\Delta }^{}/t}.`$
Note that, in addition to the ratio $`t/\mathrm{\Delta }^{}=k_\mathrm{B}T/\mathrm{\Delta }`$ of the thermal energy to the spin gap, the magnitude of $`\chi ^{}`$ in Eqs. (16) is determined by the actual value of the curvature $`c^{}`$ at the triplet one-magnon band minimum, whereas the magnitude of $`C`$ in Eq. (25) depends only on the ratio of $`c^{}`$ to $`\mathrm{\Delta }^{}`$. These formulas have been applied in the literature to model experimental data for alternating-exchange chain and two-leg spin ladder compounds. However, with one exception to our knowledge, these modeling studies have not recognized that the prefactor parameter and the spin gap are not independently adjustable parameters. For a given spin lattice, they are in fact uniquely related to each other. Their relationship for the $`S=1/2`$ two-leg Heisenberg ladder was studied in Ref. . For the alternating-exchange chain, we estimate the relationship between $`c^{}`$ and $`\mathrm{\Delta }^{}`$ below in Sec. II C 3.
### C Alternating-exchange chain
The $`S=1/2`$ AF alternating-exchange Heisenberg chain Hamiltonian is written in three equivalent ways as
$``$ $`=`$ $`{\displaystyle \underset{i}{}}J_1\stackrel{}{S}_{2i1}\stackrel{}{S}_{2i}+J_2\stackrel{}{S}_{2i}\stackrel{}{S}_{2i+1}`$ (27)
$`=`$ $`{\displaystyle \underset{i}{}}J_1\stackrel{}{S}_{2i1}\stackrel{}{S}_{2i}+\alpha J_1\stackrel{}{S}_{2i}\stackrel{}{S}_{2i+1}`$ (28)
$`=`$ $`{\displaystyle \underset{i}{}}J(1+\delta )\stackrel{}{S}_{2i1}\stackrel{}{S}_{2i}+J(1\delta )\stackrel{}{S}_{2i}\stackrel{}{S}_{2i+1},`$ (29)
where
$`J_1`$ $`=`$ $`J(1+\delta )={\displaystyle \frac{2J}{1+\alpha }},`$ (31)
$`\alpha `$ $`=`$ $`{\displaystyle \frac{J_2}{J_1}}={\displaystyle \frac{1\delta }{1+\delta }},`$ (33)
$`\delta `$ $`=`$ $`{\displaystyle \frac{J_1}{J}}1={\displaystyle \frac{J_1J_2}{2J}}={\displaystyle \frac{1\alpha }{1+\alpha }},`$ (35)
$`J`$ $`=`$ $`{\displaystyle \frac{J_1+J_2}{2}}=J_1{\displaystyle \frac{1+\alpha }{2}},`$ (37)
with AF couplings $`J_1J_20`$, $`0(\alpha ,\delta )1`$. The uniform undimerized chain corresponds to $`\alpha =1,\delta =0,J_1=J_2=J`$. The form of the Hamiltonian in Eq. (29) is most appropriate for chains showing a second-order dimerization transition at $`T_\mathrm{c}`$ with decreasing $`T`$. If the exchange modulation $`\delta 1`$ ($`\alpha 1`$), the (average) $`J`$ below $`T_\mathrm{c}`$ can be identified with the exchange coupling in the high-$`T`$ undimerized state.
In spin-Peierls systems, the spin-phonon interaction causes a lattice dimerization to occur below the spin-Peierls transition temperature, together with a spin-gap due to the formation of spin singlets on the dimers. The Hamiltonian can be mapped onto the spin Hamiltonian (II C) (with renormalized exchange constants) only in the adiabatic parameter regime, in which the relevant phonon energy is much smaller than $`J`$. If such a mapping cannot be made, dynamical phonon effects (quantum mechanical fluctuations) become important and the $`\chi (T)`$ can be significantly different from that predicted from Hamiltonian (II C). This issue will be discussed further when modeling the $`\chi (T)`$ data for NaV<sub>2</sub>O<sub>5</sub> in Sec. VIII B.
#### 1 High-temperature series expansions
##### Magnetic Susceptibility.
For the alternating-exchange chain, according to our definition one has $`J^{\mathrm{max}}=J_1`$. Then using the definition for $`\alpha `$ in Eq. (33), the $`d_n`$ HTSE coefficients in Eqs. (7) and (8) become
$$d_0=1,d_1=\frac{1+\alpha }{4},d_2=\frac{1+\alpha ^2}{8},d_3=\frac{1+\alpha ^3}{24}.$$
(38)
One can change variables from $`\alpha `$ and $`J_1`$ in $`\chi ^{}(\alpha ,t)`$ to $`\delta `$ and $`J`$ in $`\overline{\chi ^{}}(\delta ,\overline{t})`$ using Eqs. (II C) which give
$$t=\frac{\overline{t}}{1+\delta },$$
(40)
$$\overline{\chi ^{}}(\delta ,\overline{t})=\frac{1}{1+\delta }\chi ^{}(\frac{1\delta }{1+\delta },\frac{\overline{t}}{1+\delta }).$$
(41)
We write the resulting HTSE for the inverse of $`\overline{\chi ^{}}(\delta ,\overline{t})`$ as
$$\frac{1}{4\overline{\chi ^{}}\overline{t}}=\underset{n=0}{\overset{\mathrm{}}{}}\frac{\overline{d}_n}{\overline{t}^n},$$
(43)
where we find
$$\overline{d}_0=1,\overline{d}_1=\frac{1}{2},\overline{d}_2=\frac{1+\delta ^2}{4},\overline{d}_3=\frac{1+3\delta ^2}{12}.$$
(44)
An important feature of this HTSE of $`\overline{\chi ^{}}(\delta ,\overline{t})`$ is that it is an even (analytic) function of $`\delta `$ for any finite temperature. This constraint must be true in general and not just for the terms listed, because $`\overline{\chi ^{}}(\delta ,\overline{t})`$ cannot depend on the sign of $`\delta =(J_1J_2)/(2J)`$: the Hamiltonian in Eq. (29) is invariant upon such a sign change. Physically, a negative $`\delta `$ would simply correspond to relabeling all $`𝑺_i𝑺_{i+1}`$, which cannot change the physical properties. We will use this constraint that $`\overline{\chi ^{}}(\delta ,\overline{t})`$ must be an even function of $`\delta `$ to help formulate our fitting function (after a change back in variables) for our QMC and TMRG $`\chi ^{}(\alpha ,t)`$ calculations for the alternating-exchange chain. This constraint is important because it allows a fit function for $`\chi ^{}(\alpha ,t)`$ to be formulated which is accurate for $`\alpha 1`$ ($`\delta 1`$), a parameter regime relevant to compounds exhibiting second-order spin-dimerization transitions with decreasing temperature.
##### Magnetic specific heat.
Using $`J^{\mathrm{max}}`$ = $`J_1`$ and $`\alpha =J_2/J_1`$, the general HTSE expression in Eq. (9) yields the two lowest-order HTSE terms for the magnetic specific heat $`C(T)`$ of the $`S=1/2`$ AF alternating-exchange Heisenberg chain as
$$\frac{C(t)}{Nk_\mathrm{B}}=\frac{3}{32}\left[\frac{1+\alpha ^2}{t^2}+\frac{1+\alpha ^3}{2t^3}+𝒪\left(\frac{1}{t^4}\right)\right].$$
(45)
#### 2 Spin gap
The spin gap $`\mathrm{\Delta }^{}(\alpha )`$ of the alternating-exchange chain was determined to high ($`1\%`$) accuracy for $`0\alpha 0.9`$, in $`\alpha `$ increments of 0.1, using multiprecision methods by Barnes, Riera, and Tennant (BRT). They found that their calculations could be parametrized well by
$$\mathrm{\Delta }^{}(\alpha )\frac{\mathrm{\Delta }(\alpha )}{J_1}(1\alpha )^{3/4}(1+\alpha )^{1/4},$$
(47)
$$\overline{\mathrm{\Delta }^{}}(\delta )\frac{\mathrm{\Delta }(\delta )}{J}2\delta ^{3/4}.$$
(48)
The same $`\overline{\mathrm{\Delta }^{}}(\delta )`$ was found in numerical calculations by Ladavac et al. for $`0.01\delta 1`$, whereas calculations for $`0.03\delta 0.06`$ by Augier et al. yielded somewhat smaller values of $`\overline{\mathrm{\Delta }^{}}`$ than predicted by Eq. (48). The asymptotic critical behavior of $`\overline{\mathrm{\Delta }^{}}`$ as the uniform limit is approached ($`\alpha 1,\delta 0`$) has been given as
$$\overline{\mathrm{\Delta }^{}}(\delta )\frac{\delta ^{2/3}}{|\mathrm{ln}\delta |^{1/2}};$$
(49)
thus the parametrization in Eq. (48) evidently indicates that the fitted data do not reside within the asymptotic critical regime. Alternatively, Barnes, Riera, and Tennant suggested that Eq. (49) may not be the correct form for the asymptotic critical behavior. On the other hand, Uhrig et al. fitted their $`T=0`$ density matrix renormalization group (DMRG) calculations of $`\overline{\mathrm{\Delta }^{}}(\delta )`$ for $`0.004\delta 0.1`$ to a power-law behavior without the log correction and obtained $`\overline{\mathrm{\Delta }^{}}1.57\delta ^{0.65}`$. We will further discuss the above spin gap calculation results later in Sec. VII A after deriving our own $`\mathrm{\Delta }^{}(\alpha )`$ values from our TMRG $`\chi ^{}(\alpha ,t)`$ data in Sec. VI.
#### 3 One-magnon dispersion relations
Barnes, Riera, and Tennant have computed the dimer series expansion of the dispersion relation $`\epsilon (\alpha ,k)E(\alpha ,k)/J_1`$ for the one-magnon ($`S=1`$) energy $`E(\alpha ,k)`$ vs wave vector $`k`$ along the chain for the lowest-lying one-magnon band, up to fifth order in $`\alpha `$, which we write as
$$\epsilon (\alpha ,k)=\underset{n=0}{\overset{\mathrm{}}{}}a_n(\alpha )\mathrm{cos}(2nka),$$
(50)
where $`a`$ is the (average) spin-spin distance, which is 1/2 the lattice repeat distance along the chain in the dimerized state. Plots of $`\epsilon (\alpha ,k)`$ for $`\alpha =0`$, 0.2, 0.4, 0.6, and 0.8 up to fifth order in $`\alpha `$, as given in Fig. 4 of Ref. , are shown as the dashed curves in Fig. 2. The curves are symmetric about $`ka=\pi /2`$, so the same spin gap $`\mathrm{\Delta }^{}(\alpha )=_{n=0}^{\mathrm{}}a_n(\alpha )`$ occurs at $`ka=0`$ and $`\pi `$. This fifth-order approximation yields $`\mathrm{\Delta }^{}(\alpha )`$ values for $`\alpha 0.9`$ in rather close agreement with BRTs’ results discussed in the previous section. For a dimer series expansion we expect the average energy of the one-magnon band states to be nearly independent of $`\alpha `$, i.e.,
$$\frac{1}{\pi }_0^\pi \epsilon (\alpha ,ka)d(ka)=1.$$
(51)
Indeed, upon inserting BRTs’ fifth order expansion coefficients into Eq. (50) and the result into Eq. (51), we find that this sum rule is satisfied to within 1% for $`0\alpha 1`$.
Also shown as a solid curve in Fig. 2 is the exact result $`\epsilon (k)=(\pi /2)|\mathrm{sin}(ka)|`$ for the uniform chain ($`\alpha =1`$). This $`\epsilon (k)`$ has a cusp with infinite curvature (at $`ka=0`$
and $`\pi `$) which cannot be accurately approximated by a Fourier series with a small number of terms. This singular behavior is evidently closely related to the critical behavior of $`\mathrm{\Delta }^{}(\alpha 1)`$ discussed above. In order to later model our TMRG $`\chi ^{}(\alpha ,t)`$ data close to, but not in, the low-$`t`$ limit, we will need an expression for $`\epsilon (\alpha ,k)`$ which is correct in the limit $`\alpha 1`$ and which also reproduces reasonably well the $`\epsilon (\alpha ,k)`$ of BRT. We found that the simple one-parameter ($`\mathrm{\Delta }^{}`$) form suggested earlier by one of us in the context of the $`S=1/2`$ two-leg ladder
$$\epsilon (\mathrm{\Delta }^{},k)=[\mathrm{\Delta }_{}^{}{}_{}{}^{2}+f^2(\mathrm{\Delta }^{})\mathrm{sin}^2(ka)]^{1/2},$$
(53)
is satisfactory in these regards for the AF alternating-exchange chain over the entire range $`0\mathrm{\Delta }^{}1`$. The function $`f(\mathrm{\Delta }^{})`$ is determined here by the sum rule (51), which yields the condition
$$\mathrm{E}\left[\frac{f^2(\mathrm{\Delta }^{})}{\mathrm{\Delta }_{}^{}{}_{}{}^{2}}\right]=\frac{\pi }{2\mathrm{\Delta }^{}},$$
(54)
where E($`x`$) is the complete elliptic integral of the second kind. From Eq. (54), $`f`$ varies nonlinearly with $`\mathrm{\Delta }^{}`$ from $`f(\mathrm{\Delta }^{}=0)=\pi /2`$ to $`f(\mathrm{\Delta }^{}=1)=0`$, as shown in Fig. 3. From an independently determined dependence of $`\mathrm{\Delta }^{}`$ on $`\alpha `$ as in Eq. (47), one can then determine $`f(\alpha )`$ as also shown in Fig. 3. Using the fifth-order $`\mathrm{\Delta }^{}(\alpha )`$ values of BRT in Fig. 2, the resulting dispersion relations (3) for $`\alpha =0,`$ 0.2, 0.4, 0.6, and 0.8 were calculated and are shown as the solid curves in Fig. 2, where they are seen to be in close agreement with the respective dashed curves of BRT. An important difference for large $`\alpha `$, however, is that the $`\epsilon (\mathrm{\Delta }^{},k)`$ in Eqs. (3) properly reduces by construction to the exact $`\epsilon (\alpha ,k)`$ for $`\alpha 1`$, whereas the one in Eq. (50) with a finite number of terms does not.
Close to the one-magnon band minimum, the square root and the sine function in the dispersion relation in Eq. (53) can be expanded, yielding
$$\epsilon (\mathrm{\Delta }^{},k0)\mathrm{\Delta }^{}+\frac{1}{2}\frac{f^2(\mathrm{\Delta }^{})}{\mathrm{\Delta }^{}}(ka)^2.$$
(55)
A comparison of Eqs. (55) and (16) shows that the parameter $`c^{}`$ in the formulas for $`\chi ^{}(t0)`$ \[Eqs. (16)\] and $`C(t0)`$ \[Eq. (25)\] is a unique function of $`\mathrm{\Delta }^{}`$ which in our approximation is given by
$$c^{}(\mathrm{\Delta }^{})=\frac{1}{2}\frac{f^2(\mathrm{\Delta }^{})}{\mathrm{\Delta }^{}},$$
(56)
with $`f^2(\mathrm{\Delta }^{})`$ given by Eq. (54). Thus both $`\chi ^{}(t0)`$ and $`C(t0)`$ for the alternating-exchange chain only depend on the single parameter $`\mathrm{\Delta }^{}`$ (in addition to $`t`$). Explicitly, we obtain
$$\chi ^{}(t0)=\frac{1}{\sqrt{2\pi }f(\mathrm{\Delta }^{})}\left(\frac{\mathrm{\Delta }^{}}{t}\right)^{1/2}\mathrm{e}^{\mathrm{\Delta }^{}/t}.$$
(57)
As might have been anticipated, the only thermodynamic variable is the ratio $`t/\mathrm{\Delta }^{}=k_\mathrm{B}T/\mathrm{\Delta }`$ of the thermal energy to the spin gap. The numerical prefactor depends explicitly (only) on the reduced spin gap $`\mathrm{\Delta }^{}\mathrm{\Delta }/J_1`$. Similarly, the magnetic specific heat is obtained as
$`{\displaystyle \frac{C(t0)}{Nk_\mathrm{B}}}={\displaystyle \frac{3}{\sqrt{2\pi }}}`$ $`{\displaystyle \frac{\mathrm{\Delta }^{}}{f(\mathrm{\Delta }^{})}}\left({\displaystyle \frac{\mathrm{\Delta }^{}}{t}}\right)^{3/2}`$ (59)
$`\times \left[1+{\displaystyle \frac{t}{\mathrm{\Delta }^{}}}+{\displaystyle \frac{3}{4}}\left({\displaystyle \frac{t}{\mathrm{\Delta }^{}}}\right)^2\right]\mathrm{e}^{\mathrm{\Delta }^{}/t},`$
where again the same characteristics are present as just discussed for $`\chi ^{}(t)`$. The variations of the prefactors with $`\mathrm{\Delta }^{}`$ for $`\chi ^{}(t)`$ and $`C(t)`$ can both be ascertained from the plot of $`f(\mathrm{\Delta }^{})`$ in Fig. 3. In particular, when $`\alpha 1`$ ($`\delta 1`$) for which $`\mathrm{\Delta }^{}1`$, $`f`$ is nearly a constant. For our and our readers’ convenience when modeling materials showing small spin gaps, we have fitted our numerical $`f(\mathrm{\Delta }^{})`$ calculations for $`0\mathrm{\Delta }^{}0.4`$ by a third-order polynomial to within 2 parts in $`10^4`$, given by
$$f(\mathrm{\Delta }^{})=\frac{\pi }{2}0.034289\mathrm{\Delta }^{}1.18953\mathrm{\Delta }_{}^{}{}_{}{}^{2}+0.40030\mathrm{\Delta }_{}^{}{}_{}{}^{3}.$$
(60)
By a change in variables to $`(J,\delta )`$ and using the $`\mathrm{\Delta }^{}(\alpha )`$ in Eq. (II C 2), we obtain the following forms which are more useful for modeling materials with small spin gaps, especially those showing second-order spin dimerization transitions with decreasing $`T`$:
$$\overline{\chi ^{}}(\overline{t}0)=\frac{1}{(1+\delta )\sqrt{2\pi }f(\overline{\mathrm{\Delta }^{}})}\left(\frac{\overline{\mathrm{\Delta }^{}}}{\overline{t}}\right)^{1/2}\mathrm{e}^{\overline{\mathrm{\Delta }^{}}/\overline{t}}.$$
(62)
$`{\displaystyle \frac{\overline{C}(t0)}{Nk_\mathrm{B}}}=`$ $`{\displaystyle \frac{3}{(1+\delta )\sqrt{2\pi }}}{\displaystyle \frac{\overline{\mathrm{\Delta }^{}}}{\overline{f}(\overline{\mathrm{\Delta }^{}})}}\left({\displaystyle \frac{\overline{\mathrm{\Delta }^{}}}{\overline{t}}}\right)^{3/2}`$ (64)
$`\times \left[1+{\displaystyle \frac{\overline{t}}{\overline{\mathrm{\Delta }^{}}}}+{\displaystyle \frac{3}{4}}\left({\displaystyle \frac{\overline{t}}{\overline{\mathrm{\Delta }^{}}}}\right)^2\right]\mathrm{e}^{\overline{\mathrm{\Delta }^{}}/\overline{t}},`$
$$\overline{f}(\overline{\mathrm{\Delta }^{}})=\frac{\pi }{2}0.033933\overline{\mathrm{\Delta }^{}}1.19607\overline{\mathrm{\Delta }^{}}^2+0.92430\overline{\mathrm{\Delta }^{}}^3.$$
(65)
Note that in these formulas $`\mathrm{\Delta }^{}/t=\overline{\mathrm{\Delta }^{}}/\overline{t}=\mathrm{\Delta }/(k_\mathrm{B}T)`$.
## III Theory: $`𝑺\mathbf{=}\mathrm{𝟏}\mathbf{/}\mathrm{𝟐}`$ Uniform Heisenberg Chain
### A Magnetic spin susceptibility
The uniform $`S=1/2`$ chain is one limit of the alternating-exchange chain with $`J_{ij}J,\alpha =1,\delta =0`$, and with no spin gap \[the $`\chi ^{}(t0)`$ and $`C(t0)/t`$ are finite\]. The spin susceptibility was calculated accurately by Eggert, Affleck and Takahashi in 1994, and recently refined by Klümper as shown in Fig. 4(a) where only the calculations up to $`t=2`$ are shown. An expanded plot of the data for $`t0.02`$, including the exact value $`1/\pi ^2`$ at $`t=0`$, is shown in Fig. 4(b), along with a fit (Fit 2) to the data to be derived and discussed in Sec. IV C. The most recent calculations of Ref. have an absolute accuracy estimated to be $`1\times 10^9`$ and show a broad maximum at a temperature $`T^{\mathrm{max}}`$ with a value $`\chi ^{\mathrm{max}}`$. By fitting data points near the maximum by up to 6th order polynomials, we determined these numerical values to be given by
$$T^{\mathrm{max}}=\mathrm{0.6\hspace{0.17em}408\hspace{0.17em}510}(4)J/k_\mathrm{B},$$
(67)
$$\frac{\chi ^{\mathrm{max}}J}{Ng^2\mu _\mathrm{B}^2}=\mathrm{0.146\hspace{0.17em}926\hspace{0.17em}279}(1),$$
(68)
$$\chi ^{\mathrm{max}}T^{\mathrm{max}}=\mathrm{0.0\hspace{0.17em}941\hspace{0.17em}579}(1)\frac{Ng^2\mu _\mathrm{B}^2}{k_\mathrm{B}}.$$
(69)
These values are consistent within the errors with those found by Eggert et al., but are much more accurate. For one mole of spins, setting $`N=N_\mathrm{A}`$ (Avogadro’s number) in Eq. (69) yields
$$\chi ^{\mathrm{max}}T^{\mathrm{max}}=\mathrm{0.0\hspace{0.17em}353\hspace{0.17em}229}(3)g^2\frac{\mathrm{cm}^3\mathrm{K}}{\mathrm{mol}}.$$
(70)
Note that the product $`\chi ^{\mathrm{max}}T^{\mathrm{max}}`$ in Eqs. (69) and (70) is independent of $`J`$, and hence is a good initial test of whether the $`S=1/2`$ AF uniform Heisenberg chain model might be applicable to a particular compound.
#### 1 High-temperature series expansions
The coefficients $`c_n`$ of the HTSE for $`\chi ^{}(t)`$,
$$4\chi ^{}t=\underset{n=0}{\overset{\mathrm{}}{}}\frac{c_n}{t^n},$$
(72)
are given up to $`𝒪(1/t^7)`$ by
$`c_0`$ $`=`$ $`1,c_1={\displaystyle \frac{1}{2}},c_2=0,c_3={\displaystyle \frac{1}{24}},c_4={\displaystyle \frac{5}{384}},`$ (73)
$`c_5`$ $`=`$ $`{\displaystyle \frac{7}{1280}},c_6={\displaystyle \frac{133}{30720}},c_7={\displaystyle \frac{1}{4032}}.`$ (74)
Inverting the series, we obtain the corresponding $`d_n`$ coefficients in Eq. (6) as
$`d_0`$ $`=`$ $`1,d_1={\displaystyle \frac{1}{2}},d_2={\displaystyle \frac{1}{4}},d_3={\displaystyle \frac{1}{12}},d_4={\displaystyle \frac{1}{128}},`$ (75)
$`d_5`$ $`=`$ $`{\displaystyle \frac{29}{3840}},d_6={\displaystyle \frac{317}{92160}},d_7={\displaystyle \frac{11}{71680}}.`$ (76)
The $`d_n`$ coefficients with $`n=0,`$ 1, 2, and 3 are of course in agreement with Eq. (38) for $`\alpha =1`$.
#### 2 Logarithmic corrections at low temperatures
At low temperatures a simple expansion of thermodynamic properties in for instance the variable $`t`$ is not possible. Such a nonanalyticity in $`t`$ can be viewed as due to the strong correlation of the quasiparticles, i.e., the elementary excitations of the system are not strictly free; they show rather nontrivial scattering processes. Spinons with low energies $`ϵ_1`$ and $`ϵ_2`$ have a scattering phase $`\varphi (ϵ_1,ϵ_2)\varphi _0+\text{const}/|\mathrm{log}(ϵ_1ϵ_2)|`$. From this property it is clear that an expansion in the single variable $`t`$ is not possible, but has to be supplemented by a term $`1/\mathrm{log}(t)`$. Although being very intuitive, this physical picture on the basis of scattering processes of spinons has not played any important role in the investigation of logarithmic corrections until recently.
Logarithmic dependencies of physical quantities have been known for the spin-1/2 Heisenberg chain for a rather long time. Usually, a quantum chain with critical couplings leads to critical correlations only in the thermodynamic limit $`1/L=0`$ and at $`T=0`$, where $`L`$ is the length of the chain. If one of these conditions is not met the physical properties receive nonanalytic contributions in terms of $`1/L`$ or $`T`$. From the renormalization group point of view the existence of logarithmic corrections is reflected by the perturbation of the (critical) fixed point Hamiltonian by some marginal operator. Such operators usually exist only for isotropic systems.
The investigation of the size dependence of energy levels of critical quantum chains was started more than a decade ago. For the isotropic spin-1/2 Heisenberg chain, expansions in $`1/L`$ and additional logarithmic corrections ($`1/L\mathrm{log}L`$ etc.) were found in lattice approaches (Bethe ansatz) as well as in field theory \[RG study of the Wess-Zumino-Witten (WZW) model with topological term $`k=1`$ (Refs. )\].
Many of these earlier results are still relevant for the issues discussed in this section due to an equivalence of many-particle systems at $`T=0`$, $`1/L>0`$ (groundstate properties of finite chains) and those at $`T>0`$, $`1/L=0`$ (thermodynamics of the bulk). This leads to asymptotic series where $`T`$ and $`1/L`$ play very similar roles. To our knowledge the first explicit report on $`\mathrm{log}(T)`$ corrections in the magnetic susceptibility resulting in an infinite slope at $`T=0`$ was given in Ref. . Including higher order terms, the asymptotic expansion $`\chi _{\mathrm{lt}}^{}(t)`$ for $`\chi ^{}(t)`$ is
$`\chi _{\mathrm{lt}}^{}(t)`$ $`=`$ $`{\displaystyle \frac{1}{\pi ^2}}\left[1+{\displaystyle \frac{1}{2}}{\displaystyle \frac{\mathrm{ln}\left(+\frac{1}{2}\right)}{(2)^2}}+\mathrm{}\right],`$ (78)
$`=`$ $`{\displaystyle \frac{1}{\pi ^2}}\left[1+{\displaystyle \frac{1}{2}}{\displaystyle \frac{\mathrm{ln}}{(2)^2}}{\displaystyle \frac{1}{(2)^3}}+\mathrm{}\right],`$ (79)
$$\mathrm{ln}(t_0/t),$$
(80)
where $`t_0`$ is a nonuniversal (undetermined) parameter. In Ref. the field theoretical prediction on the basis of the WZW model was compared with the results of thermodynamic Bethe ansatz calculations and showed convincing agreement in an intermediate temperature regime. Using up to the first logarithmic correction term in Eq. (78), Eggert, Affleck, and Takahashi estimated $`t_07.7`$, so at low temperatures $`t0.01`$ the parameter $`1`$.
A general feature of field theoretical and lattice approaches is their restriction to “low” and “high” temperatures, respectively. Field theoretical studies suffer at high temperatures from the neglect of (infinitely many) irrelevant operators. Lattice studies show convergence problems at low temperatures as increasingly larger systems have to be studied in order to avoid finite-size effects. In addition, the comparison of field theory and lattice results can only verify or falsify the universal aspects of an asymptotic expansion. Non-universal quantities like $`t_0`$ which derive from some coupling constant of a marginal or irrelevant operator (undetermined within the field theory) can at best be fitted as done in Ref. .
The latter problem of determining the coupling constants in an effective field theory was solved by Lukyanov who used a bosonic representation of the Heisenberg chain. The coupling constants were fixed by a comparison of the susceptibility data $`\chi (T=0,h)`$ obtained by him with Bethe ansatz calculations for magnetic field $`h`$ at $`T=0`$. Eventually, the $`\chi (T>0,h=0)`$ data could be calculated without any need of a fit parameter.
Lukyanov obtained the following analytical low-temperature expansion of $`\chi ^{}(t)`$,
$`\chi _{\mathrm{lt},\mathrm{g}}^{}(t)={\displaystyle \frac{1}{\pi ^2}}\{1`$ $`+`$ $`{\displaystyle \frac{g}{2}}+{\displaystyle \frac{3g^3}{32}}+𝒪(g^4)`$ (82)
$`+`$ $`{\displaystyle \frac{\sqrt{3}}{\pi }}t^2[1+𝒪(g)]\},`$ (83)
where $`g`$ obeys the transcendental equation
$$\frac{1}{2}\mathrm{ln}g+\frac{1}{g}=$$
(84)
or equivalently
$$\sqrt{g}\mathrm{e}^{1/g}=\frac{t_0}{t},$$
(85)
with a unique value of $`t_0`$ given by
$$t_0=\sqrt{\frac{\pi }{2}}\mathrm{e}^{\gamma +(1/4)}\mathrm{2.866\hspace{0.17em}257\hspace{0.17em}058},$$
(86)
where $`\gamma \mathrm{0.577\hspace{0.17em}215\hspace{0.17em}665}`$ is Euler’s constant. Lukyanov showed that his $`\chi _{\mathrm{lt},\mathrm{g}}^{}(t)`$ is in agreement with the numerical data of Eggert, Affleck and Takahashi at low temperatures $`t0.003`$.
In the following, we will compare high-accuracy numerical Bethe ansatz calculations carried out to much lower temperatures by Klümper and Johnston with this theory in some detail because this theory is exact at low temperatures with no adjustable parameters. The calculations of Ref. are based on lattice studies, however without suffering from the usual shortcomings. By means of a lattice path integral representation of the finite temperature Heisenberg chain and the formulation of a suitable quantum transfer matrix (both quite analogous to the numerical TMRG calculations presented later in this paper) a set of numerically well-posed expressions for the free energy was derived. In more physical terms the method can be understood as an application of the familiar though often rather vague concept of quasiparticles to a quantitative description of the many particle system valid for all temperatures $`T`$ and magnetic field values $`h`$. The work can be understood as an evaluation of the full scattering theory of spinons and antispinons.
Our iterative solution of Eq. (84) yields the expansion
$$g=\frac{1}{}\left\{1\frac{\mathrm{ln}}{2}+\frac{(\mathrm{ln})^2\mathrm{ln}}{(2)^2}+𝒪\left[\frac{1}{(2)^3}\right]\right\}.$$
(87)
A log-log plot of $`g`$ vs $`t`$ obtained by numerically solving Eq. (85) is shown in Fig. 5 (solid curve), along with its lowest-order approximation $`g(t)1/=1/\mathrm{ln}(t_0/t)`$ (dashed curve). This approximation is 1.1% larger than the exact result at $`t=10^{30}`$, with the discrepancy increasing steadily to 5.8% at $`t=10^{15}`$ and 8.5% at $`t=10^7`$. Substituting Eq. (87) into (83) gives
$`\chi _{\mathrm{lt},\mathrm{log}}^{}(t)={\displaystyle \frac{1}{\pi ^2}}\{1`$ $`+`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{\mathrm{ln}}{(2)^2}}`$ (88)
$`+`$ $`{\displaystyle \frac{(\mathrm{ln})^2\mathrm{ln}+(3/4)}{(2)^3}}+𝒪\left[{\displaystyle \frac{1}{(2)^4}}\right]`$ (89)
$`+`$ $`{\displaystyle \frac{\sqrt{3}}{\pi }}t^2[1+𝒪\left({\displaystyle \frac{1}{2}}\right)]\}.`$ (90)
The first three terms are identical with those in Eq. (79), but the constant term in the numerator of the fourth term is not the same as in Eq. (79), indicating that Eq. (79) is not accurate to order higher than $`𝒪[1/(2)^2]`$.
An important issue is the accuracy to which the log expansion approximation $`\chi _{\mathrm{lt},\mathrm{log}}^{}(t)`$ in Eq. (90) reproduces the $`\chi _{\mathrm{lt},\mathrm{g}}^{}(t)`$ prediction of the original Eqs. (III A 2). We have calculated both quantities to high accuracy and plot the difference vs $`t`$, for the range $`10^{30}t0.5`$, in Fig. 6. The $`\chi _{\mathrm{lt},\mathrm{log}}^{}(t)`$ is seen to increasingly diverge from $`\chi _{\mathrm{lt},\mathrm{g}}^{}(t)`$ with increasing $`t`$.
When comparing the predictions of Lukyanov’s theory with numerical results such as obtained from the Bethe ansatz, it is important to know at what temperature the low temperature expansion in Eqs. (III A 2) ceases to be accurate (“accurate” must be defined) with increasing temperature. There are three aspects of this issue that need to be addressed. The first and second aspects concern the temperatures at which the unknown $`𝒪(g^4)`$ and $`𝒪(g)`$ terms in Eq. (83) become significant, respectively; we will return to these two issues shortly. The third aspect is whether the log expansion approximation $`\chi _{\mathrm{lt},\mathrm{log}}^{}(t)`$ in Eq. (90) can be used in this comparison. The absolute accuracy of the most recent Bethe ansatz calculations is estimated to be $`1\times 10^9`$. From Fig. 6, we see that $`\chi _{\mathrm{lt},\mathrm{log}}^{}(t)`$ approximates $`\chi _{\mathrm{lt},\mathrm{g}}^{}(t)`$ to this degree only for temperatures $`t10^{30}`$ \[we infer that the previous Eqs. (III A 2) only apply to this accuracy at similarly very low temperatures\]. Therefore, to avoid this unnecessary approximation as a source of error at higher temperatures, we will henceforth compare the numerical Bethe ansatz calculations with $`\chi _{\mathrm{lt},\mathrm{g}}^{}(t)`$ calculated from Lukyanov’s original Eqs. (III A 2).
A comparison of the low-temperature Bethe ansatz $`\chi ^{}(t)`$ calculations and Lukyanov’s theoretical $`\chi _{\mathrm{lt},\mathrm{g}}^{}`$ prediction is shown in Fig. 7(a). On the scale of this figure, the two results are identical. The (small) quantitative differences between them are shown as the filled circles in Fig. 7(b). The lower error bar on each data point in Fig. 7(b) is $`1\times 10^7`$ to indicate the scale. The upper error bar is the estimated uncertainty in $`\chi _{\mathrm{lt},\mathrm{g}}^{}`$ arising from the presence of the unknown $`𝒪(g^4)`$ and higher-order terms in Eq. (83), which was set to $`g^4(t)/\pi ^2`$; the uncertainty in the $`t^2`$ contribution, $`\sqrt{3}t^2g(t)/\pi ^3`$, is negligible at low $`t`$ compared to this. At the lower temperatures, the data agree extremely well with the prediction of Lukyanov’s theory. At the highest temperatures $`t10^3`$, higher or-
der $`t^n`$ terms also become important, as inferred from our empirical fits (Fits 1 and 2) below to the numerical data.
Irrespective of the uncertainties in the theoretical prediction at high temperatures just discussed, we can safely conclude directly from Fig. 7(b) that the Bethe ansatz $`\chi ^{}(t)`$ data are in agreement with the exact theory of Lukyanov to within an absolute accuracy of $`1\times 10^6`$ (relative accuracy $`10`$ ppm) over a temperature range spanning 18 orders of magnitude from $`t=5\times 10^{25}`$ to $`t=5\times 10^7`$. The agreement is much better than this at the lower temperatures.
### B Magnetic specific heat
The magnetic specific heat $`C`$ of the $`S=1/2`$ AF uniform Heisenberg chain was recently calculated to high accuracy by Klümper and Johnston over the temperature range $`5\times 10^{25}k_\mathrm{B}T/J5`$. The accuracy is estimated to be $`3\times 10^{10}C(t)`$. The results for $`T2J/k_\mathrm{B}`$ are shown in Fig. 8(a). The initial $`T`$ dependence is approximately (see below) linear, and is given exactly in
the $`t=0`$ limit by
$$\frac{C(t0)}{Nk_\mathrm{B}}=\frac{2}{3}t.$$
(91)
The data show a maximum with a value $`C^{\mathrm{max}}`$ at a temperature $`T_C^{\mathrm{max}}`$. By fitting 3–7 data points in the vicinity of the maximum by up to 6th order polynomials, these values were found to be
$`{\displaystyle \frac{k_\mathrm{B}T_C^{\mathrm{max}}}{J}}`$ $`=`$ $`\mathrm{0.48\hspace{0.17em}028\hspace{0.17em}487}(1),`$ (92)
$`{\displaystyle \frac{C^{\mathrm{max}}}{Nk_\mathrm{B}}}`$ $`=`$ $`\mathrm{0.3\hspace{0.17em}497\hspace{0.17em}121\hspace{0.17em}235}(2).`$ (94)
The electronic specific heat coefficient $`C(T)/T`$ is plotted vs temperature in Fig. 8(b). As expected from Eq. (91), the data approach the value $`(2/3)Nk_\mathrm{B}^2/J`$ for $`t0`$. The initial deviation from this constant value is positive and approximately (see below) quadratic in $`t`$. The data exhibit a smooth maximum with a value $`(C/T)^{\mathrm{max}}`$ at a temperature $`T_{\mathrm{C}/\mathrm{T}}^{\mathrm{max}}`$, values which we determined by fitting polynomials to the data in the vicinity of the peak to be
$`{\displaystyle \frac{k_\mathrm{B}T_{C/T}^{\mathrm{max}}}{J}}`$ $`=`$ $`\mathrm{0.30\hspace{0.17em}716\hspace{0.17em}996}(2),`$ (95)
$`{\displaystyle \frac{(C/T)^{\mathrm{max}}J}{Nk_\mathrm{B}^2}}`$ $`=`$ $`\mathrm{0.8\hspace{0.17em}973\hspace{0.17em}651\hspace{0.17em}576}(5).`$ (97)
The magnetic entropy $`S(T)`$ is determined by integrating the $`C(T)/T`$ data in Fig. 8(b) vs $`T`$ and the result, normalized by $`S(T\mathrm{})=Nk_\mathrm{B}\mathrm{ln}2`$, is plotted vs $`T`$ in Fig. 9. This figure allows one to estimate the maximum magnetic entropy that can be associated with a dimerization transition or any other magnetic transition involving $`S=1/2`$ Heisenberg chains which are weakly coupled to each other \[assuming that the (average) $`J`$ does not change at the transition\]. For example, for $`\mathrm{NaV}_2\mathrm{O}_5`$ where $`k_\mathrm{B}T_\mathrm{c}/J0.057`$, one can estimate from Fig. 9 that the magnetic entropy at $`T_\mathrm{c}`$ cannot exceed $`0.056R\mathrm{ln}2=0.32`$ J/mol K, where $`R`$ is the molar gas constant. The reason this value is the upper limit is that magnetic critical fluctuations will increase the specific heat, and hence the entropy, above $`T_\mathrm{c}`$ and thus reduce it at (and below) $`T_\mathrm{c}`$, by conservation of magnetic entropy, compared to the values for the isolated chain at the same reduced temperatures. Similarly, the $`C(T)`$ data in Fig. 8(a) allow one to estimate the minimum lattice specific heat contribution $`C^{\mathrm{lat}}(T)`$ above $`T_\mathrm{c}`$ if the $`C^{\mathrm{lat}}(T)`$ has not been determined previously from experiments and/or theory directly.
At low temperatures, the electronic specific heat coefficient $`C(T)/T`$ becomes independent of temperature (apart from logarithmic corrections, see below), as does the spin susceptibility $`\chi ^{}(t)`$, just as in a metal (Fermi liquid). Therefore it is of interest to compute a normalized ratio of these two quantities. For a metal, the relevant quantity is the Wilson-Sommerfeld ratio, which for
$`S=1/2`$ quasiparticles reads, in the notation of this paper and with $`k_\mathrm{B}`$ set to 1,
$$R_\mathrm{W}(t)=\frac{4\pi ^2\chi ^{}(t)t}{3C(t)}.$$
(98)
In a degenerate free electron gas, $`R_\mathrm{W}=1`$ and is independent of $`t`$. For exchange-enhanced metals $`1<R_\mathrm{W}10`$, for $`S=1/2`$ Kondo impurities in a metal the Wilson ratio associated with the impurities is $`R_\mathrm{W}=2`$, and for many heavy fermion metals $`R_\mathrm{W}2`$. Plotted in Fig. 10 is $`R_\mathrm{W}(t)`$ for the $`S=1/2`$ AF Heisenberg chain, where $`C(t)/t`$ and $`\chi ^{}(t)`$ were given above in Figs. 8(b) and 4, respectively. For $`t0`$, the Wilson ratio for the $`S=1/2`$ Heisenberg chain is exactly 2. With increasing $`t`$, $`R_\mathrm{W}`$ is seen to be nearly independent of $`t`$ to within $`\pm 10`$% up to $`t0.4`$, but the influence of the logarithmic corrections to both $`\chi (T)`$ and $`C(T)`$ are quantitatively important. Although the logarithmic corrections for $`\chi (T)`$ and $`C(T)`$ oppose each other in their ratio in $`R_\mathrm{W}(t)`$, the logarithmic corrections for $`\chi (t)`$ win out, giving a net $`10`$% increase in $`R_\mathrm{W}(t)`$ with increasing $`t`$ at low $`t`$. At higher $`t`$, the system crosses over to the expected local moment Heisenberg behavior where $`R_\mathrm{W}t^2`$. Thus as far as the thermodynamics is concerned, the uniform Heisenberg chain behaves at low temperatures as expected for a Fermi liquid, apart from the influence of the logarithmic corrections. This quasi-Fermi liquid behavior arises because the elementary excitations at low temperatures are $`S=1/2`$ spinons which are fermions with a Fermi surface (i.e., Fermi points in one dimension). Since the spinons carry no charge, the chain is an insulator. The deviation of the Wilson ratio from unity and the logarithmic corrections are due to spinon interactions.
#### 1 High-temperature series expansions
The HTSE for the specific heat of a spin $`S`$ AF uniform Heisenberg chain is
$$\frac{C(T)}{Nk_\mathrm{B}}=\frac{x^2}{3t^2}\left[1+\underset{n=1}{\overset{\mathrm{}}{}}\frac{c_n(x)}{t^n}\right],$$
(100)
$$x=S(S+1),t=\frac{k_\mathrm{B}T}{J},$$
(101)
$`c_1={\displaystyle \frac{1}{2}},c_2={\displaystyle \frac{1}{15}}(38x3x^2),`$
$`c_3={\displaystyle \frac{1}{36}}(316x4x^2),`$
$`c_4={\displaystyle \frac{1}{5040}}(1921432x+1123x^2+800x^3+160x^4),`$
$$c_5=\frac{1}{21600}(4143768x+6635x^2+2624x^3+480x^4).$$
(102)
Specializing Eqs. (III B 1) to $`S=1/2`$ $`(x=3/4`$) then gives
$$\frac{C(T)}{Nk_\mathrm{B}}=\frac{3}{16t^2}\left[1+\underset{n=1}{\overset{\mathrm{}}{}}\frac{c_n}{t^n}\right],$$
(104)
$$c_1=\frac{1}{2},c_2=c_3=\frac{5}{16},c_4=\frac{7}{256},c_5=\frac{917}{7680}.$$
(105)
The two $`C(T)`$ HTSE terms of order $`1/t^2`$ and $`1/t^3`$ in Eqs. (102) are in agreement with the general expression for the two lowest-order HTSE expansion terms for $`C(T)`$ of the $`S=1/2`$ alternating-exchange Heisenberg chain in Eq. (45) with alternation parameter $`\alpha =1`$.
In a later section, the Bethe ansatz $`C(T)`$ data will be fitted to obtain a function accurately representing the $`C(T)`$ of the $`S=1/2`$ AF uniform Heisenberg chain. In order that we are not required to change our fitting equations from those we use for fitting magnetic susceptibility data, the coefficients for the series inverted from that in Eq. (104) are required. We obtain
$$\frac{C(T)}{Nk_\mathrm{B}}=\frac{3}{16t^2}\left[1+\underset{n=1}{\overset{\mathrm{}}{}}\frac{d_n}{t^n}\right]^1,$$
(107)
$$d_1=\frac{1}{2},d_2=\frac{9}{16},d_3=\frac{1}{8},d_4=\frac{7}{128},d_5=\frac{7}{1920}.$$
(108)
#### 2 Low-temperature logarithmic corrections
At first sight, from Fig. 8 there appear to be no singularities in the temperature dependence of the specific heat for the $`S=1/2`$ AF uniform Heisenberg chain. However, if the electronic specific heat coefficient $`C(T)/T`$ is examined in detail, one sees anomalous behavior at low temperatures. Shown as the top curve in Fig. 11(a) is a plot of the difference between the electronic specific heat coefficient and its zero-temperature value, $`\mathrm{\Delta }C(t)/Nk_\mathrm{B}t[C(t)(2/3)t]/(Nk_\mathrm{B}t)`$ for $`0t0.1`$ \[compare with Fig. 8(b)\]. From this figure, there is still nothing particularly strange about the data. However, upon further expanding the plot to study the range $`0t0.005`$ as shown in Fig. 11(b), we see that $`\mathrm{\Delta }C/Nk_\mathrm{B}t`$ is developing an infinite slope as $`t0`$. This is the signature of the existence of logarithmic corrections to the specific heat at temperatures $`t1`$, just as it was for the magnetic susceptibility.
Klümper, Lukyanov, and others have found a logarithmic correction to the low-$`t`$ limit in Eq. (91). Lukyanov’s exact asymptotic expansion for the free energy per spin in zero magnetic field is
$`f=J\mathrm{ln}2`$ $``$ $`{\displaystyle \frac{(k_\mathrm{B}T)^2}{3J}}\left[1+{\displaystyle \frac{3}{8}}g^3+𝒪(g^4)\right]`$ (109)
$``$ $`{\displaystyle \frac{3^{3/2}(k_\mathrm{B}T)^4}{10\pi J^3}}[1+𝒪(g)],`$ (111)
where $`g(t/t_0)`$ and $`t_0`$ are the same as given in Eqs. (84) and (86), respectively, and where $`g(t)`$ was plotted in Fig. 5. The specific heat at constant volume is calculated using $`C=T^2f/T^2`$, yielding
$`{\displaystyle \frac{C_{\mathrm{lt},\mathrm{g}}(T)}{Nk_\mathrm{B}}}={\displaystyle \frac{2k_\mathrm{B}T}{3J}}[1`$ $`+`$ $`{\displaystyle \frac{3}{8}}g^3+𝒪(g^4)]`$ (112)
$`+`$ $`{\displaystyle \frac{2(3^{5/2})}{5\pi }}\left({\displaystyle \frac{k_\mathrm{B}T}{J}}\right)^3[1+𝒪(g)].`$ (114)
This formula shows that the electronic specific heat coefficient $`C(T)/T`$ increases quadradically with $`T`$ at low $`T`$ (after subtracting the logarithmic corrections). The numerical prefactor of the $`t^3`$ term is $`1.98478\mathrm{}`$. If the approximate expansion for $`g()`$ in Eq. (87) is substituted into Eq. (114), one obtains
$`{\displaystyle \frac{C_{\mathrm{lt},\mathrm{log}}(T)}{Nk_\mathrm{B}}}=`$ $`{\displaystyle \frac{2k_\mathrm{B}T}{3J}}\left\{1+{\displaystyle \frac{3}{(2)^3}}+𝒪\left[{\displaystyle \frac{1}{(2)^4}}\right]\right\}`$ (115)
$`+`$ $`{\displaystyle \frac{2(3^{5/2})}{5\pi }}\left({\displaystyle \frac{k_\mathrm{B}T}{J}}\right)^3\left[1+𝒪\left({\displaystyle \frac{1}{2}}\right)\right],`$ (117)
where the prefactor 3/8 in the logarithmic correction term was found independently by Klümper, confirming
Refs. and . The difference between $`C_{\mathrm{lt},\mathrm{log}}(T)`$ and $`C_{\mathrm{lt},\mathrm{g}}(T)`$ is plotted vs temperature in Fig. 12, where the difference becomes $`>10^{10}`$ only for $`t10^5`$.
Shown in Fig. 13 is the deviation $`\mathrm{\Delta }C/Nk_\mathrm{B}`$ ($``$) of the Bethe ansatz data from Lukyanov’s theoretical prediction in Eq. (114). For temperatures $`t10^4`$, the agreement is better than $`10^8`$. At higher temperatures, the uncertainty in the theoretical prediction due to the unknown $`𝒪(g^4)`$ and higher order correction terms becomes an important factor in the comparison. The length of the error bar on each data point in Fig. 13 has arbitrarily been set to $`(4/3)tg^4(t)`$ \[cf. Eq. (114)\]; the $`𝒪(g)`$ uncertainty in the $`T^3`$ term is negligible compared to this. Also plotted in Fig. 13 is the deviation of the numerical data from the extrapolated linear low-$`T`$ behavior ($``$). A comparison of the two data sets indicates that the $`𝒪(g^3)`$ logarithmic correction term is responsible for at least most of this latter difference for temperatures $`t0.001`$.
A more rigorous evaluation of the influence of the above logarithmic correction term is obtained by correcting for it in the plot of $`\mathrm{\Delta }C/t`$ vs $`t`$, as shown by the second curve from the top in each of Figs. 11(a) and 11(b). From the latter figure, we infer that although subtracting this correction term from the data helps to remove the zero-temperature singularity, a singularity is still present but with reduced amplitude. This means that additional logarithmic correction terms are important, within the accuracy and precision of the data. Another indication of this is shown in Fig. 14, where we have plotted $`\mathrm{\Delta }C/t^3`$ vs $`t`$. According to Eq. (114), after accounting for the logarithmic correction term(s), the result should be independent of $`t`$ at low $`t`$. Instead, both before and after accounting for the log correction term, there is a strong
upturn at low temperatures although the strength of the upturn is smaller after subtracting the influence of the log correction term.
The numerical Bethe ansatz specific heat data are sufficiently accurate and precise that we can estimate the coefficients of the next two logarithmic correction ($`g^4,g^5`$) terms in Eq. (114) from these data as follows. From Eq. (114), if we plot the numerical data as $`[C(t)/(Nk_\mathrm{B}t)(2/3)(1+3g^3/8)]/g^4`$ vs $`g`$ at low temperatures, where the $`t^3`$ term can be neglected, and fit the lowest $`t`$ data by a straight line, the $`y`$ intercept for $`g0`$ gives the coefficient of the $`g^4`$ term and the slope gives the coefficient of the $`g^5`$ term. This plot is given in Fig. 15. This type of plot places extreme demands on the accuracy of the data. Even so, we see that the data follow the required linear behavior even at the lowest temperatures. We fitted a straight line to the data from $`t=5\times 10^{25}`$ up to a maximum temperature $`t^{\mathrm{max}}`$. The fit parameters and rms deviation held nearly constant for $`t^{\mathrm{max}}=5\times 10^{15}`$ (11 data points) up to $`t^{\mathrm{max}}=5\times 10^8`$ (18 data points), but both quantities changed rapidly upon further increasing $`t^{\mathrm{max}}`$. The fit for $`t^{\mathrm{max}}=5\times 10^8`$ is shown as the straight line in Fig. 15, along with the fit parameters. From the parameters of the fit \[after accounting for the prefactor of 2/3 in Eq. (114)\], we include our estimated coefficients in Eq. (114) explicitly as
$`{\displaystyle \frac{C_{\mathrm{lt},\mathrm{g}}(T)}{Nk_\mathrm{B}}}={\displaystyle \frac{2k_\mathrm{B}T}{3J}}[1`$ $`+`$ $`{\displaystyle \frac{3}{8}}g^3+a_4g^4+a_5g^5+𝒪(g^6)]`$ (119)
$`+`$ $`{\displaystyle \frac{2(3^{5/2})}{5\pi }}\left({\displaystyle \frac{k_\mathrm{B}T}{J}}\right)^3[1+𝒪(g)],`$ (121)
$$a_4=1.5374(3),a_5=3.125(11).$$
(122)
The influences of these $`g^4`$ and $`g^5`$ logarithmic correction terms on the data in Figs. 11 and 14 are shown as the two additional data sets in each figure, where accounting for these two terms is seen to largely remove the remaining singular behavior as $`t0`$. From Fig. 14, we can now estimate that the coefficient of the $`t^3`$ term in Eq. (14) is a little larger than 2, contrary to the exact value 1.98478$`\mathrm{}`$. The magnitude of this difference is about as expected from the $`𝒪(g)`$ logarithmic correction to the $`t^3`$ term, since $`g(t0.1)0.1`$. The remaining upturn at low temperatures in Fig. 14 is due to residual logarithmic corrections which are not accounted for.
If the approximate expansion for $`g()`$ in Eq. (87) is inserted into Eq. (121), one obtains
$`{\displaystyle \frac{C_{\mathrm{lt},\mathrm{g}}}{Nk_\mathrm{B}}}`$ $`=`$ $`{\displaystyle \frac{2k_\mathrm{B}T}{3J}}\{1+{\displaystyle \frac{3}{(2)^3}}{\displaystyle \frac{9\mathrm{ln}()16a_4}{(2)^4}}`$ (123)
$`+`$ $`{\displaystyle \frac{\mathrm{ln}[18\mathrm{ln}()64a_49]+32a_5}{(2)^5}}+𝒪\left[{\displaystyle \frac{1}{(2)^6}}\right]\}`$ (124)
$`+`$ $`{\displaystyle \frac{2(3^{5/2})}{5\pi }}\left({\displaystyle \frac{k_\mathrm{B}T}{J}}\right)^3\left[1+𝒪\left({\displaystyle \frac{1}{2}}\right)\right].`$ (126)
## IV Fits to $`𝝌^{\mathbf{}}\mathbf{(}𝒕\mathbf{)}`$ and $`𝑪\mathbf{(}𝒕\mathbf{)}`$ of Heisenberg Spin Lattices
### A General $`𝝌^{\mathbf{}}\mathbf{(}𝒕\mathbf{)}`$ fit considerations
The general expression we use to fit theoretical numerical $`\chi ^{}(t)`$ data for $`S=1/2`$ Heisenberg spin lattices is
$$\chi ^{}(t)=\frac{\mathrm{e}^{\mathrm{\Delta }_{\mathrm{fit}}^{}/t}}{4t}𝒫_{(r)}^{(q)}(t),$$
(128)
$$𝒫_{(r)}^{(q)}(t)=\frac{1+_{n=1}^qN_n/t^n}{1+_{n=1}^rD_n/t^n},$$
(129)
where the orders $`q`$ and $`r`$ of the Padé approximant $`𝒫_{(r)}^{(q)}`$ are often constrained by the behavior of $`\chi ^{}(t)`$ at low $`t`$, and the fitted gap $`\mathrm{\Delta }_{\mathrm{fit}}^{}`$ is not necessarily the same as the true gap. At high $`t`$, $`\chi ^{}(t)`$ in Eqs. (IV A) approaches the Curie law $`1/(4t)`$ as required \[for a general spin $`S`$ lattice, the numerical prefactor 1/4 in Eq. (128) would be replaced by $`S(S+1)/3`$\].
The $`N_n`$ and $`D_n`$ parameters in Eq. (129) are not in general independent if one or more of the HTSE conditions in Eqs. (7) and (8) are invoked. For example, for $`n=1`$–3 one finds
$$D_1=(d_1+N_1)\mathrm{\Delta }_{\mathrm{fit}}^{},$$
(131)
$$D_2=(d_2+d_1N_1+N_2)\mathrm{\Delta }_{\mathrm{fit}}^{}(d_1+N_1)+\frac{\mathrm{\Delta }_{\mathrm{fit}}^{}{}_{}{}^{2}}{2},$$
(132)
$`D_3=(d_3`$ $`+`$ $`d_2N_1+d_1N_2+N_3)\mathrm{\Delta }^{}_{\mathrm{fit}}(d_2+d_1N_1+N_2)`$ (133)
$`+`$ $`{\displaystyle \frac{\mathrm{\Delta }_{\mathrm{fit}}^{}{}_{}{}^{2}}{2}}(d_1+N_1){\displaystyle \frac{\mathrm{\Delta }_{\mathrm{fit}}^{}{}_{}{}^{3}}{6}}.`$ (134)
In general, one has
$$D_n=\underset{p=0}{\overset{n}{}}\underset{m=0}{\overset{np}{}}\frac{(\mathrm{\Delta }_{\mathrm{fit}}^{})^p}{p!}d_mN_{npm}.$$
(135)
A fit of experimental or theoretical $`\chi ^{}(t)`$ data by Eqs. (IV A) can be constrained by inserting one or more of Eqs. (IV A) and (135) into Eq. (129). These constraints are especially useful for high-$`t`$ extrapolations when $`\chi ^{}(t)`$ data are not available for high temperatures $`t1`$, and/or to reduce the number of fitting parameters required to obtain a fit of specified precision. In the following fits to the numerical $`\chi ^{}(t)`$ data for the dimer, the uniform chain, and finally our QMC and TMRG data for the alternating-exchange chain, the three constraints in Eqs. (IV A) on $`D_1,D_2`$, and $`D_3`$, respectively, are enforced in each case, where $`d_1,d_2`$, and $`d_3`$ for the alternating-exchange chain are given in Eq. (38).
All of the fits reported in this paper were carried out on a 400 MHz Macintosh G3 (B&W) computer with 1GB of RAM. Most fits were implemented using the program Mathematica 3.0, although a few of the simpler ones (fits to experimental data) were done using KaleidaGraph 3.08c. The fits using Mathematica sometimes required prodigious amounts of memory, e.g., 930 MB for the 28-parameter fit to the combined 2551 data point QMC and TMRG $`\chi ^{}(\alpha ,t)`$ data set for the alternating-exchange chain in Sec. V below.
### B Fit to $`𝝌^{\mathbf{}}\mathbf{(}𝒕\mathbf{)}`$ of the $`𝑺\mathbf{=}\mathrm{𝟏}\mathbf{/}\mathrm{𝟐}`$ antiferromagnetic Heisenberg dimer
The spin gap of the $`S=1/2`$ Heisenberg dimer is $`\mathrm{\Delta }=J`$, where $`J`$ is the antiferromagnetic exchange constant within the dimer. The spin susceptibility and its low-temperature limit are given by Eqs. (12). The $`\chi ^{}(t)`$ is plotted in Fig. 1(a) for $`0.02t4.99`$. In order to later obtain a continuous fit function for $`\chi ^{}(\alpha ,t)`$ for the entire range $`0\alpha 1`$ of the alternating-exchange chain, it is necessary to first obtain a high accuracy fit to the exact expression (14) for the dimer by our general fitting function in Eqs. (IV A), in addition to Fit 1 obtained for the uniform chain below. The form of our fit function in Eqs. (IV A) allows both the low- and high-$`t`$ limiting forms of $`\chi ^{}(t)`$ for the dimer to be exactly reproduced. The low-$`t`$ limit in Eq. (15) requires that $`r=q`$ and that $`D_q=N_q/4`$ in the Padé approximant $`𝒫_{(r)}^{(q)}`$; we also take $`\mathrm{\Delta }_{\mathrm{fit}}=\mathrm{\Delta }`$, so the total number of fitting parameters is $`2q4`$.
We fitted the 498-point double-precision representation of $`\chi ^{}(t)`$ in Fig. 1(a) from $`t=0.02`$ to $`t=4.99`$ by Eqs. (IV A) using the above constraints. The variances of the four fits for $`q=r=4,`$ 5, 6 and 7 were $`2.5\times 10^{13}`$, $`1.17\times 10^{16}`$, $`5.3\times 10^{17}`$, and $`5.6\times 10^{19}`$, respectively, showing that Eqs. (IV A) have the potential for very high accuracy fits with a relatively small number of fitting parameters. The six $`N_n`$ ($`n=1`$–5) and $`D_4`$ parameters of the fit for $`q,r=5`$ are given in Table I, along with $`D_1,D_2`$, and $`D_3`$ computed from Eqs. (IV A) and $`D_5=N_5/4`$. The Padé approximant $`𝒫_{(5)}^{(5)}`$ in the fit function has no poles or zeros on the positive $`t`$ axis. The fit is shown by the solid curve in Fig. 1(a), and the deviation of the fit from the exact susceptibility in Eq. (14) is plotted versus $`t`$ in Fig. 1(b).
### C Fits to $`𝝌^{\mathbf{}}\mathbf{(}𝒕\mathbf{)}`$ of the $`𝑺\mathbf{=}\mathrm{𝟏}\mathbf{/}\mathrm{𝟐}`$ Antiferromagnetic Uniform Heisenberg Chain
##### Fit 1: $`0.01t5`$.
Fits to the uniform chain $`\chi ^{}(t)`$ calculated by Eggert, Affleck and Takahashi for limited temperature regions were obtained previously. Here we obtain a fit (Fit 1) to the higher accuracy data of Klümper and Johnston for the temperature region $`0.01t5`$ (999 data points) using Eqs. (IV A), the results of which will be utilized later in the fit function for $`t0.01`$ for our QMC and TMRG alternating-exchange chain $`\chi ^{}(\alpha ,t)`$ data. This uniform chain fit can be accurately extrapolated to arbitrarily high $`t`$.
The requirement that $`\chi ^{}(t0)`$ is a finite non-zero value requires $`\mathrm{\Delta }_{\mathrm{fit}}^{}=0`$ and $`r=q+1`$ in Eqs. (IV A). We found that using $`q=5`$ and $`r=6`$ produces a fit sufficiently accurate for use in the fit function for our QMC and TMRG calculations for the alternating chain. The seven $`N_n`$ ($`n=1`$–5) and $`D_n`$ ($`n=4`$–6) parameters obtained for the fit with $`q=5,r=6`$ are given in the column labeled “Fit 1” in Table I, along with $`D_1`$, $`D_2`$,
and $`D_3`$ computed from Eqs. (IV A). The Padé approximant $`𝒫_{(6)}^{(5)}`$ in the fit function has no poles or zeros on the positive $`t`$ axis. The deviation of the fit from the data is plotted in Fig. 16. The variance of the fit is $`2.97\times 10^{12}`$, and the relative rms deviation of the fit from the data in the fitted $`t`$ region is 14.5 ppm. Extrapolation of the fit to higher temperatures is very accurate.
The quality of Fit 1 does not approach the limitation imposed by the absolute accuracy of the data ($`1\times 10^9`$). For an ideal fit, the variance is expected to be $`10^{18}`$ and the relative rms deviation $`0.01`$ ppm. As can be inferred from Fit 2 in the following section, the reason that Fit 1 cannot be optimized to this extent is due to the $`t=0`$ critical point and associated logarithmic divergence in the slope of $`\chi ^{}(t)`$ as $`t0`$; this divergence cannot be fitted accurately by a finite polynomial or Padé approximant. We attempted to improve the accuracy of the fit over the same temperature range $`0.01t5`$ by replacing the Padé approximant $`𝒫_{(6)}^{(5)}`$ in the fit function by $`𝒫_{(7)}^{(6)}`$, which incorporates two additional fitting parameters. The variance improved somewhat to $`2.18\times 10^{12}`$ and the relative rms deviation improved slightly to 12.2 ppm, but the Padé approximant developed a pole at $`1/t=129.23`$, and hence this fit was discarded. Although the temperature at which this pole occurs is below the fitted temperature range, as a general rule we cannot allow poles in the fit function at low temperatures because of problems that can occur when using the fit function to model experimental data which include data at temperatures lying below the fitted temperature range of the fit function. In fact, we will encounter this situation frequently in modeling experimental data later. For exam-
ple, for $`\mathrm{NaV}_2\mathrm{O}_5`$, $`t=0.01`$ corresponds to an absolute temperature $`T7`$ K, whereas the experimental data and modeling extend down to $`2`$ K.
##### Fit 2: $`0t5`$.
We can greatly improve the accuracy of the fit compared to that of Fit 1, and extend the fit to $`t=0`$, by restricting the high-temperature limit of the fit and using in the fit function one or more low-temperature logarithmic correction terms discussed in Sec. III A 2. In particular, in this section we obtain a very high precision fit (Fit 2) to the exact $`t=0`$ value and to the calculations of Klümper and Johnston over the entire temperature range $`5\times 10^{25}t5`$ of the calculations. We do not use this fit in our formulation of the fit function for the alternating-exchange chain. However, Fit 2 will be generally useful for evaluating the accuracy of other theoretical calculations of $`\chi ^{}(t)`$ for the uniform chain, such as our TMRG calculations to be presented below, and for modeling appropriate experimental $`\chi (T)`$ data whose scaled upper temperature limit is below $`t=5`$.
We initially formulated a fit function utilizing a modified Padé approximant in which the last term of the numerator and/or denominator contained the $`\chi _{\mathrm{lt},\mathrm{log}}^{}`$ expansion in Eq. (90), such that the low-temperature expansion of the fit function yielded $`\chi _{\mathrm{lt},\mathrm{log}}^{}`$ to lowest orders in $`t`$. The best fit to the data from $`t=5\times 10^{25}`$ to 2.5 (777 data points) was unsatisfactory, with a variance $`v=2.4\times 10^{11}`$ and a relative rms deviation $`\sigma _{\mathrm{rms}}=45`$ ppm. Allowing an arbitrary $`t^2`$ coefficient in place of the exact value $`\sqrt{3}/\pi ^3`$ yielded an improved fit with $`v=1.1\times 10^{12}`$ and $`\sigma _{\mathrm{rms}}=9.6`$ ppm. However, this fit was still unsatisfactory, given the high absolute accuracy ($`1\times 10^9`$) of the data. From these results it became clear that a fit function which can fit the data to much higher accuracy over such a large temperature range would indeed have to include an expression $`\chi _{\mathrm{log}}^{}(t)`$ containing logarithmic correction terms, but where the form and/or coefficient of one or more of these terms would have to be empirically determined by trial and error. This process yielded the formulation we now describe.
The $`\chi _{\mathrm{log}}^{}(t)`$ function is incorporated into our fit function in Eqs. (IV A) as follows. As in Fit 1, the finite value of $`\chi ^{}(0)`$ requires $`\mathrm{\Delta }_{\mathrm{fit}}=0`$ in Eq. (128) and $`r=q+1`$ in the Padé approximant $`𝒫_{(r)}^{(q)}(t)`$ in Eq. (129). Since the two terms highest order in $`1/t`$ in $`𝒫_{(r)}^{(q)}(t)`$ (one each in the numerator and denominator) dominate the fit as $`t0`$ and become small for $`t1`$, relative to the other terms in the numerator and denominator, respectively, we incorporate $`\chi _{\mathrm{log}}^{}(t)`$ into the last term in the numerator of a modified $`𝒫_{(r)}^{(q)}(t)`$. Trial fits showed that to obtain the optimum accuracy of the fit required $`q=8`$ and $`r=9`$.
Our final fit function for Fit 2 is
$$\chi ^{}(t)=\left(\frac{1}{4t}\right)\frac{1+\left[_{n=1}^7N_n/t^n\right]+4N_8\chi _{\mathrm{log}}^{}(t)/t^8}{1+\left[_{n=1}^8D_n/t^n\right]+N_8/t^9},$$
(137)
$$\chi _{\mathrm{log}}^{}(t)=\frac{1}{\pi ^2}\left[1+\frac{1}{2}\frac{\mathrm{ln}\left(+\frac{1}{2}\right)N_{81}}{(2)^2}+\frac{N_{82}}{(2)^3}\right],$$
(138)
$$\mathrm{ln}(t_1/t),$$
(139)
subject to the three constraints on $`D_1,D_2`$, and $`D_3`$ in Eqs. (IV A) which are required by the HTSE. Two of the four logarithmic correction terms in Eq. (138) are identical to the first two such terms in Eq. (78). By construction, the exact $`\chi ^{}(0)=1/\pi ^2`$ is fitted exactly.
We fitted all of the numerical $`\chi ^{}(t)`$ data, calculated over the range $`5\times 10^{25}t5`$ (1119 data points), by Eqs. (IV C). The 19 fitting parameters of the fit function (IV C), which are $`N_n`$ ($`n=1`$–8), $`D_n`$ ($`n=4`$–8), $`N_{81},N_{82}`$ and $`t_1`$, are given in the column labeled “Fit 2” in Table I, along with $`D_1,D_2`$, and $`D_3`$ computed from Eqs. (IV A). The data to parameter ratio is 59. The denominator of the modified Padé approximant in Eq. (137) has no zeros for any real positive $`t`$. The fit is shown in the low-temperature region $`0t0.02`$ in Fig. 4(b) \[over the larger $`t`$ range plotted in Fig. 4(a), the fit is indistinguishable from the data and is therefore not plotted there\].
The deviation of Fit 2 from the numerical data for $`10^{25}t5`$ is plotted vs $`\mathrm{log}_{10}t`$ in Fig. 17(a), and an expanded plot at the higher temperatures is shown in Fig. 17(b). Due to a logarithmic divergence in $`\chi _{\mathrm{log}}^{}(t)`$ at $`t=t_1=5.696`$, Fit 2 should not be used (e.g., for modeling experimental data) at temperatures $`t5`$. The variance of the fit is $`9.8\times 10^{17}`$, and the relative rms deviation is $`\sigma _{\mathrm{rms}}=0.087`$ ppm. These values are both much smaller than for Fit 1 above. The relatively large number of fitting parameters in Fit 2 is justified a posteriori by the extremely high quality of the fit over a temperature range spanning 25 orders of magnitude.
### D Fit to $`𝑪\mathbf{(}𝒕\mathbf{)}`$ for the $`𝑺\mathbf{=}\mathrm{𝟏}\mathbf{/}\mathrm{𝟐}`$ antiferromagnetic uniform Heisenberg chain
The logarithmic corrections to the magnetic specific heat $`C(t)`$ at low temperatures, discussed above in Sec. III B 2, do not pose as serious a problem for fitting the data as for $`\chi ^{}(t)`$, because the strength of these log corrections is much smaller for $`C(t)`$ than for $`\chi ^{}(t)`$. In addition, since here we fit $`C(t)`$, and not the electronic specific heat coefficient $`C(t)/t`$, the influence of the log corrections is ameliorated by the multiplicative leading order $`t^1`$ dependence of $`C(t)`$. Even so, in order to obtain the optimum fit to the highly accurate Bethe ansatz $`C(t)`$ data, we found it necessary to take the influence of the logarithmic corrections into account.
Our fit to the Bethe ansatz $`C(t)`$ data, some of which were shown previously in Fig. 8(a), was carried out in two stages. First, the data from $`t=0.01`$ to the maximum temperature $`t=5`$ of the calculations were fitted by the Padé approximant $`𝒫_{(r)}^{(q)}`$ in Eq. (129) with a prefactor $`3/(16t^2)`$ to satisfy the HTSE in Eqs. (III B 1) to lowest order in $`1/t`$. The orders $`q`$ and $`r`$ of $`𝒫_{(r)}^{(q)}`$ were chosen to satisfy $`r=q+3`$ so that $`C(t0)t`$. To obtain a fit of the required accuracy (see the fit deviations given below) we found that $`q=6`$ and $`r=9`$ are of sufficiently high order. Due to the presence of the log corrections at very low $`t`$, we did not require the parameters $`N_6`$ and $`D_9`$ to yield the exact coefficient $`\gamma =2/3`$ in the expression $`C(t)/Nk_\mathrm{B}=\gamma t`$, in a low-$`t`$ expansion of the fit function. We also found that to obtain the best fit, only the one additional HTSE constraint (on $`D_1`$) in Eq. (131) (with $`\mathrm{\Delta }_{\mathrm{fit}}^{}=0`$) could be used. It was quite difficult to find the region in parameter space in which the absolute minimum in the variance of the fit resided; the initial starting parameters usually flowed to regions with local variance minima in them with much larger values (by two to four orders of magnitude) than the smallest variance we ultimately found. Then the deviation of the fit from all the data for $`5\times 10^{25}t5`$ was computed. The fit deviations for $`t0.01`$ were very small \[$`𝒪(10^8)`$\], but the log corrections which become most important at lower
temperatures resulted in fit deviations at $`t<0.01`$ an order of magnitude larger. We therefore fitted the fit deviation versus $`t`$ for $`0<t0.1`$ by a separate empirically determined function $`F(t)`$, so the net fit function consists of the Padé approximant fit function minus the fit function to the low-$`t`$ fit deviations. In the final fitting cycles the two functions were refined simultaneously.
Our final fit function for $`C(t)`$ in the range $`0t5`$ is
$$\frac{C(t)}{Nk_\mathrm{B}}=\frac{3}{16t^2}𝒫_{(9)}^{(6)}(t)F(t),$$
(141)
$$𝒫_{(9)}^{(6)}(t)=\frac{1+\left[_{n=1}^6N_n/t^n\right]}{1+\left[_{n=1}^9D_n/t^n\right]},$$
(142)
$$F(t)=a_1t\mathrm{sin}\left(\frac{2\pi }{a_2+a_3t}\right)\mathrm{e}^{a_4t}+a_5t\mathrm{e}^{a_6t},$$
(143)
subject to the constraint on $`D_1`$ in Eq. (131) which is required by the HTSE. By construction, the exact $`C(0)=0`$ is fitted exactly. The 20 fitting parameters, $`N_n(n=16),D_n(n=29)`$ and $`a_n(n=16)`$, are given in Table I, together with the constrained parameter $`D_1`$ computed from Eq. (131) with $`\mathrm{\Delta }_{\mathrm{fit}}^{}=0`$ and $`d_1`$ given in Eq. (108). The deviation of the fit from the data is shown in a semilog plot vs temperature in Fig. 18. The maximum deviations of $`\pm 4\times 10^8`$ occur at $`t0.3`$. The absolute rms deviation of the fit from all the data (1119 data points), which extend over the temperature range $`5\times 10^{25}t5`$, is $`1.34\times 10^8`$, and the relative rms deviation for $`0.01t5`$ (999 data points) is 0.50 ppm.
At high temperatures, our $`C(t)`$ fit function reduces by construction to the lowest order $`1/t^2`$ and $`1/t^3`$ terms of the HTSE of $`C(t)`$ in Eqs. (III B 1), so extrapolation of our $`C(t)`$ fit function to arbitrarily higher temperatures should be very accurate (see Fig. 18). In particular, even though our fit was to $`C(t)`$ and hence not optimized as a fit to the electronic specific heat coefficient $`C(t)/t`$, the magnetic entropy $`S`$ at $`t=\mathrm{}`$ computed from our $`C(t)`$ fit function is
$$\frac{S(t=\mathrm{})}{Nk_\mathrm{B}}_0^{\mathrm{}}\frac{C(t)}{Nk_\mathrm{B}t}𝑑t=\mathrm{0.693\hspace{0.17em}147\hspace{0.17em}235},$$
(144)
which is the same as the exact value $`\mathrm{ln}2=\mathrm{0.693\hspace{0.17em}147\hspace{0.17em}181}`$ to within 8 parts in $`10^8`$. This agreement reflects well on our fit function, and of course also strongly confirms the high accuracy of the Bethe ansatz $`C(t)`$ data.
### E Fit function for the $`𝑺\mathbf{=}\mathrm{𝟏}\mathbf{/}\mathrm{𝟐}`$ AF alternating-exchange Heisenberg chain $`𝝌^{\mathbf{}}\mathbf{(}𝜶\mathbf{,}𝒕\mathbf{)}`$
Here we formulate a single two-dimensional ($`\alpha ,t`$) function to accurately fit numerical calculations of $`\chi ^{}(\alpha ,t)`$ for the $`S=1/2`$ alternating-exchange Heisenberg chain for the entire range $`0\alpha 1`$, and for the entire temperature range $`t0.01`$ over which our Fit 1 for $`\chi ^{}(t)`$ of the uniform chain is most accurate, subject to four general requirements as follows. (i) The HTSE of the $`\chi ^{}(\alpha ,t)`$ fit function must give the correct result to $`𝒪(1/t^4)`$, as satisfied by the fit functions for the dimer and uniform chain (Fit 1) susceptibilities above, so that the fit can be accurately extrapolated to higher temperatures. (ii) We require the $`\chi ^{}(\alpha ,t)`$ fit function to become identical with those found above for the isolated dimer and for the uniform chain (Fit 1) when $`\alpha =0`$ and $`\alpha =1`$, respectively. As discussed above in Sec. II C, at any finite temperature, $`\overline{\chi ^{}}(\delta ,\overline{t})`$ in the variables $`\delta `$ and $`\overline{t}`$ is an even (analytic) function of $`\delta `$. Therefore, as a minimum accommodation of this fact, (iii) we require that the fit function for $`\chi ^{}(\alpha ,t)`$, when transformed to the form $`\overline{\chi ^{}}(\delta ,\overline{t})`$, must have the property $`\overline{\chi ^{}}(\delta ,\overline{t})/\delta |_{\delta =0}=0`$ at all finite temperatures. This requirement is clearly the minimum necessary in order to accurately interpolate the fit vs $`\alpha `$ for $`\alpha 1`$ at each $`t`$, and to thereby accurately model the susceptibility of materials which are in or near this limit. Finally, the QMC and TMRG calculations of $`\chi ^{}(\alpha ,t)`$ to be presented below are sufficiently accurate and cover sufficiently large ranges of $`\alpha `$ and $`t`$ with sufficient resolution that (iv) we require the nonanalytic energy gap $`\mathrm{\Delta }(\alpha )`$ \[see Eqs. (II C 2) and (49)\] to be included in the fit function in order to fit the data for $`\alpha 1`$ at $`t1`$, so as to avoid the alternate necessity of including high-order power series in $`\alpha `$ and $`t`$ in the fit function. We note that according to Eq. (48) or (49), $`\overline{\mathrm{\Delta }^{}}(\delta )/\delta |_{\delta =0}=\mathrm{}`$. The major obstacle we faced in formulating the fit function for $`\chi ^{}(\alpha ,t)`$ was to simultaneously satisfy both requirements (iii) and (iv), which at first sight seem to require mutually exclusive forms for the fit function.
We found that these four requirements can all be satisfied by an extension of the form of the fit function in Eqs. (IV A) which was used above for the isolated dimer and for the uniform chain Fit 1. This extension consists of using a modified Padé approximant $`𝒫_{\mathrm{m}}^{}{}_{(8)}{}^{(7)}`$ in the fit function in place of the former $`𝒫_{(r)}^{(q)}`$. The fit function is
$$\chi ^{}(\alpha ,t)=\frac{\mathrm{e}^{\mathrm{\Delta }_{\mathrm{fit}}^{}(\alpha )/t}}{4t}𝒫_{\mathrm{m}}^{}{}_{(8)}{}^{(7)}(\alpha ,t),$$
(146)
$`𝒫_{\mathrm{m}}^{}{}_{(8)}{}^{(7)}(\alpha ,t)=`$
$$\frac{[_{n=0}^6N_n/t^n]+(N_{71}\alpha +N_{72}\alpha ^2)(\mathrm{\Delta }_0^{}/t)^y/t^7}{[_{n=0}^7D_n/t^n]+(D_{81}\alpha +D_{82}\alpha ^2)(\mathrm{\Delta }_0^{}/t)^z\mathrm{e}^{(\mathrm{\Delta }_0^{}\mathrm{\Delta }_{\mathrm{fit}}^{})/t}/t^8},$$
(147)
$$\mathrm{\Delta }_{\mathrm{fit}}^{}(\alpha )=1\frac{1}{2}\alpha 2\alpha ^2+\frac{3}{2}\alpha ^3,$$
(148)
$`\mathrm{\Delta }_0^{}(\alpha )=(1`$ $``$ $`\alpha )^{3/4}(1+\alpha )^{1/4}`$ (149)
$`+`$ $`g_1\alpha (1\alpha )+g_2\alpha ^2(1\alpha )^2,`$ (150)
$$N_0=D_0=1,$$
(151)
$$N_n(\alpha )=\underset{m=0}{\overset{4}{}}N_{nm}\alpha ^m(n=16),$$
(152)
$$D_n(\alpha )=\underset{m=0}{\overset{4}{}}D_{nm}\alpha ^m(n=17).$$
(153)
To satisfy requirement (i), $`D_1(\alpha ),D_2(\alpha )`$, and $`D_3(\alpha )`$ are determined from the $`N_1(\alpha ),N_2(\alpha )`$, and $`N_3(\alpha )`$ fitting parameters according to the three constraints in Eqs. (IV A) demanded by the HTSE. In order to satisfy requirement (ii), the $`\{N_{n0},D_{n0}\}`$ parameters are set to be identical with those determined above for the dimer, and we require $`\{N_n(1),D_n(1)\}`$ to be identical with the corresponding fit parameters determined above in Fit 1 for the uniform chain. In order to satisfy requirement (iii), the $`N_{nm}`$ and $`D_{nm}`$ coefficients must satisfy
$$\underset{m=0}{\overset{4}{}}(n2m)(N_{nm}\mathrm{or}D_{nm})=0,$$
(154)
so that no $`\delta ^1`$ term appears in the Taylor series expansions in $`\delta `$ of the transformed $`\{\overline{N}_n(\delta ),\overline{D}_n(\delta )\}`$. These various constraints on the $`\{N_{nm},D_{nm}\}`$ parameters reduce the number of independent fitting parameters within this set from 50 to 20. Together with the parameters $`N_{71},N_{72},N_{81},N_{82},y,z`$ in Eq. (147) and $`g_1,g_2`$ in Eq. (150), the total number of independent fitting parameters in the fit function is 28.
The quantity $`\mathrm{\Delta }_{\mathrm{fit}}^{}(\alpha )`$ in the exponential prefactor to $`𝒫_{\mathrm{m}}^{}{}_{(8)}{}^{(7)}`$ in Eq. (146) cannot be set equal to the true nonanalytic gap $`\mathrm{\Delta }^{}(\alpha )`$, because this prefactor affects the fit at all $`t`$, and would not allow requirement (iii) above to be fulfilled. In addition, the nonanalytic critical behavior of $`\mathrm{\Delta }^{}(\alpha 1)`$ in practice only becomes manifest in $`\chi ^{}(\alpha ,t)`$ at low temperatures $`t1`$. Therefore, we separated the spin gap into an analytic part $`\mathrm{\Delta }_{\mathrm{fit}}^{}(\alpha )`$ which goes into the argument of the exponential prefactor in Eq. (146), and a nonanalytic part $`\mathrm{\Delta }_0^{}(\alpha )`$ \[satisfying requirement (iv)\] which is placed into the argument of the exponential in the last term of the denominator of $`𝒫_{\mathrm{m}}^{}{}_{(8)}{}^{(7)}`$ in Eq. (147) and which therefore only becomes important at low temperatures. The first two terms of $`\mathrm{\Delta }_{\mathrm{fit}}^{}(\alpha )`$ (to order $`\alpha ^1`$) in Eq. (148) are the first two terms of the exact dimer series expansion up to $`𝒪(\alpha ^9)`$ given by Barnes, Riera and Tennant for the AF alternating-exchange chain, and the last two are included so that $`\overline{\mathrm{\Delta }_{\mathrm{fit}}^{}}(\delta )/\delta |_{\delta =0}=0`$, in accordance with requirement (iii). The nonanalytic $`\mathrm{\Delta }_0^{}(\alpha )`$ in Eq. (150) contains the behavior in Eq. (47) proposed by Barnes, Riera, and Tennant, plus two analytic terms which are included to adjust the $`\alpha `$ dependence for $`\alpha 1`$ but which make no contribution at $`\alpha =0`$ or $`\alpha =1`$. Provided that the inequality $`y,z>4/3`$ is satisfied by the powers $`y`$ and $`z`$ in Eq. (147), the last term in each of the numerator and denominator of $`𝒫_{\mathrm{m}}^{}{}_{(8)}{}^{(7)}(\alpha ,t)`$, when transformed to the variables $`(\delta ,\overline{t})`$, has a partial derivative with respect to $`\delta `$ which is zero at $`\delta =0`$.
We have now shown that at $`\delta =0(\alpha =1)`$, the partial derivative of each part of $`\overline{\chi ^{}}(\delta ,\overline{t})`$ with respect to $`\delta `$ is zero (if $`y,z>4/3`$, which is confirmed in the actual fit later). Hence, the entire fit function has the property $`\overline{\chi ^{}}(\delta ,\overline{t})/\delta |_{\delta =0}=0`$ at all finite temperatures, thus satisfying requirement (iii), despite the fact that the fit function contains the nonanalytic $`\mathrm{\Delta }_0^{}(\alpha )`$ as required by requirement (iv).
At the lowest temperatures, the last term in each of the numerator and denominator of $`𝒫_{\mathrm{m}}^{}{}_{(8)}{}^{(7)}`$ in Eq. (147) should dominate the fit, together with the exponential prefactor to $`𝒫_{\mathrm{m}}^{}{}_{(8)}{}^{(7)}`$ in Eq. (146), so in this limit our fit function for $`0<\alpha <1`$ becomes
$$\chi ^{}(\alpha ,t0)=\frac{N_{71}\alpha +N_{72}\alpha ^2}{4(D_{81}\alpha +D_{82}\alpha ^2)}\left[\frac{\mathrm{\Delta }_0^{}(\alpha )}{t}\right]^{yz}\mathrm{e}^{\mathrm{\Delta }_0^{}(\alpha )/t}.$$
(155)
This expression has the form of Eq. (18) (with $`\gamma =yz`$) as required in the low-$`t`$ limit. In fact, the forms of the last term in each of the numerator and denominator of $`𝒫_{\mathrm{m}}^{}{}_{(8)}{}^{(7)}`$ were designed to result in the form of Eq. (18) in the low-$`t`$ limit, with $`\mathrm{\Delta }_0^{}`$ and $`t`$ entering the prefactor only as their ratio as in Eq. (57), in addition to being consistent with requirements (iii) and (iv). One might expect the fitted $`y`$ and $`z`$ powers to satisfy $`yz=\gamma =1/2`$ as in Eq. (19). However, if a fit of $`\chi ^{}(t)`$ data by Eq. (18) is not carried out completely within the low-$`t`$ limit, an effective exponent $`\gamma 1`$ is often inferred \[see, e.g., Eq. (181) and subsequent discussion, and Fig. 36 below\]. Similarly, since many of our calculated $`\chi ^{}(\alpha ,t)`$ data sets for different $`\alpha `$ in the fitted temperature range $`t0.01`$ are not, or do not contain extensive data, in the low-$`t`$ limit, we did not impose the constraint $`yz=1/2`$. On the basis of the above discussion we expect the actual fitted values of $`y`$ and $`z`$ to yield $`yz1`$. In fact, as will be seen in the next section, our fitted parameters $`y`$ and $`z`$ give $`yz=1.14`$.
## V QMC and TMRG $`𝝌^{\mathbf{}}\mathbf{(}𝜶\mathbf{,}𝒕\mathbf{)}`$ Calculations and Fit for the $`𝑺\mathbf{=}\mathrm{𝟏}\mathbf{/}\mathrm{𝟐}`$ AF Alternating-Exchange Heisenberg Chain
QMC simulations of $`\chi ^{}(\alpha ,t)`$ were carried out on $`S=1/2`$ alternating-exchange chains containing 100 spins for $`\alpha =0.05,0.1,0.15,\mathrm{},`$ 0.9, 0.92, 0.94, 0.96, 0.97, 0.98, and 0.99 in various temperature ranges spanned by $`0.01t4`$.
Complementary TMRG calculations of $`\chi ^{}(\alpha ,t)`$ of $`S=1/2`$ alternating-exchange chains were carried out for $`\alpha =0.80`$, 0.82, …, 0.96, 0.97, 0.98, 0.99, 0.995 and 1, where the number of states kept was $`m=150`$ or 256. The calculations were carried out for temperatures given by $`1/t=0.1,0.2,\mathrm{},(1/t)^{\mathrm{max}}`$, with $`(1/t)^{\mathrm{max}}500`$ increasing with increasing $`\alpha `$, and comprised a total of 22 370 ($`\alpha ,t`$) parameter combinations. The details of the calculational method are given in Refs. and . It should be noted that the TMRG calculations by their nature are explicitly in the thermodynamic limit.
The reason for doing TMRG calculations for the uniform chain ($`\alpha =1`$) was to enable comparison of the results with the values computed with the Bethe ansatz which have a high absolute accuracy of $`1\times 10^9`$. This comparison was done using the above very accurate and precise Fit 2 for the Bethe ansatz data. The relative deviation of the TMRG data from Fit 2 is shown in Fig. 19(a), and an expanded plot for the higher temperature region $`t0.01`$ is shown in Fig. 19(b). This comparison indicates that the accuracy of the TMRG calculations for both $`m=150`$ and 256 in the range $`t0.01`$ is better than 0.1 %, which is the same as the estimate made previously for $`m=80`$. However, the accuracy of these calculations deteriorates rapidly at lower $`t`$, to about 3 % at the lowest temperatures $`t0.002`$ for $`m=150`$.
Since the TMRG calculations extend close to the $`t=0`$ limit for most of the above-stated $`\alpha `$ values, the spin gaps can be estimated from these data. Comparisons with previous work can then be made of the dependence of the spin gap on $`\alpha `$. An important question, not yet
answered in previous work, is the approximate $`\alpha `$ value at which the asymptotic critical region is entered upon approaching the uniform limit. Performing these estimates and comparisons will be postponed to the following sections. In the present section, we present the QMC and TMRG $`\chi ^{}(\alpha ,t)`$ data and obtain a fit to these combined data by the fit function formulated in the previous section.
Some of the results for $`t2`$ are shown as the filled symbols without error bars in Fig. 20(a) (the error bars are smaller than the data symbols); an expanded plot of data for $`t0.4`$ is shown in Fig. 20(b). \[A log-log plot of the TMRG $`\chi ^{}(\alpha ,t)`$ data at low $`t`$ is shown below in Fig. 27.\] Also shown in both figures as the two bounding solid curves with no data points are the fits we obtained above to $`\chi ^{}(t)`$ for the dimer and uniform chain (Fit 1), respectively. The data points plotted for a given $`\alpha `$ value are the subset below the upper temperature limits of the figures, of the subset of available data points which were fitted by our fit function as described below.
We fitted a combined QMC and TMRG $`\chi ^{}(\alpha ,t)`$ data set containing 2551 selected data points over the temperature range $`0.01t10`$. The 802 QMC data points covered the ranges $`0.01t4`$ and $`0.05\alpha 0.99`$.
The average estimated absolute accuracy of these QMC data is $`1.7\times 10^4`$. The best estimated accuracy among these QMC data is $`7.7\times 10^6`$ and the worst is $`1.5\times 10^3`$, with the better accuracies occurring at the highest temperatures. The 1749 TMRG data points covered the ranges $`0.01t10`$ and $`0.8\alpha 0.995`$. We did not use all 22 370 TMRG data points in the available data set, because this would have weighted the region $`\alpha 1`$ too heavily in the fit, and in any case a large fraction of these are for temperatures below our low-temperature fitting limit of $`t=0.01`$. We used the low-temperature data to determine the spin gaps as described in the following section.
We fitted this $`\chi ^{}(\alpha ,t)`$ data set by Eqs. (IV E), with the constraints on the parameters discussed above. Obtaining a reliable 28-parameter two-dimensional fit to these data over the full above-cited ranges of $`t`$ and $`\alpha `$, with no poles in the fit, posed a very difficult challenge. The particular choice of starting parameters and the detailed sequence of refinements were found to be important to avoiding poles in the final fit. Since there are a total of 28 parameters in the fit function for 2551 data points, the data to parameter ratio is 91. The number of fitting parameters seems large, until it is realized that we are simultaneously fitting $`\chi ^{}(t)`$ data for 29 different $`\alpha `$ values, so on average a $`\chi ^{}(t)`$ data set for a given $`\alpha `$ value is fitted by a single parameter. A weighting function was not included during the variance minimization, because we were interested in obtaining a fit which treated all the data points the same on an absolute scale; this choice optimizes the fit for use in modeling experimental data.
The parameters of the fit are given in Table II, where we have also included the constrained parameters for completeness and for ease of implementation of our fit function by the reader. From Eqs. (IV A), the constrained parameters $`D_2`$ and $`D_3`$ contain products of the third-order (in $`\alpha `$) polynomial $`\mathrm{\Delta }_{\mathrm{fit}}^{}`$ with itself and/or with the fourth-order $`N_1`$ fitting polynomial, so $`D_2`$ and $`D_3`$ are of seventh and tenth-order, respectively. The two-dimensional fit is shown as the set of solid curves in Fig. 20. The variance of the fit is $`v=3.77\times 10^8`$. The absolute rms deviation $`\sqrt{v}1.9\times 10^4`$ is about the same as the average estimated accuracy of the QMC data noted above, indicating that the fit function is appropriate and that the fit is a reliable representation of
the data. The fit deviations from the 802 QMC and 1749 TMRG data are shown separately in Figs. 21(a) and 21(b), respectively. A comparison of the two figures shows that the TMRG data are, on average, significantly more precise at a given temperature.
After the parameters in the present $`\chi ^{}(\alpha ,t)`$ fit function were finalized, as a check on the accuracy of the fit function for $`\alpha `$ values close to the uniform limit, we carried out QMC $`\chi ^{}(t)`$ simulations for alternating-exchange chains of length $`L=400`$ and 800, factors of four and eight longer than the chains for which QMC data were combined with TMRG data to determine the fit function, respectively. The simulations were carried out for $`\alpha =0.98,0.985,0.99`$, and 0.995 at temperatures $`0.01t4`$. Overall, the fit function was found to be in extremely good agreement with the QMC data. For $`0.4t4`$, the $`\chi ^{}(\alpha ,t)`$ fit function agreed with the simulation data to within about $`\pm 5\times 10^5`$ or better. The deviations of the fit function from the data for $`0.01t0.4`$ are shown in Fig. 22, along with the error bars on the QMC data. As can be seen from the figure, the only significant deviation of the fit function from the QMC data in this $`t`$ range is at the lowest temperature $`t=0.01`$ for each of the four $`\alpha `$
values. Because the fit deviations at this temperature remain upon increasing the length of the simulated chain from $`L=400`$ to $`L=800`$, these fit deviations are most likely due to inaccuracies in the fit function, as expected at this lowest fitted temperature.
For compounds showing spin-Peierls or other types of second-order spin dimerization transitions, it is more appropriate to scale $`\chi `$ by $`1/J`$ and $`T`$ by $`J`$, where $`J`$ is the average of $`J_1`$ and $`J_2`$, in which case the appropriate alternation parameter is $`\delta `$ rather than $`\alpha `$. It is straightforward to convert our $`\chi ^{}(\alpha ,t)`$ fit function to the form $`\overline{\chi ^{}}(\delta ,\overline{t})`$, where $`\overline{t}k_\mathrm{B}T/J`$, using Eq. (41). We have done this and plot the $`\overline{\chi ^{}}(\delta ,\overline{t})`$ fit function versus temperature for a series of $`\delta `$ values in Fig. 23(a). An appealing monotonic progression of $`\overline{\chi ^{}}(\delta ,\overline{t})`$ with increasing $`\delta `$ is seen in Fig. 23(a); an expanded plot at lower temperatures is shown in Fig. 23(b). This formulation of the fit function allows accurate estimates to be made of the temperature-dependent spin gap in compounds exhibiting spin-dimerization transitions, provided that the nearest neighbor $`S=1/2`$ AF alternating-exchange Heisenberg model is appropriate to them. An illustration of the procedure and the results to be gained will be given later when we model the $`\chi (T)`$ data for NaV<sub>2</sub>O<sub>5</sub>.
## VI Spin Gap from TMRG $`𝝌^{\mathbf{}}\mathbf{(}𝜶\mathbf{,}𝒕\mathbf{)}`$
According to Eq. (23), if highly precise $`\chi ^{}(t)`$ data in the low-$`t`$ limit are available, the spin gap $`\mathrm{\Delta }^{}`$ can in principle be computed directly from the derivative of these data with respect to inverse temperature. However, in general the maximum temperature of the low-$`t`$ limit region is ill defined since it depends on how precise and accurate the data are and the accuracy to which $`\mathrm{\Delta }^{}`$ is to be determined. Therefore, in practice one could define a temperature-dependent effective spin gap $`\mathrm{\Delta }_{\mathrm{eff}}^{}`$ from Eq. (23) as
$$\mathrm{\Delta }_{\mathrm{eff}}^{}(t)=\frac{\mathrm{ln}(\chi ^{}\sqrt{t})}{(1/t)},$$
(156)
and then try to ascertain $`\mathrm{\Delta }^{}`$ from the extrapolated zero-temperature limit $`\mathrm{\Delta }^{}=lim_{t0}\mathrm{\Delta }_{\mathrm{eff}}^{}(t)`$. Using Eq. (22) would be less desirable and precise because a fit of this type typically averages $`\mathrm{\Delta }_{\mathrm{eff}}^{}(t)`$ over a rather large temperature range.
An overview of $`\mathrm{\Delta }_{\mathrm{eff}}^{}(\alpha ,t)`$ determined from our TMRG $`\chi ^{}(\alpha ,t)`$ data for $`0.8\alpha 0.995`$ using Eq. (156) is shown in Fig. 24(a). At the lowest temperatures, and for $`\alpha `$ not too close to 1, the $`\mathrm{\Delta }_{\mathrm{eff}}^{}(\alpha ,t)`$ data do approach a constant value with decreasing $`t`$, confirming the applicability of Eqs. (16) and prior assumptions and hence Eqs. (1) and (156) to the alternating chain, and the approximate values of $`\mathrm{\Delta }^{}(\alpha )`$ can be estimated from the figure. Closer inspection reveals that $`\mathrm{\Delta }_{\mathrm{eff}}^{}(\alpha ,t)`$ shows a weak maximum before decreasing by $`\frac{1}{2}`$% to $`\mathrm{\Delta }^{}`$ as $`t0`$, as illustrated in Fig. 24(b) for $`\alpha =0.8`$. For this among other reasons, we will not use Eq. (156) to extract the spin gaps from our TMRG $`\chi ^{}(\alpha ,t)`$ data. On the other hand, we need to know whether such behavior is expected, since it could conceivably arise from systematic errors in the TMRG calculations. Therefore, in the next section we study the $`\mathrm{\Delta }_{\mathrm{eff}}^{}(\alpha ,t)`$ expected at low temperatures for the alternating-exchange chain. As part of this study, we formulate and discuss the fit function which we will use in Sec. VI B to extract $`\mathrm{\Delta }^{}(\alpha )`$ from our TMRG $`\chi ^{}(\alpha ,t)`$ data at low temperatures.
### A Effective spin gap $`𝚫_{\mathrm{𝐞𝐟𝐟}}^{\mathbf{}}\mathbf{(}𝚫^{\mathbf{}}\mathbf{,}𝒕\mathbf{)}`$ for the alternating-exchange chain
From our definition of $`\mathrm{\Delta }_{\mathrm{eff}}^{}`$ in Eq. (156), a discussion of how this quantity varies with $`t`$ at low $`t`$ requires an independent estimate of $`\chi ^{}(\alpha ,t)`$ for the alternating-exchange
chain, which must include at least the leading order correction to the low-$`t`$ limit in Eqs. (16). As a first attempt, we used the general expression for $`\chi ^{}(t)`$ in Eqs. (II B), which requires as input the one-magnon dispersion relation $`\epsilon (k)`$ for the alternating chain. For this we used the explicit $`\epsilon (\mathrm{\Delta }^{},k)`$ for $`0\alpha 1`$ in Eqs. (3) that we presented and discussed previously. The resultant $`\chi ^{}(t)`$ is plotted for eleven $`\mathrm{\Delta }^{}`$ values in Fig. 25(a), where the results are designated by $`\chi ^{(1)}`$ in the figure. Although $`\chi ^{(1)}(\mathrm{\Delta }^{},t)`$ is exact in both the low- and high-$`t`$ limits, the results are only qualitatively correct at intermediate temperatures, as can be seen by comparing Fig. 25(a) with the QMC and TMRG data and fit in Fig. 20. Troyer, Tsunetsugu and Würtz obtained a very good fit of $`\chi ^{(1)}(\mathrm{\Delta }^{},t)`$ to QMC $`\chi ^{}(t)`$ simulation data over a large temperature range for the $`S=1/2`$ two-leg Heisenberg ladder with spatially isotropic exchange; however, they assumed a $`\epsilon (\mathrm{\Delta }^{},k)`$ in the fit function which was later found to be inaccurate over much of the Brillouin zone.
We formulated an approximation \[designated as $`\chi ^{(2)}`$\] which is more accurate in the low-temperature range, and which we will use in the next section as a fit function to fit our TMRG $`\chi ^{}(t)`$ data at low $`t`$ to extract $`\mathrm{\Delta }^{}(\alpha )`$. The function $`\chi ^{(2)}`$ was obtained by summing the susceptibilities of isolated dimers with a distribution of singlet-trip-
let energy gaps given by our one-parameter dispersion relation $`\epsilon (\mathrm{\Delta }^{},k)`$ for $`0\mathrm{\Delta }^{}1`$ in Eqs. (3), which takes into account the interdimer interactions. Thus from Eq. (14) we simply obtain
$$\chi ^{(2)}(\mathrm{\Delta }^{},t)=\frac{1}{\pi t}_0^\pi \frac{dk}{3+\mathrm{e}^{\epsilon (\mathrm{\Delta }^{},k)/t}}.$$
(157)
Note that we make no assumptions here about the form of the function $`\mathrm{\Delta }^{}(\alpha )`$, since only $`\mathrm{\Delta }^{}`$ appears in the expression. This $`\chi ^{(2)}(\mathrm{\Delta }^{},t)`$ is exact in both the low- and high-$`t`$ limits, as is $`\chi ^{(1)}(\mathrm{\Delta }^{},t)`$, and both reproduce $`\chi ^{}(t)`$ for the isolated dimer ($`\mathrm{\Delta }^{}=1`$) exactly, but $`\chi ^{(2)}(\mathrm{\Delta }^{},t)`$ is more accurate at intermediate temperatures for $`\alpha 1`$ as shown in Fig. 25(b). In addition, by comparing $`\chi ^{(1)}(\mathrm{\Delta }^{},t)`$ and $`\chi ^{(2)}(\mathrm{\Delta }^{},t)`$ with the TMRG $`\chi ^{}(\alpha ,t)`$ calculations at low $`t`$, we found that the low-$`t`$ corrections to the low-$`t`$ limit in Eqs. (16) are much more accurately given by $`\chi ^{(2)}(t)`$ than by $`\chi ^{(1)}(t)`$. We will therefore not discuss $`\chi ^{(1)}(t)`$ further here.
At low temperatures, the approximation $`\chi ^{(2)}(\mathrm{\Delta }^{},t)`$ is expected to accurately describe the leading-order $`t`$ corrections to the low-$`t`$ limit only as long as the average number of magnons $`n_\mathrm{m}`$ occupying a state near the minimum in the one-magnon band is much less than unity. Using the expression for the boson occupation number
for this case,
$$n_\mathrm{m}=\frac{1}{\mathrm{e}^{\mathrm{\Delta }^{}/t}1},$$
(158)
yields $`t/\mathrm{\Delta }^{}=0.22`$ and 0.42 for $`n_\mathrm{m}=0.01`$ and 0.1, respectively. Thus, when fitting our low-$`t`$ TMRG $`\chi ^{}(\alpha ,t)`$ data by the fit function $`\chi ^{(2)}(\mathrm{\Delta }^{},t)`$ in the following section, we expect $`\chi ^{(2)}(\mathrm{\Delta }^{},t)`$ to be sufficiently accurate only for $`t/\mathrm{\Delta }^{}0.4`$. For this reason, our fits will be limited to this maximum scaled temperature.
We have computed $`\mathrm{\Delta }_{\mathrm{eff}}(\mathrm{\Delta }^{},T)/\mathrm{\Delta }`$ from $`\chi ^{(2)}(\mathrm{\Delta }^{},t)`$ in Eq. (157), using the definition in Eq. (156), and plot the results vs $`k_\mathrm{B}T/\mathrm{\Delta }`$ in Fig. 26. For the dimer ($`\mathrm{\Delta }^{}=1`$), one finds analytically that $`\mathrm{\Delta }_{\mathrm{eff}}(T)/\mathrm{\Delta }=12k_\mathrm{B}T/\mathrm{\Delta }`$ to lowest order in $`T`$. On the other hand, for $`0<\mathrm{\Delta }^{}<1`$, the initial dependence is positive and quadratic in $`T`$, and a maximum is seen in $`\mathrm{\Delta }_{\mathrm{eff}}(T)/\mathrm{\Delta }`$, which for $`\mathrm{\Delta }^{}0.4`$ occurs at $`t/\mathrm{\Delta }^{}k_\mathrm{B}T/\mathrm{\Delta }0.14`$ with a height of $`0.5`$%. This height is quantitatively consistent with the data in Fig. 24(b) derived from the TMRG $`\chi ^{}(t)`$ for $`\alpha =0.8`$. Thus the weak maximum seen in that figure is not a spurious effect.
### B Fits to the low-$`𝒕`$ TMRG $`𝝌^{\mathbf{}}\mathbf{(}𝜶\mathbf{,}𝒕\mathbf{)}`$ data
We were tempted to fit $`\mathrm{\Delta }_{\mathrm{eff}}^{}(\alpha ,t)`$ derived from the low-$`t`$ TMRG $`\chi ^{}(\alpha ,t)`$ data, as discussed above, to obtain the spin gaps $`\mathrm{\Delta }^{}(\alpha ,t)`$. However, this procedure would have weighted the $`\chi ^{}(\alpha ,t)`$ data in an ill advised way. We therefore decided to do conventional fits of the low-$`t`$ $`\chi ^{}(\alpha ,t)`$ data by the fit function $`\chi ^{(2)}(\mathrm{\Delta }^{},t)`$ in Eq. (157). For a given $`\alpha `$, this is a one-parameter ($`\mathrm{\Delta }^{}`$) fit function and the fits are therefore stringent tests of both the appropriateness of the fit function and the precision and accuracy of the data. Because the temperature dependence of the accuracy of the calculations is unknown except for the uniform chain data (see Fig. 19), we assumed that all data for a given $`\alpha `$ in a given fitted temperature range have the same accuracy. Thus in the nonlinear least-squares fits for each $`\alpha `$ we minimized the square of the relative rms deviation of the fit from the data
$$\sigma _{\mathrm{rms}}^2=\frac{1}{N_\mathrm{p}}\underset{i=1}{\overset{N_\mathrm{p}}{}}\frac{[\chi ^{(2)}(t_i)\chi ^{}(t_i)]^2}{[\chi ^{}(t_i)]^2},$$
(159)
where $`N_\mathrm{p}`$ is the number of data points fitted, which was usually between 250 and 1500.
Due to the presence of the spin gap $`\mathrm{\Delta }^{}`$ in the exponential of the fit function, $`\sigma _{\mathrm{rms}}`$ is extremely sensitive to the precise value of $`\mathrm{\Delta }^{}`$ when low-$`t`$ fits are carried out. For example, close to the optimum $`\mathrm{\Delta }^{}`$ fit value, a change in $`\mathrm{\Delta }^{}`$ by only 0.0001 ($`0.1\%`$) can change $`\sigma _{\mathrm{rms}}`$ by up to $`300`$ %. Thus a few percent accuracy in the $`\chi ^{}(t)`$ data at low $`t`$ is sufficient to allow $`\mathrm{\Delta }^{}`$ to be determined for a given fit to a precision better than 0.0001. For a given $`\alpha `$, the obtained $`\mathrm{\Delta }^{}`$ was found to be insensitive, typically to within $`0.000`$2, to the $`t`$ range of the fit, as long as the maximum fitted temperature satisfied $`t/\mathrm{\Delta }^{}0.4`$, consistent with the above discussion of the boson occupation number. This lack of sensitivity of the value of the fitted $`\mathrm{\Delta }^{}`$ to the precise fit range demonstrated that the fit function $`\chi ^{(2)}(t)`$ is an appropriate one. Depending on the $`\alpha `$ value and the $`t`$ range of the fit, $`\sigma _{\mathrm{rms}}`$ was typically between 0.1 % and several percent.
The $`\mathrm{\Delta }^{}(\alpha )`$ values obtained from the fits are listed in Table III, together with the estimated accuracies in parentheses. Note that a quoted accuracy is associated with variations in $`\mathrm{\Delta }^{}`$ in fits to a specific set of data for a given $`\alpha `$ over various temperature ranges as discussed above, and does not include possible systematic errors due to, e.g., the finite fixed number of states kept in the TMRG calculations. Also included in Table III are literature data which will be compared with the present results in the next section. Log-log plots of the low-$`t`$ data and fits are shown in Fig. 27, where on the scale of this figure, for most $`\alpha `$ values the data and fit are identical (cannot be distinguished) within the fitted temperature range. Extrapolations of the fits to higher and lower temperatures are also shown for comparison with the data.
## VII Comparisons of the Calculations with Previous Work
### A Spin gap
Our $`\overline{\mathrm{\Delta }^{}}(\delta )`$ spin gap data determined by fitting our TMRG $`\chi ^{}(t)`$ data by Eq. (157) are plotted in Fig. 28(a),
along with the results of previous workers listed in Table III. The solid curve is the function $`\overline{\mathrm{\Delta }^{}}=2\delta ^{3/4}`$ in Eq. (48) proposed by Barnes, Riera, and Tennant (BRT). The overall behavior of the data in Fig. 28(a) is well described by this function, but significant deviations of the data from the curve occur as illustrated in the expanded plot for $`\delta 0.1`$ in Fig. 28(b). The error bars are included with each plotted data point in Fig. 28(b), except for the data of Ref. which were not available, but they are all small and not clearly seen. Our values for $`\delta 0.1`$ are significantly smaller than those of Uhrig et al., where the differences are far outside the combined limits of error, and are larger than those of Augier et al.
As will be seen explicitly in Sec. VIII C below, our $`\overline{\chi ^{}}(\delta ,\overline{t})`$ fit function allows $`\delta (T)`$ to be determined for real materials by using the fit function to model experimental $`\chi (T)`$ data. However, if one would like to determine the spin gap $`\mathrm{\Delta }(T)`$ from the derived $`\delta (T)`$, an expression is needed for $`\mathrm{\Delta }(\delta )`$ over the entire range $`0\delta 1`$ in order to be generally useful and applicable. At present, the only extant expression is that of BRT in Eqs. (II C 2). As seen in Fig. 28 and in Table IV below, this expression is only an approximation that fits neither BRTs’ $`\mathrm{\Delta }^{}(\alpha )`$ data for $`0.1\alpha 0.9`$ nor our TMRG spin gap data for $`0.8\alpha 0.995`$ to within the respective error bars. To formulate a more flexible expression, we modify BRTs’ formula to read
$$\overline{\mathrm{\Delta }^{}}(\delta )\frac{\mathrm{\Delta }(\delta )}{J}=2\delta ^{y(\delta )},$$
(161)
so the $`\delta `$-dependent power $`y`$ is
$$y(\delta )=\frac{\mathrm{ln}[\mathrm{\Delta }(\delta )/2J]}{\mathrm{ln}\delta }.$$
(162)
The numerical prefactor “2” in Eq. (161) must be retained in order to reproduce the exact $`\overline{\mathrm{\Delta }^{}}(\delta =1)=2`$. Shown in Fig. 29(a) is a semilog plot of $`y`$ versus $`\delta `$ for the same numerical data as in Fig. 28. This plot \[and Fig. 29(b) below\] explicitly shows, from BRTs’ data, that the exponent deviates significantly from the value 3/4 even for $`\delta 1`$. The plot also clearly differentiates the various numerical data for small $`\delta `$ by the different groups, and shows that one of our two data points from the TMRG for $`\delta =0.0025`$ (the one derived from $`m=256`$ data at high $`t`$) is not in agreement with the trend of the remainder of our data. This data point will not be included in the plot and fit to be discussed in the next paragraph.
Our $`y(\delta )`$ data at small $`\delta `$ are in agreement with both the magnitude and trend of BRTs’ data at larger $`\delta `$. The
$`y(\delta )`$ for these two data sets from Fig. 29(a), with the exception of one of our two data points for $`\delta =0.0025`$ just noted above, are plotted together on an expanded vertical scale in Fig. 29(b) where a rather smooth behavior of $`y(\delta )`$ is seen over the combined range of the two calculations $`0.0025\delta 1`$. With the behavior in Fig. 29(b) in mind, we formulated a five-parameter fit function for these two combined $`y(\delta )`$ data sets that yields the correct limits $`\overline{\mathrm{\Delta }^{}}(\delta 0)=0`$ and $`\overline{\mathrm{\Delta }^{}}(\delta 1)=2`$, with the property $`lim_{\delta 0}y(\delta )=`$ const, given by
$`y(\delta )=y(1)`$ $`+`$ $`n_1\mathrm{tanh}\left[{\displaystyle \frac{\mathrm{ln}\delta }{m_1}}\mathrm{ln}\left({\displaystyle \frac{\mathrm{ln}\delta }{m_2}}\right)\right]`$ (163)
$`+`$ $`n_2\mathrm{tanh}^2\left[{\displaystyle \frac{\mathrm{ln}\delta }{m_1}}\mathrm{ln}\left({\displaystyle \frac{\mathrm{ln}\delta }{m_2}}\right)\right].`$ (164)
An unweighted fit of this expression to all the data in Fig. 29(b) yielded the parameters
$`y(1)=0.74922,n_1=0.00776,n_2=0.00685,`$
$$m_1=3.3297,m_2=2.2114,$$
(165)
so that
$$\underset{\delta 0}{lim}y(\delta )=y(1)n_1+n_2=0.7346.$$
(166)
The fit is plotted as the solid curve in Fig. 29(b). As can be seen from the figure, our data are fitted to within our error bars. In addition, when the $`y(\delta )`$ fit function in Eqs. (164) and (165) is inserted into Eq. (161), the predicted values of $`\overline{\mathrm{\Delta }^{}}(\delta )`$ are in agreement with each of the values of BRT at larger $`\delta `$ to within 0.0001, which is sufficient for modeling experimental data. The $`\delta =0`$ limit of $`y(\delta )`$ in Eq. (166) is in agreement with the theoretical effective value $`y(0)=0.72(3)`$, which was obtained without the log correction term by Singh and Weihong from an eleventh-order dimer series expansion of the triplet dispersion relation. We will use Eqs. (27) to compute $`\mathrm{\Delta }(T)`$ from the experimentally derived $`\delta (T)`$ for $`\mathrm{NaV}_2\mathrm{O}_5`$ in Sec. VIII C below.
In order to test the critical behavior prediction $`\overline{\mathrm{\Delta }^{}}=A\delta ^{2/3}/|\mathrm{ln}\delta |^{1/2}`$ in Eq. (49), which need only hold in the asymptotic critical regime $`\delta 0`$ in contrast to the fit function for $`0\delta 1`$ in Eqs. (27), in Fig. 30 is plotted $`\mathrm{\Delta }/J`$ vs $`\delta ^{2/3}/|\mathrm{ln}\delta |^{1/2}`$ in the region $`\delta 0.06`$ for the same data and symbols as in Fig. 28. A proportionality appears to be developing in our data for $`\delta 0.005`$, as shown by the straight line with slope $`A=3.3`$ passing through the origin of the figure, suggesting that the asymptotic critical regime begins with decreasing $`\delta `$ below $`\delta 0.005(\alpha 0.99)`$. High-accuracy $`\overline{\mathrm{\Delta }^{}}(\delta )`$ data for $`\delta 0.001`$ are needed to test this conjecture. From Fig. 30, the slope 3.3 of the line drawn is evidently a lower limit of the prefactor $`A`$ within the actual asymptotic critical regime.
### B Numerical $`𝝌^{\mathbf{}}\mathbf{(}𝜶\mathbf{,}𝒕\mathbf{)}`$ results
Barnes and Riera previously carried out exact diagonalizations of Hamiltonian (II C) for $`S=1/2`$ alternating chains of length up to 16 spins using the Lanczos technique. Their computed $`\chi ^{}(t)`$ values for $`\alpha =0.2,`$ 0.4, 0.6, 0.7, and 0.8 were extrapolated to the bulk limit and the results are shown as the symbols in Fig. 31(a). Our fit function as in Fig. 20 for the same $`\alpha `$ values is plotted as the solid curves in Fig. 31(a), which are seen to be in good overall agreement with the calculations of Barnes and Riera. The deviations of the data of Barnes and Riera from our fit function are plotted vs temperature in Fig. 31(b). The average deviation of their data from our fit function is very small for each data set: $`0.41,+0.33,0.40,0.26`$, and $`+0.79\times 10^4`$ for $`\alpha =0.2`$, 0.4, 0.6, 0.7, and 0.8, respectively. The absolute rms deviations $`\sigma _{\mathrm{rms}}`$ of their data from our fit function for $`\alpha =0.2`$, 0.4, 0.6, 0.7, and 0.8 are (in units of $`10^4`$) 1.73, 1.43, 0.73, 0.78, and 3.76, respectively. We conclude that their data are in good quantitative agreement with our data and fit, with the exception of their data point for $`\alpha =0.8`$ at their lowest temperature $`t=0.05`$.
### C Bulaevskii Theory
Bulaevskii calculated $`\chi ^{}(t)`$ analytically in the Hartree-Fock approximation. He first obtained an integral equation for the magnon spectrum $`E(k)`$:
$`\epsilon (k){\displaystyle \frac{E(k)}{J_1}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}[\sqrt{1+\alpha ^22\alpha \mathrm{cos}k}`$ (167)
$`+`$ $`{\displaystyle \frac{C_1+\alpha C_2(\alpha C_1+C_2)\mathrm{cos}k}{\sqrt{1+\alpha ^22\alpha \mathrm{cos}k}}}],`$ (168)
$`C_1(t)`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}{\displaystyle _0^\pi }dk{\displaystyle \frac{1\alpha \mathrm{cos}k}{\sqrt{1+\alpha ^22\alpha \mathrm{cos}k}}}\mathrm{tanh}{\displaystyle \frac{\epsilon (k)}{2t}},`$ (169)
$`C_2(t)`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}{\displaystyle _0^\pi }dk{\displaystyle \frac{\alpha ^2\alpha \mathrm{cos}k}{\sqrt{1+\alpha ^22\alpha \mathrm{cos}k}}}\mathrm{tanh}{\displaystyle \frac{\epsilon (k)}{2t}},`$ (171)
where $`k`$ is measured in units of $`2\pi /d`$. $`d1`$ is the lattice repeat distance along the chain, which is twice the average distance between spins. He then expressed $`\chi ^{}(t)`$ in terms of $`\epsilon (k)`$:
$`\chi ^{}(t)`$ $``$ $`{\displaystyle \frac{F(t)}{2+(1+\alpha )F(t)}},`$ (172)
$`F(t)`$ $`=`$ $`{\displaystyle \frac{1}{2\pi t}}{\displaystyle _0^\pi }{\displaystyle \frac{\mathrm{d}k}{\mathrm{cosh}^2[\epsilon (k)/(2t)]}}.`$ (174)
At $`t=0`$ and $`\alpha 0`$, from Eqs. (LABEL:EqEpsK2) we obtain
$`C_1(\alpha )`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}\{(1+\alpha )\mathrm{E}\left[{\displaystyle \frac{4\alpha }{(1+\alpha )^2}}\right]`$ (176)
$`+(1\alpha )\mathrm{K}\left[{\displaystyle \frac{4\alpha }{(1+\alpha )^2}}\right]\},`$
$`C_2(\alpha )`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}\{(1+\alpha )\mathrm{E}\left[{\displaystyle \frac{4\alpha }{(1+\alpha )^2}}\right]`$ (179)
$`(1\alpha )\mathrm{K}\left[{\displaystyle \frac{4\alpha }{(1+\alpha )^2}}\right]\},`$
where K$`(y)`$ and E$`(y)`$ are, respectively, the complete elliptic integrals of the first and second kinds. The dispersion relations versus $`\alpha `$ at $`t=0`$ are obtained by inserting Eqs. (LABEL:EqEpsK3) into (168) and a selection of results is shown in Fig. 32. From Eqs. (168) and (LABEL:EqEpsK3), at $`t=0`$ the spin-gap $`\mathrm{\Delta }_{k=0}(\alpha )`$ at $`k=0`$ is given by
$$\mathrm{\Delta }_{k=0}(\alpha )=\frac{1\alpha }{2}\left\{1+\frac{2(1\alpha )}{\pi }\mathrm{K}\left[\frac{4\alpha }{(1+\alpha )^2}\right]\right\}.$$
(180)
This expression gives the actual spin-gap for $`0<\alpha 0.79`$. However, for $`0.79\alpha <1`$, the minimum in the dispersion relation does not occur at $`k=0`$, as illustrated in an expanded plot of $`E(k)`$ for $`\alpha =0.9`$ in the inset to Fig. 32. The wave vector $`k_\mathrm{G}`$ at which the minimum spin gap $`\mathrm{\Delta }_\mathrm{B}`$ occurs is plotted versus $`\alpha `$ in Fig. 33(b). The $`\mathrm{\Delta }_\mathrm{B}`$ from Bulaevskii’s theory at $`t=0`$ is plotted versus $`\alpha `$ as the solid curve in Fig. 33(a), and a few representative values are given in Table IV. The predictions of Bulaevskii’s theory are in very good agreement with those of Barnes, Riera, and Tennant for $`\alpha 0.3`$, but the agreement becomes progressively worse as $`\alpha `$ increases further. From Eqs. (168) and (LABEL:EqEpsK2), $`E(k)`$ is temperature dependent. In addition, in the range $`0.79\alpha <1`$ for which $`k_\mathrm{G}0`$ at $`t=0`$, we find that $`k_\mathrm{G}`$ depends on $`t`$, as shown in Fig. 34. From Fig. 34, $`k_\mathrm{G}0`$ at $`t0.083`$, 0.122, 0.131, 0.125, and 0.111 for $`\alpha =0.8`$, 0.85, 0.9, 0.95, and 0.99, respectively.
We computed $`\chi ^{}(t)`$ by inserting $`\epsilon (k)`$ in Eq. (168) into Eqs. (LABEL:EqEpsK2), numerically solving the latter two simultaneous equations for $`C_1`$ and $`C_2`$ at each $`t`$, and then inserting the resulting $`\epsilon (k)`$ into Eqs. (LABEL:EqChiBul). The progression of $`\chi ^{}(t)`$ with increasing $`\alpha `$ from 0.001 to 0.99 is shown in Fig. 35. As noted by Bulaevskii, the values of $`\chi ^{}`$ at the maxima are too large and the temperatures at which these occur are too small by $`5`$–10% (compare Fig. 35 with Fig. 20).
At low temperatures $`0.033t1/4`$, Bulaevskii fitted $`\chi ^{}(\alpha ,t)`$, calculated from Eqs. (LABEL:EqChiBul), by the two-parameter form
$$\chi ^{}(\alpha ,t)=\frac{a(\alpha )}{t}\mathrm{e}^{\mathrm{\Delta }(\alpha )/(J_1t)},$$
(181)
and obtained values of $`a_\mathrm{B}`$ and $`\mathrm{\Delta }_{\mathrm{B},\mathrm{Fit}}/J_1`$ for $`0\alpha 0.9`$ which are reproduced in Table IV; $`\mathrm{\Delta }_{\mathrm{B},\mathrm{Fit}}(\alpha )/J_1`$ is plotted as the open squares in Fig. 33(a). Note that the temperature exponent in the prefactor to the exponential is $`\gamma =1`$, contrary to the $`\gamma =1/2`$ in Eq. (19) which is expected in the low-$`t`$ limit for any 1D $`S=1/2`$ Heisenberg spin system with a spin gap (and with a nondegenerate one-magnon band with a parabolic minimum). We have confirmed that over the temperature range fitted by Bulaevskii, one indeed obtains $`\gamma 1`$ for the best fit of Eq. (181) to numerical calculations of $`\chi ^{}(\alpha ,t)`$. We infer that the discrepancy between Bulaevskii’s $`\gamma =1`$ and the expected $`\gamma =1/2`$ arises because the fits were not carried out completely within the low-$`t`$ limit. This issue is discussed in more detail below.
Equation (181) together with Bulaevskii’s table of {$`a_\mathrm{B}(\alpha ),\mathrm{\Delta }_{\mathrm{B},\mathrm{Fit}}(\alpha )/J_1`$} values were subsequently used extensively in the analysis of experimental $`\chi (T)`$ data for compounds exhibiting spin-Peierls transitions to determine the alternation parameter $`\alpha `$ at low temperatures $`TT_\mathrm{c}`$ where the experimental spin-gap is nearly independent of $`T`$. However, from Table IV and Fig. 33(a), the $`\mathrm{\Delta }_{\mathrm{B},\mathrm{Fit}}(\alpha )`$ values of Bulaevskii are in generally poor agreement with the actual spin gaps $`\mathrm{\Delta }_\mathrm{B}(\alpha )`$ of his theory and with those \[$`\mathrm{\Delta }_{\mathrm{BRT}}(\alpha )`$\] calculated for the same $`\alpha `$ values by Barnes, Riera, and Tennant. Therefore, one should consider the $`\mathrm{\Delta }_{\mathrm{B},\mathrm{Fit}}`$ parameters as fitting parameters only, with no direct relation to the actual spin gap.
According to Eq. (14) for $`\chi ^{}(t)`$ of the isolated dimer which is a zero-dimensional spin system, the form of $`\chi ^{}(t)`$ in Eq. (181) with $`\gamma =1`$ is correct for $`\alpha =0`$ and $`t0`$. On the other hand, for one-dimensional spin systems such as the two-leg spin ladder (and the alternating-exchange chain) at temperatures $`k_\mathrm{B}T\mathrm{\Delta }`$ and $`k_\mathrm{B}T`$ one-magnon bandwidth, Eqs. (16) apply, with $`\gamma =1/2`$, assuming that the triplet one-magnon dispersion relation $`E(k)`$ is parabolic at the minimum. In this case one expects $`\gamma =1/2`$ at sufficiently low $`t`$ for any finite $`\alpha `$. Thus, in the temperature region of validity of Eq. (18), a plot of the left-hand-side of Eq. (21) vs ln$`t`$ should give a straight line with slope $`\gamma `$. Shown in Fig. 36 are such plots, obtained using our $`\chi ^{}(\alpha ,t)`$ calculated from Bulaevskii’s theory as described above, for $`\alpha =0.001`$ to 0.99. For $`\alpha =0.001`$, a crossover is clearly evident from $`\gamma =1`$ to $`\gamma =1/2`$ with decreasing $`t`$. The other curves also exhibit signs of a crossover, with $`\gamma 1/2`$ at the lowest temperatures, with the exception of the curve for $`\alpha =0.8`$. For this $`\alpha `$ value, which is just above the value $`\alpha 0.79`$ at which $`k_\mathrm{G}`$ becomes nonzero at $`t=0`$ \[see Fig. 33(b)\], the $`\gamma `$ at the lowest $`t`$ is intermediate between the values of 1/2 and 1, and the assumption of a parabolic form for $`E(k)`$ at the band minimum is evidently not satisfied (see Fig. 32). In fact, Troyer, Tsunetsugu and Würtz calculated the low-$`t`$ limit of $`\chi ^{}(t)`$ for 1D systems with general dispersion relation $`\epsilon (k)=\mathrm{\Delta }^{}+c^{}|ka|^n`$, where $`k`$ is the deviation of the wave vector from that at the band minimum. They found the same form $`\chi ^{}(t)=(A_n/t^\gamma )\mathrm{exp}(\mathrm{\Delta }^{}/t)`$ as for the parabolic case $`n=2`$, but where $`\gamma =1(1/n)`$. Thus, e.g., $`\gamma =2/3`$ and 3/4 for $`n=3`$ and 4, respectively. This range of $`\gamma `$ is consistent with the slope of the data at the lowest temperatures for $`\alpha =0.8`$ in Fig. 36.
The predictions of Bulaevskii’s theory for $`\chi ^{}(t)`$ from Fig. 35 are compared with our $`\chi ^{}(\alpha ,t)`$ fit function (solid
curves) for $`\alpha =0.2`$, 0.4, 0.6, 0.8 and 0.99 (as in Fig. 20) in Fig. 37, where the Bulaevskii prediction for each of these $`\alpha `$ values is shown as the corresponding dashed curve. The disagreement between the two calculations becomes progressively more severe as temperature decreases and as the uniform chain limit is approached with increasing $`\alpha `$. Therefore, the accuracies of the $`\alpha `$ and $`J_1`$ values previously extracted from experimental data at low $`T`$ for compounds with $`\alpha 1`$ using Bulaevskii’s theory are unclear. Our $`\chi ^{}(\alpha ,t)`$ fit function now provides a much more accurate and reliable means of extracting exchange constants and spin gaps from experimental $`\chi (T)`$ data.
## VIII Magnetic Susceptibility of NaV<sub>2</sub>O<sub>5</sub>
Crystals of $`\mathrm{Na}_{0.996(3)}\mathrm{V}_2\mathrm{O}_{5.00(6)}`$ were grown at the Max-Planck-Institut für Festkörperforschung, Stuttgart, in a Pt crucible in flowing Ar atmosphere by a self-flux method from a 5:1:1 mixture of $`\mathrm{NaVO}_3`$, $`\mathrm{V}_2\mathrm{O}_3`$ and $`\mathrm{V}_2\mathrm{O}_5`$. The flux was dissolved by boiling the solidified melt in distilled water. X-ray powder diffraction patterns collected with a STOE diffractometer yielded the lattice parameters $`a=11.3187(8)`$ Å, $`b=3.6111(3)`$ Å, and $`c=4.8007(5)`$ Å. Chemical analyses on two independent representative samples of the batch were performed with a standard AAS analysis technique for V and Pt and ICP emission spectroscopy for the Na content. The oxygen content was determined by measuring with IR spectroscopy the amount of CO generated when the sample is fused in a graphite crucible at 2700 $`{}_{}{}^{\mathrm{o}}\mathrm{C}`$ in vacuo. Platinum impurities above the level of sensitivity of the analysis (500 ppm with respect to V) could not be detected.
At Ames Laboratory, single crystals of $`\mathrm{NaV}_2\mathrm{O}_5`$ were grown out of the ternary melt. Powders of $`\mathrm{V}_2\mathrm{O}_5`$ and V<sub>2</sub>O<sub>3</sub> were prepared by oxidizing and reducing NH<sub>4</sub>VO<sub>3</sub> at 600$`^{}`$C and 900 C, respectively. The resulting $`\mathrm{V}_2\mathrm{O}_5`$ is reacted with Na<sub>2</sub>CO<sub>3</sub> at 550 C yeilding NaVO<sub>3</sub>. About 10 grams of NaVO<sub>3</sub>, $`\mathrm{V}_2\mathrm{O}_5`$, and V<sub>2</sub>O<sub>3</sub> in the molar ratio 32:1:1 were placed in a Pt crucible and sealed in an evacuated quartz tube. The melt was then slowly cooled from 800 to 660 C over 50 h and the remaining liquid was decanted. Small amounts of solidified melt remaining on the crystals were dissolved with hot water. Typical dimensions of the ribbon-shaped crystals grown in this manner are $`0.5\times 1.5\times 11`$ mm<sup>3</sup> with the $`c`$ axis perpendicular to the plane of the ribbon, the $`b`$ axis along the length of the ribbon and the $`a`$ axis along the width of the ribbon, with lattice parameters $`a11.303`$ Å, $`b3.611`$ Å, and $`c4.752`$ Å. The crystal denoted as AL1 has a mass of 8.2 mg and approximate planar dimensions $`1.5\times 2.5`$ mm<sup>2</sup>.
The magnetic susceptibility $`\chi (T)M(T)/H`$ of the crystals was measured using Quantum Design SQUID magnetometers at Stuttgart and Ames. The measurements on eight crystals of NaV<sub>2</sub>O<sub>5</sub> in Stuttgart were carried out in a field $`H=1`$ T along the V ladder ($`b`$) axis direction in various temperature ranges between 2 and 750 K. Measurements of the anisotropy of $`\chi (T)`$ along the $`a`$, $`b`$, and $`c`$ axis directions were carried out from 2 to 300 K in Ames on crystal AL1 in $`H=2`$ T.
The results for two of the crystals up to 750 K are shown in Fig. 38. The data illustrate the variabilities we have observed between measurements along the same axis on different crystals. Above $`50`$ K, the two data sets are nearly parallel, with the difference between them being $`34\times 10^5`$ cm<sup>3</sup>/mol; we have no explanation for this difference, and no comments have been made in the literature about such variabilities and/or their origins in $`\chi (T)`$ along the same axis in different crystals that we are aware of. The data from $`T_\mathrm{c}33`$–34 K up to 300 K are in approximate agreement with the single crystal data of Isobe, Kagumi and Ueda taken in this $`T`$ range along the same axis in $`H=5`$ T. A variable Curie-Weiss-like contribution $`\chi ^{\mathrm{CW}}(T)`$ to $`\chi (T)`$ occurs below $`20`$ K which is attributed to paramagnetic defects, impurities, inclusions and/or intergrowths in the crystals. The “Fit” shown in the figure will be discussed later in Sec. VIII B.
The experimental data are analyzed with the general expression
$$\chi (T)=\chi _0+\chi ^{\mathrm{CW}}(T)+\chi ^{\mathrm{spin}}(T),$$
(183)
$$\chi _0=\chi ^{\mathrm{core}}+\chi ^{\mathrm{VV}},$$
(184)
$$\chi ^{\mathrm{CW}}(T)=\frac{C_{\mathrm{imp}}}{T\theta },$$
(185)
$$\chi ^{\mathrm{spin}}(T)=\frac{Ng^2\mu _\mathrm{B}^2}{J}\overline{\chi ^{}}\left(\frac{k_\mathrm{B}T}{J}\right),$$
(186)
where $`\chi _0`$ is the sum of a temperature independent and (nearly) isotropic orbital diamagnetic core contribution and a usually anisotropic and temperature independent orbital paramagnetic Van Vleck contribution. We estimate $`\chi ^{\mathrm{core}}`$ using the values $``$5, $``$7, $``$4, and $`12\times 10^6`$ cm<sup>3</sup>/mol for Na<sup>+1</sup>, V<sup>+4</sup>, V<sup>+5</sup>, and O<sup>-2</sup>, respectively, yielding the isotropic value
$$\chi ^{\mathrm{core}}=7.8\times 10^5\frac{\mathrm{cm}^3}{\mathrm{mol}\mathrm{NaV}_2\mathrm{O}_5}.$$
(187)
The second term in Eq. (183) is the above-noted Curie-Weiss impurity and/or defect contribution and the last term is the intrinsic spin susceptibility, each of which may or may not be anisotropic. For a Heisenberg spin system, $`\overline{\chi ^{}}`$ is isotropic, and therefore so is $`\chi ^{\mathrm{spin}}`$ apart from anisotropy in the $`g`$ factor. The impurity Curie-Weiss term $`\chi ^{\mathrm{CW}}(T)`$ can be anisotropic if the impurities are defects or intergrowths in the crystals which have atomic coordination principal axes which are fixed with respect to the crystal axes rather than being randomly oriented. We model our $`\chi (T)`$ data according to Eq. (183) in terms of the $`\overline{\chi ^{}}(t)`$ in Eq. (186), which are (fit functions to) theoretical susceptibility calculations presented in previous sections. Before moving on to do that, we first experimentally examine the anisotropy in $`\chi (T)`$ of $`\mathrm{NaV}_2\mathrm{O}_5`$ and its implications in the next section.
### A Anisotropy of the magnetic susceptibility
The magnetic susceptibilities of $`\mathrm{NaV}_2\mathrm{O}_5`$ crystal AL1 along the $`a`$, $`b`$, and $`c`$ axes are plotted vs temperature in Fig. 39(a), where the $`a`$ and $`c`$ axes are perpendicular to the V chains which run along the $`b`$ axis, and the $`c`$ axis is perpendicular to the trellis layers that the V chain/ladders reside in. The data are similar to the aniso-
tropic $`\chi (T)`$ data reported by Isobe, Kagami and Ueda, although the anisotropies we measure at both room temperature and at low temperatures are somewhat larger than they reported. The anisotropies at low temperatures are seen more clearly if the respective impurity term $`\chi ^{\mathrm{CW}}(T)`$ in Eq. (185) is subtracted from each data set, as shown in Fig. 39(b). The impurity Curie constant $`C_{\mathrm{imp}}`$ and Weiss temperature $`\theta `$ for each direction of the applied field were determined by the requirement that $`\chi (T)`$ become independent of $`T`$ for $`T0`$. The fitted values of $`C_{\mathrm{imp}}`$ were found to be slightly anisotropic and are given in Table VI below. The values of $`C_{\mathrm{imp}}`$ are equivalent to the contribution of only 0.07 mol% of $`S=1/2`$ impurities with $`g=2`$; if the impurity spin is actually greater than 1/2, the concentration of paramagnetic impurities could be much less than this estimate. From a comparison of Figs. 39(a) and 39(b), $`\chi ^{\mathrm{CW}}(T)`$ is seen to make a negligible contribution to the measured $`\chi (T)`$ above $`100`$ K. Since in the presence of a spin gap $`\chi ^{\mathrm{spin}}=0`$ at the lowest temperatures, from Eqs. (38) and Fig. 39(b) we obtain
$`\chi _b^{\mathrm{VV}}`$ $`=`$ $`18.7\times 10^5{\displaystyle \frac{\mathrm{cm}^3}{\mathrm{mol}}},\chi _c^{\mathrm{VV}}=13.3\times 10^5{\displaystyle \frac{\mathrm{cm}^3}{\mathrm{mol}}},`$ (188)
$`\chi _a^{\mathrm{VV}}`$ $`=`$ $`20.0\times 10^5{\displaystyle \frac{\mathrm{cm}^3}{\mathrm{mol}}}(TT_\mathrm{c}).`$ (190)
From Fig. 39(b), the anisotropies of $`\chi (T)`$ are seen to be quite temperature-dependent upon heating through $`T_\mathrm{c}=33.4`$ K. These results are surprising, because $`\chi ^{\mathrm{spin}}`$ is expected to be isotropic (apart from the small anisotropy due to the anisotropic $`g`$ factor), with $`\chi ^{\mathrm{spin}}(T0)=0`$ because of the spin gap, and the anisotropic $`\chi ^{\mathrm{VV}}`$ values are expected to be temperature independent for our $`S=1/2`$ system over the temperature range of our measurements. Thus one expects the difference $`\chi _\alpha (T)\chi _\beta (T)(\alpha ,\beta =a,b,c)`$ to be nearly independent of temperature compared with the magnitude of either, where a subscript refers to the crystallographic axis along which the magnetic field is applied.
To be more quantitative, we define the anisotropy in the intrinsic susceptibility as
$$\mathrm{\Delta }\chi _{\alpha \beta }(T)[\chi _\alpha \chi _\alpha ^{\mathrm{CW}}](T)[\chi _\beta \chi _\beta ^{\mathrm{CW}}](T),$$
(191)
which eliminates extrinsic anisotropy in the Curie-Weiss impurity contribution from the values calculated from the experimental data. For example, according to this definition, $`\mathrm{\Delta }\chi _{ac}(T)`$ is the difference between the uppermost and lowermost data sets in Fig. 39(b). The three $`\mathrm{\Delta }\chi _{\alpha \beta }(T)`$ anisotropies are plotted in Fig. 40. It seems to us that the only reasonable explanation for the strong temperature-dependent anisotropies in Fig. 40 for two of the three data sets is that one or more of the $`\chi _\alpha ^{\mathrm{VV}}`$ susceptibilities is strongly temperature dependent near $`T_\mathrm{c}`$, contrary to our initial expectations. Such a temperature dependence may be associated with the crystallographic and charge-ordering transitions which occur at or near the same temperature as the spin dimerization transition, as discussed in the Introduction.
One can make rather strong general statements about the magnetic susceptibility anisotropies and their temperature dependences as follows. Defining the Van Vleck susceptibility anisotropy $`\mathrm{\Delta }\chi _{\alpha \beta }^{\mathrm{VV}}=\chi _\alpha ^{\mathrm{VV}}\chi _\beta ^{\mathrm{VV}}`$ and similarly the spin susceptibility anisotropy $`\mathrm{\Delta }\chi _{\alpha \beta }^{\mathrm{spin}}=\chi _\alpha ^{\mathrm{spin}}\chi _\beta ^{\mathrm{spin}}`$, from Eqs. (38) one obtains an expression for the anisotropy $`\mathrm{\Delta }\chi _{\alpha \beta }(T)`$ for a spin system in which the only anisotropy in $`\chi ^{\mathrm{spin}}`$ arises from anisotropy in the $`g`$ factor, given by
$$\mathrm{\Delta }\chi _{\alpha \beta }(T)=\mathrm{\Delta }\chi _{\alpha \beta }^{\mathrm{VV}}+\frac{N\mu _\mathrm{B}^2(g_\alpha ^2g_\beta ^2)}{J}\overline{\chi ^{}}\left(\frac{k_\mathrm{B}T}{J}\right).$$
(192)
The reduced spin susceptibility $`\overline{\chi ^{}}(\overline{t})`$ is necessarily positive, and it is isotropic for a Heisenberg spin system as noted above. Thus, if $`\chi _\alpha ^{\mathrm{VV}}`$ and $`\chi _\beta ^{\mathrm{VV}}`$ and therefore $`\mathrm{\Delta }\chi _{\alpha \beta }^{\mathrm{VV}}`$ are independent of temperature, the slope $`\mathrm{\Delta }\chi _{\alpha \beta }(T)/T`$ must have the same sign as the difference $`g_\alpha ^2g_\beta ^2`$. As discussed in the next subsection, for $`\mathrm{NaV}_2\mathrm{O}_5`$, this difference has been reported to be positive for
$`\alpha \beta =ac`$ and $`bc`$ and near zero for $`\alpha \beta =ab`$, consistent with the slopes in Fig. 40. However, in a simple ionic crystalline electric field model and with a positive spin-orbit coupling parameter for V one would predict that a $`\chi _\alpha ^{\mathrm{VV}}`$ should increase with the negative deviation of $`g_\alpha `$ from the free electron value $`g=2`$. Thus, a particularly visible and puzzling discrepancy is that since $`(2g_a)(2g_b)<(2g_c)`$ according to the reported $`g_\alpha `$ values below, on this basis one strongly expects $`\chi _a^{\mathrm{VV}}\chi _b^{\mathrm{VV}}<\chi _c^{\mathrm{VV}}`$; thus two of the three $`\chi _\alpha ^{\mathrm{VV}}`$ values should be about the same and smaller than the third one. Qualitatively contrary to this expectation, for $`TT_\mathrm{c}`$ we observe in Eq. (190) that $`\chi _a^{\mathrm{VV}}\chi _b^{\mathrm{VV}}>\chi _c^{\mathrm{VV}}`$.
We will not emphasize or further discuss these puzzling discrepancies with expectation with respect to their possible influence on our theoretical modeling of our $`\chi (T)`$ data in Secs. VIII B and VIII C, since at present there is no way to model, e.g., a temperature dependent Van Vleck susceptibility which changes rapidly near $`T_\mathrm{c}`$, but the anisotropic susceptibility results and the above discussion should be kept in mind. In the following two subsections the reported anisotropies in the $`g`$ factor as measured using electron spin resonance (ESR) and in the Van Vleck susceptibility as deduced from nuclear magnetic resonance (NMR) measurements will be discussed, respectively, in light of our anisotropic $`\chi (T)`$ data.
##### Anisotropy in the $`g`$ factor from ESR.
Many ESR measurements have recently been reported for NaV<sub>2</sub>O<sub>5</sub>. Each study found a signal with $`g2`$ which was attributed to bulk V species, and the $`g`$ values found in the various studies were the same within the errors, e.g.,
$$g_a=g_b=1.972(2),g_c=1.938(2).$$
(193)
The powder-average value is $`g=\sqrt{(g_a^2+g_b^2+g_c^2)/3}=1.961(2)`$. The $`g`$ values were found to be independent of $`T`$ down to 20 K, which is below $`T_\mathrm{c}`$. From all these measurements, there is no indication that the $`S=1/2`$ Heisenberg Hamiltonian is not appropriate to the spin system in NaV<sub>2</sub>O<sub>5</sub>. Unfortunately, given the sensitivity of the ESR technique, we cannot be certain that these ESR results are representative of the bulk spin species in NaV<sub>2</sub>O<sub>5</sub>, because no quantitative measurements of the concentration of spin species observed in these measurements were reported. Although the (uncalibrated) ESR intensity versus temperature measurements approximately mirror the bulk susceptibility behavior in most (but not all) of these studies, it is still possible that the signal arises from a minority spin species that is coupled to the bulk spin system. An interesting related issue which has not been discussed in the literature is why the presumed bulk $`S=1/2`$ species in NaV<sub>2</sub>O<sub>5</sub> are observable to low temperatures $`T0.03J/k_\mathrm{B}`$ by ESR, where the AF exchange constant is $`J/k_\mathrm{B}580`$ K (see below), whereas the bulk Cu<sup>+2</sup> spins 1/2 in the high-$`T_\mathrm{c}`$ cuprates are not observable by ESR up to 1100 K, which is $`0.7J/k_\mathrm{B}`$ where $`J/k_\mathrm{B}1600`$ K is only a factor of 2.8 larger.
In Ref. the authors estimated the $`\chi ^{\mathrm{VV}}`$ values using the reported anisotropic $`g`$ values obtained from ESR measurements, obtaining $`\chi _a^{\mathrm{VV}}=\chi _b^{\mathrm{VV}}=2.4\times 10^5`$ cm<sup>3</sup>/mol and $`\chi _c^{\mathrm{VV}}=6.6\times 10^5`$ cm<sup>3</sup>/mol, which were stated to be in agreement with the values from their $`K`$-$`\chi `$ analysis discussed in the following subsection. These values do not agree with our $`T=0`$ values in Eq. (190). In addition, from the $`\chi ^{\mathrm{VV}}`$ values of Ohama et al., one obtains $`\mathrm{\Delta }\chi _{ca}^{\mathrm{VV}}=\mathrm{\Delta }\chi _{cb}^{\mathrm{VV}}=4.2\times 10^5`$ cm<sup>3</sup>/mol, which are similar in magnitude but opposite in sign to our data in Eq. (190). If the strong change in each of $`\mathrm{\Delta }\chi _{ac}`$ and $`\mathrm{\Delta }\chi _{bc}`$ below $`T_\mathrm{c}`$ in Fig. 40 is due to a respective $`\mathrm{\Delta }\chi _{\alpha \beta }^{\mathrm{VV}}`$ which is strongly temperature dependent in this temperature range, an effect similar to that reported to occur from NMR measurements discussed in the next subsection, it is hard to understand why this change is not reflected in a distinct change in the reported temperature dependent anisotropy of the $`g`$ values at $`T_\mathrm{c}`$ if these $`g`$-value measurements are recording the characteristics of the bulk phase.
##### Anisotropy in the Van Vleck susceptibility from NMR.
From a so-called $`K`$-$`\chi `$ analysis using NMR paramagnetic nuclear resonance shift $`K(T)`$ data, combined with $`\chi (T)`$ measurements, under certain assumptions $`\chi ^{\mathrm{VV}}`$ can be obtained if $`K`$ is proportional to $`\chi `$, with $`T`$ as an implicit parameter. In this way, $`\chi ^{\mathrm{VV}}`$ values have been obtained by Ohama and coworkers for NaV<sub>2</sub>O<sub>5</sub> using <sup>23</sup>Na (Ref. ) and <sup>51</sup>V (Ref. ) NMR measurements on the same aligned powder sample. The former <sup>23</sup>Na study yielded $`\chi _b^{\mathrm{VV}}=23\times 10^5`$ cm<sup>3</sup>/mol below $`T_\mathrm{c}`$ and $`16\times 10^5`$ cm<sup>3</sup>/mol above $`T_\mathrm{c}`$, corresponding to a decrease of $`7\times 10^5`$ cm<sup>3</sup>/mol at $`T_\mathrm{c}`$. Their low temperature value is quite similar to our value in Eq. (190).
The <sup>51</sup>V NMR study, carried out above $`T_\mathrm{c}`$, yielded $`\chi _b^{\mathrm{VV}}=2(1)\times 10^5`$ cm<sup>3</sup>/mol, roughly an order of magnitude smaller than obtained in the authors’ first study (no comment was made about this discrepancy), and in addition gave $`\chi _a^{\mathrm{VV}}=1(1)\times 10^5`$ cm<sup>3</sup>/mol and $`\chi _c^{\mathrm{VV}}=4(1)\times 10^5`$ cm<sup>3</sup>/mol. These values are significantly smaller than our values. We note that a $`K`$-$`\chi `$ analysis on the $`d^1\mathrm{V}^{+4}`$ compound VO<sub>2</sub> yielded $`\chi ^{\mathrm{VV}}=6.5\times 10^5`$ cm<sup>3</sup>/mol.
### B Modeling the susceptibility of $`\mathrm{𝐍𝐚𝐕}_\mathrm{𝟐}𝐎_\mathrm{𝟓}`$ above $`𝑻_𝐜`$
Turning now to the experimental $`\chi (T)`$ data in Fig. 38, we have $`T^{\mathrm{max}}370`$ K. Assuming the validity of the Hamiltonian (1), Eq. (67) for the uniform Heisenberg chain yields the exchange constant $`J/k_\mathrm{B}580`$ K. Then the $`g_b`$ value in Eq. (193) and our $`\chi _0`$ values at $`T=0`$ in Table VI below, together with Eqs. (70) and (38), predict that the measured $`\chi ^{\mathrm{max}}40\times 10^5`$ cm<sup>3</sup>/mol, which is similar to the measured values of $`44`$ and $`48\times 10^5`$ cm<sup>3</sup>/mol for the two crystals in Fig. 38, respectively. We therefore proceeded to try to fit the data by the uniform chain model. The “Fit”, shown as the solid curve in Fig. 38, is a plot of Eqs. (38), with $`\chi ^{}(t)`$ being the susceptibility of the uniform chain (Fit 2 above) and with the parameters
$`\chi _0=8\times 10^5{\displaystyle \frac{\mathrm{cm}^3}{\mathrm{mol}}},C_{\mathrm{imp}}=0,`$ (194)
(195)
$`g=1.972,{\displaystyle \frac{J}{k_\mathrm{B}}}=580\mathrm{K}.(\mathrm{`}\mathrm{`}\mathrm{Fit}\mathrm{"})`$ (196)
This “Fit” is not really a fit, since we just set the $`g`$ and $`J`$ values to those estimated above and then set $`\chi _0`$ so that the calculated curve is in the vicinity of the data, because no small change in the parameters can bring the theory in agreement with the data. It is clear that adjusting $`\chi _0`$ further will not improve the agreement, nor will including a nonzero impurity Curie constant $`C_{\mathrm{imp}}`$. However, the shapes of the curve and the data are similar, so the agreement can be improved considerably (not shown) by simultaneously decreasing $`\chi _0`$ to $`10\times 10^5\mathrm{cm}^3/\mathrm{mol}`$, which is not possible physically according to Eqs. (38) because it would require the Van Vleck susceptibility to be negative, and increasing $`g`$ to the unphysically large value of $`2.4`$, while keeping $`J`$ constant. These results are in disagreement with the conclusion of Isobe and Ueda who found that the Bonner-Fisher prediction fitted their powder susceptibility data from 50 to 700 K very well assuming $`g=2`$. We can only note that their $`\chi (T)`$ data have not been quantitatively reproduced in either their or others’ subsequent measurements on NaV<sub>2</sub>O<sub>5</sub>, including ours, and that the Bonner-Fisher prediction is not accurate at temperatures below $`J/(4k_\mathrm{B})145`$ K as discussed in the Introduction.
Lohmann et al. and Hemberger et al. also previously concluded that the $`\chi (T)`$ of NaV<sub>2</sub>O<sub>5</sub> is not described (below 250 K) by the prediction for the $`S=1/2`$ Heisenberg chain, based on their fits by the Bonner-Fisher prediction to their $`\chi (T)`$ deduced from ESR measurements up to 650 K. They suggested that additional exchange couplings may be required to explain the observed $`\chi (T)`$. We consider this possibility here by modeling the influence of possible interchain spin coupling. Because there are no accurate and generally applicable numerical calculations for this case that we are aware of, we utilize the following simple molecular field theory (MFT) prediction for the spin susceptibility
$$\frac{1}{\chi ^{}(t)}=\frac{1}{\chi _{\mathrm{chain}}^{}(t)}+\frac{z^{}J^{}}{J},$$
(197)
where $`\chi _{\mathrm{chain}}^{}(t)`$ is the reduced spin susceptibility of the isolated quantum $`S=1/2`$ uniform Heisenberg chain (our Fit 2 above). The parameter $`z^{}`$ is the effective number of spins on other chains to which a spin in a given chain is coupled with effective (or average) exchange constant $`J^{}`$. To be consistent with our sign convention for the intrachain exchange constant $`J`$, $`J^{}`$ is positive for AF interactions and negative for ferromagnetic (FM) interactions. Equation (197) is very accurate when $`|z^{}J^{}/J|1`$.
We fitted the $`\chi (T)`$ data above 50 K for the two crystals in Fig. 38 by Eqs. (38) and (197), where we fixed $`g_b=1.972`$ and $`C_{\mathrm{imp}}=0`$ and allowed $`\chi _0`$, $`J`$ and the product $`z^{}J^{}`$ to vary. Very good fits were obtained, for which the fitting parameters are given in Table V. The fits are plotted as the solid curves in Fig. 41. For the parameters of the two crystals taken together, the fitted $`J/k_\mathrm{B}`$ = 584(9) K is the same as deduced above (580 K) from the temperature of the maximum in $`\chi (T)`$, and the fitted $`\chi _0=1.4(16)\times 10^5`$ cm<sup>3</sup>/mol is similar to our results at low temperatures in Fig. 39(b). A moderately large and negative (FM) interchain coupling $`z^{}J^{}/J=1.26(5)`$ was obtained. This coupling is sufficiently strong that long-range magnetic ordering might be expected, but which is not observed, possibly due to magnetic frustration effects. If the present mean-field interchain coupling analysis is correct, this interchain coupling should be evident in the magnon dispersion relations observable by inelastic magnetic neutron scattering measurements. Indeed, moderately strong dispersions of 1.4 meV in each of two bands perpendicular to the chains have in fact been observed by Yosihama et al. in such measurements on single crystals. It remains to be seen whether the magnitude and sign of the interchain exchange coupling that we infer in the mean-field analysis are consistent with the dispersion relations deduced from the neutron scattering data.
An alternative and/or additional mechanism which can produce a strong deviation of the measured $`\chi (T)`$ of a uniform chain compound from that predicted for Heisenberg uniform and alternating chains is the spin-phonon interaction. At low $`T`$ this interaction can lead to a spin-Peierls transition and can strongly modify $`\chi (T)`$ above $`T_\mathrm{c}`$ from that expected for the Heisenberg chain. Sandvik, Singh, and Campbell carried
out a detailed QMC investigation of a spin-Peierls model in which the spin 1/2 interactions were modified by the presence of dynamical (quantum mechanical) dispersionless Einstein phonons. For particular values of the spin-phonon coupling constant and phonon frequency, they found that the effective exchange constant $`J_{\mathrm{eff}}`$ decreases strongly with increasing $`T`$, and at $`T=0`$ is about 27.3% larger than the bare $`J`$. Perhaps surprisingly, they found however that if the bare $`g`$ factor is reduced by $`7`$% and the bare $`J`$ by $`18`$% in the $`\chi ^{}(t)`$ predicted for the Heisenberg model, this model was then in good agreement with their QMC simulations for temperatures above $`T_\mathrm{c}`$. A recent important extensive study of many finite-temperature properties of the same model using QMC simulations was carried out by Kühne and Löw. They found that for not too low temperatures, the susceptibilities for various Einstein phonon frequencies and spin-phonon coupling constants can all be scaled onto a universal curve, given by that for the uniform Heisenberg chain, using only a suitably defined effective exchange constant $`J_{\mathrm{eff}}>J`$. Contrary to the result of Ref. , they found that a rescaling of the $`g`$ factor was not necessary.
Our experimental results for $`\mathrm{NaV}_2\mathrm{O}_5`$ are not consistent with either of these theoretical studies, because as discussed below Eq. (195) above, to force-fit the Heisenberg chain $`\chi (T)`$ prediction onto the data requires an unphysically large negative value of $`\chi _0`$, as well as an unphysically large increase in $`g`$.
On the other hand, our observed $`\chi (T)`$ does not agree with the Heisenberg chain model (with a temperature-independent $`J`$), and in the next section we simultaneously model the data both above and below $`T_\mathrm{c}`$ within the context of the Heisenberg chain model with a temperature-dependent $`J`$, where we find that $`J(T)`$ above $`T_\mathrm{c}`$ is very similar in form to that deduced in the calculations of Refs. and . Thus it may be the case that the spin-phonon interaction is indeed important to determining $`\chi (T)`$ in NaV<sub>2</sub>O<sub>5</sub>, but where the effects on $`\chi (T)`$ are somewhat different than calculated in the models. In particular, the theoretical predictions may be substantially modified if phonon spectra appropriate to real materials were to be used in the calculations instead of dispersionless Einstein phonons.
### C Simultaneous modeling of the susceptibility of NaV<sub>2</sub>O<sub>5</sub> below and above $`𝑻_𝐜`$
Previous modeling of $`\chi (T)`$ of NaV<sub>2</sub>O<sub>5</sub> to extract the spin gap has usually been done at the lowest temperatures without reference to the magnitude of $`\chi `$ above $`T_\mathrm{c}`$. Here we utilize our fit to the $`\chi ^{}(t)`$ for the Heisenberg chain to extract $`J`$ above $`T_\mathrm{c}`$ from the experimental data. Clearly, since the measured $`\chi (T)`$ above $`T_\mathrm{c}`$ cannot be modeled within this framework using a temperature-independent $`J`$ as shown in the previous section, it follows that if we are to remain within this framework, $`J`$, which is then evidently an effective exchange constant incorporating additional physics of the material, must be temperature dependent. Then with $`J(T)`$ fixed, we derive the $`T`$-dependent spin gap $`\mathrm{\Delta }(T)`$ and exchange alternation parameter $`\delta (T)`$ near and below $`T_\mathrm{c}`$ directly from the measured $`\chi (T)`$ data, which has not, to our knowledge, been carried out before for any system showing a spin-dimerization transition, using our $`\chi ^{}(\alpha ,t)`$ fit function for the alternating chain.
The specific procedure we adopted for modeling the $`\chi _b(T)`$ measurement on each crystal consists of the following six steps, where we fixed $`g_b=1.972`$ in steps 3–5.
(1) The $`\chi (T)`$ from 2 to 10 K is fitted by Eqs. (38), setting $`\chi ^{\mathrm{spin}}=0`$ because of the presence of the spin gap, thereby obtaining the parameters $`\chi _0,C_{\mathrm{imp}}`$, and $`\theta `$.
(2) Using these $`\chi _0,C_{\mathrm{imp}}`$ and $`\theta `$ parameters, we solve for $`J(T)`$ for $`T60`$ K, or for $`T=50`$ K only, using our “Fit 1” function for $`\overline{\chi ^{}}(\overline{t})`$ of the Heisenberg chain, which is one end-point function of our $`\overline{\chi ^{}}(\delta ,\overline{t})`$ fit function, and fit the resulting $`J(T)`$ by a polynomial in $`T`$ for extrapolation below $`T_\mathrm{c}`$; we used the extrapolation function $`J(T)=J(0)+aT^2+bT^3`$.
(3) With this $`J(T)`$, or using $`J(50`$ K) only, we fitted $`\chi (T)`$ from 2 to 20 K, now including $`\chi ^{\mathrm{spin}}(T)`$ for the alternating-exchange chain \[i.e., using our alternating chain $`\overline{\chi ^{}}(\delta ,\overline{t})`$ fit function\] assuming a $`T`$-independent $`\delta `$ (and $`\mathrm{\Delta }`$), and obtain a new set of $`\chi _0,C_{\mathrm{imp}}`$ and $`\theta `$ parameters \[in addition to $`\delta (0)`$\].
(4) Steps 2 and 3 are repeated until convergence is achieved, which takes in practice only one additional iteration. Note that we implicitly assume that $`\chi _0,C_{\mathrm{imp}}`$ and $`\theta `$ are independent of $`T`$, i.e., that the transition(s) at $`T_\mathrm{c}`$ do not affect them.
(5) The experimentally determined molar spin susceptibility $`\chi ^{\mathrm{spin}}(T)`$ is now computed by inserting the final $`\chi _0,C_{\mathrm{imp}}`$ and $`\theta `$ fit parameters into Eq. (183). Then using the fitted $`J(T)`$ or $`J(50`$ K), the $`\delta (T)`$ is computed using our fit function $`\overline{\chi ^{}}(\delta ,\overline{t})`$ for the alternating-exchange chain by finding the root for $`\delta `$, at each data point temperature $`T`$, of
$$\chi _b^{\mathrm{spin}}(T)=\frac{N_\mathrm{A}g_b^2\mu _\mathrm{B}^2}{J_b(T)}\overline{\chi ^{}}[\delta ,\frac{k_\mathrm{B}T}{J_b(T)}].$$
(198)
(6) In a separate step not associated with the fitting procedure in steps 1–5, the spin gap $`\mathrm{\Delta }(T)`$ is computed from $`\delta (T)`$ determined in step 5 using an independently known function $`\overline{\mathrm{\Delta }^{}}(\delta )\mathrm{\Delta }(\delta )/J`$ and our $`J(T)`$ or $`J(50`$ K). We used our $`\overline{\mathrm{\Delta }^{}}(\delta )`$ fit function in Eqs. (27) for this purpose.
We measured $`\chi _b(T)`$ for nine different crystals from four different batches of NaV<sub>2</sub>O<sub>5</sub> and now present illustrative results obtained in each of the above modeling steps 2 to 4 (final iteration), 5 and 6 for three representative crystals. We will follow in graphical form the data modeling through successive steps for these three crystals to show how differences in one property between the crystals may or may not propagate through the next step(s)
of the analysis, but we present the fitting parameters for all of the crystals in Table VI.
The measured $`\chi (T)`$ data below 50 K for the three crystals are shown in Fig. 42(a), where the fits below 20 K in step 4 are shown as the solid curves with parameters in Table VI. Crystals E097A and AL1 are seen to have much lower levels of paramagnetic impurities than E083EF, as reflected in the impurity Curie constant, i.e., the magnitude of the impurity Curie-Weiss upturn at low $`T`$. By subtracting the $`\chi _0`$ and impurity Curie-Weiss terms from the data, the spin susceptibility $`\chi ^{\mathrm{spin}}(T)`$ is obtained for each crystal as shown in Fig. 42(b). These data show good consistency below $`T_\mathrm{c}`$ for the three crystals, despite the differences in the $`\chi _0`$ values, the magnitudes of the Curie-Weiss impurity term and in the $`\chi (T)`$ above $`T_\mathrm{c}`$. The $`J(T)`$ determined for the three crystals in step 2 are shown up to 300 K in Fig. 43. $`J`$ is found to decrease by $`10`$–20 % upon increasing $`T`$ from 60 to 300 K, which when $`T`$ is scaled by $`J`$ is similar to the fractional decrease predicted by Sandvik et al. due to the spin-phonon interaction. It is noteworthy that crystal E083EF, with by far the highest level of paramagnetic defects and/or impurities, also has the largest $`J(T)`$ and the largest change in $`J`$ with $`T`$.
Figures 44(a) and 44(b) show the spin dimerization parameter $`\delta (T)`$ and spin gap $`\mathrm{\Delta }(T)`$ determined for each of the three crystals in the final modeling steps 5 and 6, respectively. Several features of these data are of note. First, there is a rather large variation in the dimerization parameter, $`\delta (0)=0.028`$–0.040, between the three crystals, despite the fact that $`T_\mathrm{c}=33`$–34 K is nearly the same for the different crystals; the most impure crystal E083EF has the smallest $`\delta (0)`$, as might have been expected. Despite this variability, these $`\delta (0)`$ values are all significantly smaller than the three values reported for various samples by different groups as determined using different techniques, which are listed in Table VII along
with other related information. On the other hand, the
corresponding range of $`\mathrm{\Delta }(0)/k_\mathrm{B}=103(2)`$ K for the three crystals is fractionally much smaller than that of $`\delta (0)`$. We infer that some of the discrepancies between the $`\mathrm{\Delta }(0)`$ values in Table VII reported for NaV<sub>2</sub>O<sub>5</sub> by different groups may arise from differences in, e.g., the types of measurements which are used to determine $`\mathrm{\Delta }(0)`$ and in the different analyses of those data, rather than from different $`\mathrm{\Delta }(0)`$ values in the samples. The variability in $`\delta (0)`$ between the crystals in Fig. 44(a), compared with the lack of much variability in $`\mathrm{\Delta }(0)`$ in Fig. 44(b), evidently arises because $`\delta `$ must be combined with $`J`$ to obtain $`\mathrm{\Delta }`$, and the variations in the first two parameters must largely cancel. Thus, not surprisingly, the low-$`T`$ $`\chi ^{\mathrm{spin}}(T)`$ is governed by the spin gap $`\mathrm{\Delta }`$ and not by $`\delta `$ or $`J`$ separately.
The $`\delta (T)`$ data for our best crystals show very sharp, nearly vertical increases with decreasing $`T`$ at $`T_\mathrm{c}`$. We cannot extract a precise critical exponent $`\beta `$ from our $`\delta (T)`$ data due to the large temperature-dependent background above $`T_\mathrm{c}`$, to be discussed shortly. However, rough fits below $`T_\mathrm{c}`$ by the expression $`\delta (T)(1T/T_\mathrm{c})^\beta `$ gave $`\beta `$ values consistent with the values $`\beta =0.25(10)`$ (Ref. ) from infrared reflectivity measurements, 0.34(8) (Ref. ) from sound velocity measurements along the chain axis and 0.35(8) (Ref. ) from thermal expansion measurements along that axis. We note that these values are a factor of two larger than the value of $`0.15`$ (Ref. ) inferred from x-ray diffuse scattering measurements.
The data in Figs. 44(a) and 44(b) clearly show the existence of spin dimerization fluctuations and a spin pseudogap above $`T_\mathrm{c}`$ for each crystal, respectively, irrespective of the crystal quality as judging from the Curie-Weiss impurity term in the low-$`T`$ $`\chi (T)`$, with magnitudes just above $`T_\mathrm{c}`$ of about 20 % and 40 % of $`\delta (0)`$ and $`\mathrm{\Delta }(0)`$, respectively. This is a robust result, which was obtained
for each of the nine crystals we measured, which does not depend on the precise value of $`J`$ \[and resultant $`\chi ^{\mathrm{spin}}(T,\delta =0)`$\] or the details of how $`J`$ is determined above $`T_\mathrm{c}`$, or even on the detailed formulation of the $`\chi ^{}(\alpha ,t)`$ fit function for the alternating-exchange chain. For example, setting $`J`$ to be a constant, equal to the value at 50 K, yields nearly the same $`\mathrm{\Delta }(T)`$ near $`T_\mathrm{c}`$ as determined using a $`T`$-dependent $`J`$. Similarly, deleting the impurity Curie-Weiss term in the fit to the data above $`T_\mathrm{c}`$ changes the derived $`\chi _0`$ and $`J(T)`$ or $`J`$(50 K) somewhat as well as the detailed temperature dependence of the pseudogap $`\mathrm{\Delta }(T)`$ above $`T_\mathrm{c}`$ but has little influence on the magnitude of $`\mathrm{\Delta }`$ near $`T_\mathrm{c}`$. Further, in a previous version of the QMC and TMRG $`\chi ^{}(\alpha ,t)`$ fit function (not otherwise discussed in this paper), we did not enforce the requirement (iii) in Sec. IV E that the transformed $`\overline{\chi ^{}}(\delta ,\overline{t})`$ satisfy $`\overline{\chi ^{}}(\delta ,\overline{t})/\delta |_{\delta =0}=0`$, and the same fluctuation effects above $`T_\mathrm{c}`$ were found using that fit function as using the present one, although these fluctuations were somewhat reduced in magnitude compared to the present results. Finally, these fluctuations are observable directly in the measured $`\chi (T)`$ data in Fig. 42(a) as a rounding of the susceptibility curves above $`T_\mathrm{c}`$.
From Fig. 44, the fluctuation effects persist up to high temperatures $`T>50`$ K, although the fluctuation amplitudes decrease with increasing $`T`$. Precursor effects above $`T_\mathrm{c}`$ have been reported in x-ray diffuse scattering measurements up to $`90`$ K, in ultrasonic sound velocity and optical absorption measurements up to $`70`$ K, and in specific heat measurements up to $`40`$–50 K, so it is not surprising that spin dimerization parameter fluctuations in Fig. 44(a), and a spin pseudogap in Fig. 44(b) reflecting fluctuations in the spin gap, are found above $`T_\mathrm{c}`$.
### D Specific heat of NaV<sub>2</sub>O<sub>5</sub>
In order to correlate the magnetic effects discussed above in NaV<sub>2</sub>O<sub>5</sub> with thermal effects, we have carried out specific heat vs temperature $`C_\mathrm{p}(T)`$ measurements on the same crystal AL1, and a crystal E097 from the same batch as E097A, for which $`\chi (T)`$ data were presented and modeled above. The results from 2 K to 50 K for crystals E097 and AL1 are shown in Fig. 45(a). Over this temperature range, the $`C_\mathrm{p}(T)`$ data for the two crystals agree extremely well, except in the range 33.0–34.2 K, i.e., in the vicinity of the transitions as will be discussed shortly. The shapes of the specific heat anomalies at $`T_\mathrm{c}`$ are not mean-field-like specific heat jumps as observed in, e.g., conventional superconductors, but instead are $`\lambda `$-shaped anomalies. Thus, any attempt to define a (mean-field) “specific heat jump at $`T_\mathrm{c}`$” is fraught with ambiguity. These shapes are retained in plots of $`C_\mathrm{p}(T)/T`$ vs $`T`$ as shown in Fig. 46(a). This $`\lambda `$ shape has been observed previously, and variously attributed to fluctuation effects or a possible smeared-out first order transition. In view
of the coupled structural, charge-ordering and spin dimerization transitions at $`T_\mathrm{c}`$ in NaV<sub>2</sub>O<sub>5</sub> as discussed in the Introduction, their relative contributions to the specific heat anomalies are not clear, if indeed their contributions can be uniquely distinguished.
Expanded plots of $`C_\mathrm{p}(T)`$ and $`C_\mathrm{p}(T)/T`$ versus $`T`$, shown in Figs. 45(b) and 46(b), respectively, reveal a sharp high peak at 33.4 K for crystal AL1, which is slightly split by $`0.1`$ K in spite of the fact that the overall height of the anomaly is much larger than pre-
viously reported for any crystal of NaV<sub>2</sub>O<sub>5</sub>. Two peaks are also observed for crystal E097, at 33.4 K and 33.8 K, which are more widely separated than for crystal AL1. From Fig. 46(b), the entropy under the anomaly(ies) for each crystal is about the same (see below). Comparing these results with the $`\delta (T)`$ and $`\mathrm{\Delta }(T)`$ data in Fig. 44, the larger splitting of the $`C_\mathrm{p}(T)`$ peak for crystal E097 does not result in any major difference in the magnetic order parameter properties between the two crystals, although the transition onset is slightly rounded for crystal E097A compared to AL1. By using the Fisher relation, $`[\chi (T)T]/TC(T)`$ where $`C(T)`$ is the magnetic contribution to the specific heat, one obtains results which show the same features near $`T_\mathrm{c}`$ as does the specific heat, as shown in Fig. 45(c). Thus, careful scrutiny of the magnetic properties can reveal the fine detail observed in the specific heat near $`T_\mathrm{c}`$. In particular, this comparison suggests that both anomalies in the specific heat near $`T_\mathrm{c}`$ for each crystal are associated with and/or reflected by magnetic effects.
The splitting of the transition into two apparent transitions that we report here was previously observed in thermal expansion, but not seen in their specific heat, measurements of a crystal by Köppen et al. The detailed origin of the transition splitting, and more fundamentally whether the splitting is instrinsic to ideal crystallographically ordered NaV<sub>2</sub>O<sub>5</sub>, remain to be clarified. An essential feature that any explanation must account for is that the temperature splitting between the two transitions in a crystal varies from crystal to crystal.
##### Modeling.
In this section we will only consider the model utilized above for analyzing our $`\chi (T)`$ data, in which $`\mathrm{NaV}_2\mathrm{O}_5`$ consists, effectively, of isolated $`S=1/2`$ uniform or (below $`T_\mathrm{c}`$) alternating-exchange Heisenberg chains, where the (average) exchange constant $`J`$ shows, at most, only a smooth and relatively small change below $`T_\mathrm{c}`$. For reasons which will become clear below, unfortunately we cannot use our specific heat data to extract detailed information about the magnetic subsystem in $`\mathrm{NaV}_2\mathrm{O}_5`$. However, other types of important information about the thermodynamics will be derived using various of the theoretical results presented and discussed previously in this paper.
There have been two reports deriving the spin gap from $`C_\mathrm{p}(T)`$ data at $`T15`$ K. We first discuss the limits of this type of analysis. Using $`J(0)/k_\mathrm{B}=600`$ K, $`\delta (0)=0.040`$, and $`\mathrm{\Delta }/k_\mathrm{B}=100`$ K (see Table VI), Eqs. (64) and (65) predict that the magnetic specific heat $`C(T)`$ in the dimerized phase at low temperatures $`T(\mathrm{\Delta }/k_\mathrm{B},T_\mathrm{c})`$ is
$`C(T)=`$ $`1.0{\displaystyle \frac{\mathrm{J}}{\mathrm{mol}\mathrm{K}}}\left({\displaystyle \frac{100}{T}}\right)^{3/2}`$ (200)
$`\times \left[1+{\displaystyle \frac{T}{100}}+{\displaystyle \frac{3}{4}}\left({\displaystyle \frac{T}{100}}\right)^2\right]\mathrm{e}^{100/T},`$
with $`T`$ in units of K. Equation (200) predicts that $`C(15\mathrm{K})=0.026`$ J/mol K, which is about 40 times smaller than the observed $`C_\mathrm{p}(15\mathrm{K})1`$ J/mol K (which must therefore be due to the lattice contribution) and hence is unresolvable at such low temperatures. Within this model, we must therefore conclude that the previous estimates of the spin gap based upon modeling the low temperature specific heat were most likely artifacts of modeling the lattice specific heat. This can happen if one does not utilize the fact that the prefactor to the activated exponential term of the magnetic contribution $`C(T)`$ is not an independently adjustable parameter, but is rather determined by the spin gap itself as we have previously demonstrated and emphasized in Sec. II C 3.
A related question is whether the entropy associated with the transition(s) at $`T_\mathrm{c}`$ can be associated solely with the magnetic subsystem. The minimum possible estimate of the entropy of the transition is obtained from the $`C_\mathrm{p}(T)/T`$ vs $`T`$ data in Fig. 46(b) by drawing a horizontal line from the data at the $`C_\mathrm{p}(T)/T`$ minimum at $`35.0`$ K, just above $`T_\mathrm{c}`$, to the data that the line intersects with below $`T_\mathrm{c}`$ at $`30.6`$ K, and then computing the area between the line and the peak(s) above the line. In this way we obtain a value of 0.397 J/mol K for crystal E097 and 0.375 J/mol K for crystal AL1. On the other hand, the maximum magnetic entropy of the $`S=1/2`$ uniform chain subsystem at $`T_\mathrm{c}`$, using rough values $`J/k_\mathrm{B}=600`$ K
and $`T_\mathrm{c}=34`$ K, is roughly $`S(T_\mathrm{c})(2R/3)(k_\mathrm{B}T_\mathrm{c}/J)=0.31`$ J/mol K. Thus, the specific heat $`\lambda `$ anomaly at $`T_\mathrm{c}`$ cannot arise solely from the magnetic subsystem, since the minimum possible entropy of the transition is significantly greater than the maximum possible magnetic entropy at $`T_\mathrm{c}`$. At the least, the remaining entropy must therefore be due to the crystallographic and/or charge-ordering transitions which occur at or close to the spin dimerization transition temperature as discussed in the Introduction.
A potentially definitive and effective way to proceed from this point would be to quantitatively determine the magnetic contribution $`C(T)`$ to the measured specific heat $`C_\mathrm{p}(T)`$ from $`\chi (T)`$ at and near $`T_\mathrm{c}`$, using a relationship between $`\chi (T)`$ and $`C(T)`$ such as the Fisher relation cited above, and then compare this result with $`C_\mathrm{p}(T)`$. From a comparison of Figs. 45(b) and 45(c), it seems clear that such a relation must exist, at least for temperatures near $`T_\mathrm{c}`$, but the relationship between $`\chi (T)`$ and $`C(T)`$ near spin dimerization transitions has not yet been worked out theoretically.
In the absence of such a formulation, we proceed to estimate the change in the specific heat associated with the transition(s). In order to do this modeling, we must fit $`C_\mathrm{p}(T)`$ to higher temperatures than we have been discussing so far. The $`C_\mathrm{p}(T)`$ data from 2 to 100 K for $`\mathrm{NaV}_2\mathrm{O}_5`$ crystal E097 are shown as the open circles in Fig. 47. As noted above, except in the immediate vicinity of $`T_\mathrm{c}`$ the $`C_\mathrm{p}(T)`$ data for crystal AL1 are nearly identical with those for crystal E097 up to at least 50 K, so it will suffice to model the data for crystal E097. The four modeling steps and the assumptions we employed are as follows.
(1) We assume that critical and other order parameter fluctuations associated with the transition(s) at $`T_\mathrm{c}`$ make a negligible contribution to $`C_\mathrm{p}(T)`$ over some specified high temperature ($`TT_\mathrm{c}`$) range. By subtracting the known magnetic contribution $`C(T)`$ due to isolated Heisenberg chains \[obtained using our fit function for $`C(k_\mathrm{B}T/J)`$\] in this temperature range, we obtain the background lattice contribution $`C^{\mathrm{lat}}(T)`$ in the high temperature region. Also, since we have shown that $`C(T)`$ is negligible for $`T15`$ K, the measured $`C_\mathrm{p}(T)`$ in this $`T`$ range is assumed to be identical to $`C^{\mathrm{lat}}(T)`$ at these temperatures (we again neglect the possible but unknown specific heats associated with possible order parameter fluctuations in this range). Thus we obtain background lattice specific heats $`C^{\mathrm{lat}}(T)`$ in high and low temperature ranges which are assumed unaffected by the transition(s) and associated order parameter fluctuations.
(2) We interpolate between the $`C^{\mathrm{lat}}(T)`$ determined in step 1 in the low- and high-temperature ranges to obtain, in the intermediate temperature range, what $`C^{\mathrm{lat}}(T)`$ would have been in the absence of the transition(s) and associated order parameter fluctuations.
(3) We add the $`C(T)`$ for isolated chains, used in step 1, back to the $`C^{\mathrm{lat}}(T)`$ derived in step 2 over the full temperature range of the measurements. This is the total background specific heat that would have occurred in the absence of the transition(s) and associated order parameter fluctuations. Then we subtract the total calculated background specific heat from the measured $`C_\mathrm{p}(T)`$ data. This difference $`\mathrm{\Delta }C(T)`$ should hopefully be a reasonable estimate of the change in the specific heat associated with the transition and order parameter fluctuations, including all lattice, charge and spin contributions. $`\mathrm{\Delta }C(T)`$ must go to zero, by construction, at temperatures above the lower end of the high temperature region fitted in step 2.
(4) Finally we integrate $`\mathrm{\Delta }C/T`$ with $`T`$ up to and beyond $`T_\mathrm{c}`$ to obtain the change in entropy $`\mathrm{\Delta }S(T)`$ associated with the transition and order parameter fluctuations. $`\mathrm{\Delta }S(T)`$ must become constant, by construction, at temperatures above the lower end of the high temperature region fitted in step 2.
In the following we will present and discuss the results in each of the four steps of our modeling program described above.
##### Step 1.
Here we first use our $`C(T)`$ fit function for the numerical $`C(T)`$ data, which was given in Eqs. (18), to extract $`C^{\mathrm{lat}}(T)`$ in the high-temperature region above $`T_\mathrm{c}`$. For consistency with our analysis of the susceptibility in Sec. VIII C, we use the temperature-dependent $`J(T)`$ derived in that section for crystal E097A (see Fig. 43) when computing $`C(T)`$. The background $`C(T)`$ thus estimated for crystal E097, i.e., the values which would have been observed if no transition(s) at $`T_\mathrm{c}`$ or associated order parameter fluctuations had occurred, is shown in Fig. 48. Comparison of these data with the measured $`C_\mathrm{p}(T)`$ data in Figs. 45 and 47 shows that this $`C(T)`$ is a small, but non-negligible ($`1`$%), fraction of $`C_\mathrm{p}(T)`$ above $`T_\mathrm{c}`$. On the other hand, $`C(T)`$ is much larger than the observed $`C_\mathrm{p}(T)`$ at low temperatures, because in this temperature range $`CT`$ whereas $`C_\mathrm{p}(T)C^{\mathrm{lat}}(T)T^3`$. The $`C^{\mathrm{lat}}(T)=C_\mathrm{p}(T)C(T)`$ in the high temperature (60–100 K) region is shown in Fig. 49, together with
$`C^{\mathrm{lat}}(T)C_\mathrm{p}(T)`$ in the low temperature (2–15 K) region.
##### Step 2.
In this step we must interpolate $`C^{\mathrm{lat}}(T)`$ between the low- and high-temperature regions, i.e., in a broad temperature range spanning the transition region. The best way to do this would be to determine $`C^{\mathrm{lat}}(T)`$ directly from $`C_\mathrm{p}(T)`$ measurements on a suitably chosen reference compound, but such measurements have not yet been done. At first sight, a physically realistic possibility might be to interpolate the low and high temperature $`C^{\mathrm{lat}}(T)`$ data using the Debye specific heat function; however, this method is questionable because the Debye temperature $`\mathrm{\Theta }_\mathrm{D}`$ in real materials can be rather strongly temperature dependent within the temperature range of interest here. The Debye function for the molar lattice specific heat at constant volume $`C^{\mathrm{Debye}}(T)`$ is given by
$$C^{\mathrm{Debye}}(T)=9rR\left(\frac{T}{\mathrm{\Theta }_\mathrm{D}}\right)^3_0^{\mathrm{\Theta }_\mathrm{D}/T}\frac{x^4\mathrm{e}^x}{(\mathrm{e}^x1)^2}𝑑x,$$
(201)
where $`r`$ is the number of atoms per formula unit ($`r=8`$ here) and $`R`$ is the molar gas constant. We attempted to fit our $`C^{\mathrm{lat}}(T)`$ data for the temperature ranges 2–15 K and 40–100 K to 80–100 K by Eq. (201). The fits parametrized the data very poorly. We obtained a more reasonable fit by allowing $`r`$ to be a fitting parameter, yielding a fitted value $`r4`$, but the data were still poorly fitted, due to too much curvature in the Debye function in the high temperature region. Therefore, we were led to interpolating between the low- and high-temperature regions using a simple polynomial interpolation function.
To obtain the background lattice specific heat interpolation function, we fitted the combined $`C^{\mathrm{lat}}(T)`$ data (a total of 141 data points) in the low and high temperature ranges 2–15 K and 60–101 K, respectively, by polynomials of the form
$$C^{\mathrm{lat}}(T)=\underset{n=3}{\overset{n^{\mathrm{max}}}{}}c_nT^n.$$
(202)
The minimum summation index $`n=3`$ is set by the expected Debye low-temperature $`T^3`$ behavior of the lattice specific heat. The maximum value $`n^{\mathrm{max}}`$ was varied to see how the fit parameters and variance changed. In addition, for checking the final fits we fitted the $`C^{\mathrm{lat}}(T)`$ data in the 2–15 K low-$`T`$ range together with $`C^{\mathrm{lat}}(T)`$ data in a high-temperature range varying from 40–101 K to 90–101 K. We found that the most stable fits were for $`n^{\mathrm{max}}=7`$ and 8. For both values, the fit did not visibly change when the lower limit of the upper temperature range of the fitted data was varied from 60 to 70 K. We chose to use the fit for $`n^{\mathrm{max}}=7`$ because in this case the fit was also stable for lower limits of 50 and 80 K. This stability allows one to be confident that the interpolation of the fit between the fitted low- and high-temperature ranges is an accurate representation of the background lattice specific heat in the interpolated intermediate temperature range. The fit for the temperature ranges 2–15 K and 60–101 K is shown as the solid curve in Fig. 49. The absolute rms deviation of this fit from the fitted data is quite small, $`\sigma _{\mathrm{rms}}=0.046`$ J/mol K. The curve over the full temperature range 2–101 K represents the background lattice specific heat $`C^{\mathrm{lat}}(T)`$ expected in the absence of any transitions or order parameter fluctuations.
##### Step 3.
Adding the magnetic background specific heat contribution $`C(T)`$ obtained in step 1 to the lattice background specific heat contribution $`C^{\mathrm{lat}}(T)`$ obtained in step 2 gives the total background specific heat, which is plotted as the solid curve in Fig. 47. We reiterate that this background is interpreted as the specific heat that would have been observed had the transition(s) at $`T_\mathrm{c}`$
and associated order parameter fluctuations not occurred. The difference $`\mathrm{\Delta }C`$ between the measured $`C_\mathrm{p}(T)`$ and the total background specific heat is plotted versus temperature in Fig. 50(a). As would have been qualitatively anticipated, $`\mathrm{\Delta }C`$ is negative below about 16 K due to the loss of magnetic specific heat at low temperatures arising from the opening of the spin gap at $`T_\mathrm{c}`$. This negative $`\mathrm{\Delta }C`$ does not arise from a problem in our polynomial interpolation $`C^{\mathrm{lat}}(T)`$ fit function or from our $`C(T)`$ function; these functions are both positive for all $`T>0`$. Since the magnetic background contribution is proportional to $`T`$ and the lattice background contribution \[which is assumed not to change below 15 K due to the occurrence of the transition(s)\] is proportional to $`T^3`$ at low $`T`$, opening a spin gap at $`T_\mathrm{c}`$ must necessarily lead to a negative $`\mathrm{\Delta }C`$ at sufficiently low temperatures since the magnetic contribution then becomes exponentially small there.
##### Step 4.
Finally, we can compute the change $`\mathrm{\Delta }S`$ in the total entropy of the system versus temperature due to the transition(s) and associated order parameter fluctuations by integrating $`\mathrm{\Delta }C(T)`$ from step 3 according to $`\mathrm{\Delta }S(T)=_0^T[\mathrm{\Delta }C(T)/T]𝑑T`$. The result is shown in Fig. 50(b). The entropy change is negative below about 22 K, due to the loss of magnetic entropy at low temperatures associated with the loss of magnetic specific heat as just discussed. From conservation of magnetic entropy, this lost entropy must reappear at higher temperatures.
By construction, step 2 requires that $`\mathrm{\Delta }C(T>60`$ K) = 0 and consequently $`\mathrm{\Delta }S(T>60`$ K) = const. This requirement is not desirable, but we had to enforce it to ensure that the $`\mathrm{\Delta }C(T)`$ and $`\mathrm{\Delta }S(T)`$ derived at lower temperatures were accurate. Since the effects of the order parameter fluctuations are likely to continue to be present at temperatures higher than 60 K, the $`\mathrm{\Delta }C(T)`$ and $`\mathrm{\Delta }S(T)`$ at temperatures at and near 60 K in Fig. 50 are lower limits.
The net change in the entropy at 60 K in Fig. 50(b) due to the occurrence of the transition(s) at $`T_\mathrm{c}34`$ K and associated order parameter fluctuations above and below $`T_\mathrm{c}`$ is $`\mathrm{\Delta }S(60`$ K) = 2.28 J/mol K. This is far larger than the maximum possible change $`\mathrm{\Delta }S_{\mathrm{mag}}^{\mathrm{max}}=0.556`$ J/mol K in the magnetic entropy at this temperature obtained from Fig. 48, where this value is just the maximum possible entropy of the magnetic subsystem at this temperature, confirming our qualitative conclusion above based on very rough arguments. In particular, our quantitative analysis indicates that at least 76% of the entropy change at 60 K must arise from the lattice and charge degrees of freedom, and only a minor fraction ($`<24`$ %) from the magnetic degrees of freedom. Similarly, at $`T_\mathrm{c}=33.7`$ K, we obtain $`\mathrm{\Delta }S=1.38`$ J/mol K and $`\mathrm{\Delta }S_{\mathrm{mag}}^{\mathrm{max}}=0.311`$ J/mol K, yielding $`\mathrm{\Delta }S_{\mathrm{mag}}^{\mathrm{max}}/\mathrm{\Delta }S23`$ % at $`T_\mathrm{c}`$.
As a closing remark for this section, it is clear from Fig. 50 and the discussion in the above two paragraphs that $`\mathrm{\Delta }C`$ and $`\mathrm{\Delta }S`$ do not saturate to their respective high temperature limiting values until a temperature of at least 60 K is reached, which is almost twice $`T_\mathrm{c}`$. The present analysis of the thermal behavior of $`\mathrm{NaV}_2\mathrm{O}_5`$ thus lends strong support to our independent analysis and interpretation of our magnetic susceptibility data for this compound in Sec. VIII C.
## IX Summary and Concluding Discussion
We have shown that the high-accuracy numerical Bethe ansatz calculations of the magnetic susceptibility $`\chi ^{}(t)`$ for the $`S=1/2`$ uniform Heisenberg chain by Klümper and Johnston are in excellent agreement with the theory of Lukyanov over 18 decades of temperature at low temperatures. An independent high precision empirical fit to these data was obtained over 25 decades of temperature which we found useful to determine the accuracy of our TMRG $`\chi ^{}(t)`$ calculations. The magnetic specific heat data for the uniform chain at very low temperatures was also compared with the theoretical predictions of Lukyanov, and extremely good agreement was found over many decades in temperature. We formulated an empirical fit function for these data which is highly accurate over a temperature range spanning 25 orders of magnitude; the infinite temperature entropy calculated using this fit function is within 8 parts in $`10^8`$ of the exact value. We used both of the above fit functions to model our respective experimental data for $`\mathrm{NaV}_2\mathrm{O}_5`$ in later sections of the paper. We expect that they will be useful to other theorists and experimentalists as well.
We have carried out extensive QMC simulations and TMRG calculations of $`\chi ^{}(\alpha ,t)`$ for the spin $`S=1/2`$ antiferromagnetic alternating-exchange Heisenberg chain for reduced temperatures $`tk_\mathrm{B}T/J_1`$ from 0.002 to 10 and alternation parameters $`\alpha J_2/J_1`$ from 0.05 to 1, where $`J_1(J_2)`$ is the larger (smaller) of the two alternating exchange constants. An accurate global two-dimensional ($`\alpha ,t`$) fit to these combined data was obtained, constrained by the fitting parameters for the accurately known $`\chi ^{}(t)`$ for the $`\alpha `$ parameter end points, the dimer ($`\alpha =0`$) and the uniform chain ($`\alpha =1`$), resulting in an accurate fit function over the entire range $`0\alpha 1`$ of the alternation parameter. Our fit function incorporates the first four terms of the exact high-temperature series expansion in powers of $`1/t`$, which allows accurate extrapolation to arbitrarily high temperatures. This function should prove useful for many applications including the modeling of experimental $`\chi (T)`$ data as we have shown.
Our $`\chi ^{}(\alpha ,t)`$ fit function for the alternating chain can be easily transformed (as we have done) into an equivalent fit function $`\overline{\chi ^{}}(\delta ,\overline{t})`$ in the two variables $`\delta (J_1J_2)/(2J)`$ and $`\overline{t}k_\mathrm{B}T/J`$, where the average exchange constant is $`J=(J_1+J_2)/2`$. This is a more appropriate function for analyzing experimental $`\chi (T)`$ data for $`S=1/2`$ Heisenberg chain compounds showing dimerization transitions (such as a spin-Peierls transition) which result in an alternating-exchange chain with a small value of $`\delta `$ at low temperatures. Once $`J`$ has been determined by fitting our function for $`\delta =0`$ to the experimentally determined spin susceptibility $`\chi ^{\mathrm{spin}}(T)`$ data above the transition temperature, the alternation parameter $`\delta `$ is uniquely determined by our fit function at each temperature below the transition temperature from the value of $`\chi ^{\mathrm{spin}}`$ at that temperature. One can then find the spin gap $`\mathrm{\Delta }(T)`$ using an independently known $`\overline{\mathrm{\Delta }^{}}(\delta )`$.
Our QMC and TMRG data and fit for $`\chi ^{}(\alpha ,t)`$ are in good agreement with previous calculations based on exact diagonalization of the nearest neighbor Heisenberg Hamiltonian for short chains with $`\alpha =0.2`$, 0.4, 0.6, 0.7, and 0.8, extrapolated to the thermodynamic limit, by Barnes and Riera. However, the numerical and analytical theoretical predictions of Bulaevskii, which have been used extensively in the past by experimentalists to model their $`\chi (T)`$ data for weakly-dimerized chain compounds, are found to be in poor agreement with our results and should be abandoned for such use in favor of our fit function. Similarly, the previously used fit function for the Bonner-Fisher calculation of $`\chi ^{}(t)`$ for the uniform chain ($`\alpha =1`$) should be replaced by one of our two fit functions for the most accurate calculation to date of $`\chi ^{}(t)`$ for the uniform chain.
An important theoretical issue in the study of the alternating exchange chain is how the spin gap $`\overline{\mathrm{\Delta }^{}}(\delta )`$ evolves as the uniform chain limit is approached ($`\delta 0,\alpha 1`$). We formulated a fit function for the temperature dependence of our TMRG susceptibility $`\chi ^{}(\alpha ,t)`$ calculations at low temperatures, which was used to extract the dependence $`\overline{\mathrm{\Delta }^{}}(\delta )`$ in this regime. We find that the asymptotic critical regime is not entered until, at least, $`\delta 0.005`$ ($`\alpha 0.99`$). We compared our spin gap data with many literature data. We formulated a fit function for our spin gap data together with those of Barnes, Riera, and Tennant which quite accurately covers the entire range $`0\delta 1`$.
In the remainder of this paper, we showed how the above theoretical results could be used to obtain detailed information about real systems. As a specific illustration, we carried out a detailed analysis of our experimental $`\chi (T)`$ and specific heat $`C_\mathrm{p}(T)`$ data for $`\mathrm{NaV}_2\mathrm{O}_5`$ crystals. This compound shows a transition to a spin dimerized state below the transition temperature $`T_\mathrm{c}34`$ K. We used one of our two $`\chi ^{}(t)`$ fit functions for the uniform Heisenberg chain to model the $`\chi (T)`$ above $`T_\mathrm{c}`$, where we found that the experimental $`\chi (T)`$ is not in quantitative agreement with the prediction for the uniform Heisenberg chain. A model incorporating a mean-field ferromagnetic interchain coupling between quantum $`S=1/2`$ Heisenberg chains fits the experimental data very well with reasonable parameters. It remains to be seen whether the inelastic neutron scattering measurements of the magnon dispersion relations are consistent with our derived intrachain and interchain exchange constants.
In an alternate description, we modeled the deviation in the measured $`\chi (T)`$ of NaV<sub>2</sub>O<sub>5</sub> above 60 K $`>T_\mathrm{c}`$ from the Heisenberg chain model (with fixed exchange constant $`J`$) as due to a temperature-dependent $`J`$. We found that this $`J`$ decreases with increasing $`T`$ up to 300 K in a manner very similar to $`J_{\mathrm{eff}}(T)`$ predicted by Sandvik, Singh and Campbell and Kühne and Löw for the spin-Peierls chain. Our $`J(T)`$ cannot however be compared directly with their $`J_{\mathrm{eff}}(T)`$ because the two quantities are defined differently. They found that by defining an appropriate effective exchange constant $`J_{\mathrm{eff}}`$, their resulting susceptibility $`\chi (k_\mathrm{B}T/J_{\mathrm{eff}})`$ is universal at the higher temperatures for various Einstein phonon frequencies and spin-phonon coupling constants. This function agrees well with the $`\chi (k_\mathrm{B}T/J)`$ for the $`S=1/2`$ AF uniform Heisenberg chain at these temperatures. As we discussed, these $`\chi (T)`$ calculations are not applicable to NaV<sub>2</sub>O<sub>5</sub>, possibly because the calculations do not incorporate realistic phonon spectra.
Below $`T_\mathrm{c}`$, we used the $`J(T)`$ extrapolated from above 60 K and our global $`\chi ^{}(\alpha ,t)`$ fit function for the alternating Heisenberg chain to determine the temperature-dependent alternation parameter $`\delta (T)`$, and then the spin gap $`\mathrm{\Delta }(T)`$ from $`\delta (T)`$, directly from the $`\chi (T)`$ data. We find that the $`\mathrm{\Delta }(0)/k_\mathrm{B}`$ values for nine single crystals of NaV<sub>2</sub>O<sub>5</sub> are in the range 103(2) K. This result is in agreement, within the errors, with many previous analyses of data from various types of measurements for this compound by other groups. However, our values of $`\delta (0)=0.034(6)`$ for various crystals are significantly smaller than previous estimates. We note that the two estimates with $`\delta (0)0.1`$ in Table VII were obtained using Bulaevskii’s theory for the alternating-exchange chain, which we have shown is not accurate at low temperatures in the relevant alternation parameter range.
The dispersion of two one-magnon branches perpendicular to the chains observed in the neutron scattering measurements has been recently explained quantitatively by Gros and Valenti assuming that a zig-zag charge ordering transition occurs at $`T_\mathrm{c}`$. They also predict that $`\delta (0)0.034`$. This is within our range of $`\delta (0)`$ values in spite of the fact that we assumed that $`J(T)`$ is either constant or increases slightly with decreasing $`T`$ below $`T_\mathrm{c}`$, contrary to their prediction that $`J`$ decreases below $`T_\mathrm{c}`$. Gros and Valenti made no predictions for $`\chi (T)`$, $`\delta (T)`$, $`\mathrm{\Delta }(T)`$ or $`C(T)`$, so comparisons with our results for these quantities are not possible. We note that Klümper, Raupach, and Schönfeld obtained a good fit to the $`\chi (T)`$ data below $`T_\mathrm{c}`$ for the spin-Peierls compound CuGeO<sub>3</sub> within the context of a spin-Peierls model containing frustrating second-neighbor interactions and static spin-phonon coupling.
We discovered that $`\mathrm{\Delta }(T)`$ \[and $`\delta (T)`$\] of NaV<sub>2</sub>O<sub>5</sub> does not go to zero at $`T_\mathrm{c}`$, indicating the existence of a spin pseudogap above $`T_\mathrm{c}`$ with a large magnitude just above $`T_\mathrm{c}`$ of $`40`$% of $`\mathrm{\Delta }(0)`$; the pseudogap is present up to at least 50 K with a magnitude decreasing with increasing $`T`$ above $`T_\mathrm{c}`$. To our knowledge, this pseudogap has not been reported previously, and there are as yet no theoretical predictions for the magnitude or temperature dependence of this pseudogap. The pseudogap is strongly reminiscent of the spin pseudogap derived by one of us using $`\chi (T)`$ measurements above the transition temperature of inorganic quasi-one-dimensional charge density wave compounds, as predicted theoretically by Lee, Rice, and Anderson long before those observations were made. Similar to that case, in the present system one may think of the pseudogap as the rms fluctuation in the spin gap above $`T_\mathrm{c}`$, with an associated reduction in the magnon density of states at low energy. In this interpretation, the pseudogap in NaV<sub>2</sub>O<sub>5</sub> should be observable in high resolution quasielastic neutron scattering and other spectroscopic measurements probing the low energy magnetic excitations.
Finally, we carried out an extensive modeling study of our specific heat data for NaV<sub>2</sub>O<sub>5</sub> crystals, using the same model that we used to analyze our susceptibility data. The most important part of this study is that we have been able to determine a limit on the relative contributions of the magnetic and lattice/charge degrees of freedom to the entropy associated with the transition(s) at $`T_\mathrm{c}`$. We find that at least 77 % of the change in the entropy at $`T_\mathrm{c}`$ must arise from the lattice and/or charge degrees of freedom, to which the spin degrees of freedom must of course be coupled, and that the spin degrees of freedom themselves contribute less than 23 % of this entropy change. Our results also indicate that order parameter fluctuation effects are important in the specific heat up to at least 60 K, strongly confirming the above similar and independent conclusion based on our modeling of our magnetic susceptibility data for the same crystals.
## Acknowledgments
We thank E. Brücher and C. Lin for help with the sample preparation, and C. Song for assistance with Laue x-ray diffraction measurements. We are grateful to S. Eggert and T. Barnes for providing the numerical $`\chi ^{}(t)`$ calculation results for the uniform chain in Ref. and the alternating chain in Ref. , respectively, and to D. Poilblanc and G. S. Uhrig for sending us their $`\overline{\mathrm{\Delta }^{}}(\delta )`$ data in Refs. and , respectively. We are grateful to M. Greven and X. Zotos for helpful discussions, and to A. A. Zvyagin for helpful correspondence. One of us (D.C.J.) thanks the Max-Planck-Institut für Festkörperforschung, Stuttgart, where this work was started, for kind hospitality. Ames Laboratory is operated for the U.S. Department of Energy by Iowa State University under Contract No. W-7405-Eng-82. The work at Ames was supported by the Director for Energy Research, Office of Basic Energy Sciences. The QMC program was written in C++ using a parallelizing Monte Carlo library developed by one of the authors. The QMC simulations by M.T. were performed on the Hitachi SR2201 massively parallel computer of the University of Tokyo and on the IBM SP-2 of the Competence Center for Computational Chemistry of ETH Zürich. X.W. acknowledges Swiss National Funding Grant No. 20-49486.96. A.K. acknowledges financial support by the Deutsche Forschungsgemeinschaft under Grant No. Kl 645/3 and support by the research program of the Sonderforschungsbereich 341, Köln-Aachen-Jülich.
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# Theorem 2.1
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# Rigidity of measurable structure for ℤ^𝑑–actions by automorphisms of a torus
## 1. Introduction; description of results
In the course of the last decade various rigidity properties have been found for two different classes of actions by higher–rank abelian groups: on the one hand, certain Anosov and partially hyperbolic actions of $`^d`$ and $`^d,d2`$, on compact manifolds () and, on the other, actions of $`^d,d2`$, by automorphisms of compact abelian groups (cf. e.g. ). Among these rigidity phenomena is a relative scarcity of invariant measures which stands in contrast with the classical case $`d=1`$ ().
In this paper we make the first step in investigating a different albeit related phenomenon: rigidity of the measurable orbit structure with respect to the natural smooth invariant measure.
In the classical case of actions by $``$ or $``$ there are certain natural classes of measure–preserving transformations which possess such rigidity: ergodic translations on compact abelian groups give a rather trivial example, while horocycle flows and other homogeneous unipotent systems present a much more interesting one . In contrast to such situations, individual elements of the higher–rank actions mentioned above are Bernoulli automorphisms. The measurable orbit structure of a Bernoulli map can be viewed as very “soft”. Recall that the only metric invariant of Bernoulli automorphisms is entropy (); in particular, weak isomorphism is equivalent to isomorphism for Bernoulli maps since it implies equality of entropies. Furthermore, description of centralizers, factors, joinings and other invariant objects associated with a Bernoulli map is impossible in reasonable terms since each of these objects is huge and does not possess any discernible structure.
In this paper we demonstrate that some very natural actions of $`^d,d2`$, by Bernoulli automorphisms display a remarkable rigidity of their measurable orbit structure. In particular, isomorphisms between such actions, centralizers, and factor maps are very restricted, and a lot of algebraic information is encoded in the measurable structure of such actions (see Section 5).
All these properties occur for broad subclasses of both main classes of actions of higher–rank abelian groups mentioned above: Anosov and partially hyperbolic actions on compact manifolds, and actions by automorphisms of compact abelian groups. However, at present we are unable to present sufficiently definitive general results due to various difficulties of both conceptual and technical nature. Trying to present the most general available results would lead to cumbersome notations and inelegant formulations. To avoid that we chose to restrict our present analysis to a smaller class which in fact represents the intersection of the two, namely the actions of $`^d,d2`$, by automorphisms of the torus. Thus we study the measurable structure of such actions with respect to Lebesgue (Haar) measure from the point of view of ergodic theory.
Our main purpose is to demonstrate several striking phenomena by means of applying to specific examples general rigidity results which are presented in Section 5 and are based on rigidity of invariant measures developed in (see for further results along these lines including rigidity of joinings). Hence we do not strive for the greatest possible generality even within the class of actions by automorphisms of a torus. The basic algebraic setup for irreducible actions by automorphisms of a torus is presented in Section 3. Then we adapt further necessary algebraic preliminaries to the special but in a sense most representative case of Cartan actions, i.e. to $`^{n1}`$–actions by hyperbolic automorphisms of the $`n`$–dimensional torus (see Section 4).
The role of entropy for a smooth action of a higher–rank abelian group $`G`$ on a finite-dimensional manifold is played by the entropy function on $`G`$ whose values are entropies of individual elements of the action (see Section 2.2 for more details) which is naturally invariant of isomorphism and also of weak isomorphism and is equivariant with respect to a time change.
In Section 6 we produce several kinds of specific examples of actions by ergodic (and hence Bernoulli) automorphisms of tori with the same entropy function. These examples provide concrete instances when general criteria developed in Section 5 can be applied. Our examples include:
* actions which are not weakly isomorphic (Section 6.1),
* actions which are weakly isomorphic but not isomorphic, such that one action is a maximal action by Bernoulli automorphisms and the other is not (Section 6.2),
* weakly isomorphic, but nonisomorphic, maximal actions (Section 6.3).
Once rigidity of conjugacies is established, examples of type (i) appear in a rather simple–minded fashion: one simply constructs actions with the same entropy data which are not isomorphic over $``$. This is not surprising since entropy contains only partial information about eigenvalues. Thus one can produce actions with different eigenvalue structure but identical entropy data.
Examples of weakly isomorphic but nonisomorphic actions are more sophisticated. We find them among Cartan actions (see Section 4). The centralizer of a Cartan action in the group of automorphisms of the torus is (isomorphic to) a finite extension of the acting group, and in some cases Cartan actions isomorphic over $``$ may be distinguished by looking at the index of the group in its centralizer (type (ii); see Examples 2a and 2b). The underlying cause for this phenomenon is the existence of algebraic number fields $`K=(\lambda )`$, where $`\lambda `$ is a unit, such that the ring of integers $`𝒪_K[\lambda ]`$. In general finding even simplest possible examples for $`n=3`$ involves the use of data from algebraic number theory and rather involved calculations. For examples of type (ii) one may use some special tricks which allow to find some of these and to show nonisomorphism without a serious use of symbolic manipulations on a computer.
A Cartan action $`\alpha `$ of $`^{n1}`$ on $`𝕋^n`$ is called maximal if its centralizer in the group of automorphisms of the torus is equal to $`\alpha (^{n1})\times \{\pm \mathrm{Id}\}`$. A maximal Cartan action turns out to me maximal in the above sense: it cannot be extended to any action of a bigger abelian group by Bernoulli automorphisms.
Examples of maximal Cartan actions isomorphic over $``$ but not isomorphic (type (iii)) are the most remarkable. Conjugacy over $``$ guarantees that the actions by automorphisms of the torus $`𝕋^n`$ arising from their centralizers are weakly isomorphic with finite fibres. The mechanism providing obstructions for algebraic isomorphism in this case involves the connection between the class number of an algebraic number field and $`GL(n,)`$–conjugacy classes of matrices in $`SL(n,)`$ which have the same characteristic polynomial (see Example 3). In finding these examples the use of computational number–theoretic algorithms (which in our case were implemented via the Pari-GP package) has been essential.
One of our central conclusions is that for a broad class of actions of $`^d,d2`$, (see condition $`()`$ in Section 2.2) the conjugacy class of the centralizer of the action in the group of affine automorphisms of the torus is an invariant of measurable conjugacy. Let $`Z_{\mathrm{𝑚𝑒𝑎𝑠}}(\alpha )`$ be the centralizer of the action $`\alpha `$ in the group of measurable automorphisms. As it turns out in all our examples but Example 3b, the conjugacy class of the pair $`(Z_{\mathrm{𝑚𝑒𝑎𝑠}}(\alpha ),\alpha )`$ is a distinguishing invariant of the measurable isomorphism. Thus, in particular, Example 3b shows that there are weakly isomorphic, but nonisomorphic actions for which the affine and hence the measurable centralizers are isomorphic as abstract groups.
We would like to acknowledge a contribution of J.-P. Thouvenot to the early development of ideas which led to this paper. He made an important observation that rigidity of invariant measures can be used to prove rigidity of isomorphisms via a joining construction (see Section 5.1).
## 2. Preliminaries
### 2.1. Basic ergodic theory
Any invertible (over $``$) integral $`n\times n`$ matrix $`AM(n,)GL(n,)`$ determines an endomorphism of the torus $`𝕋^n=^n/^n`$ which we denote by $`F_A`$. Conversely, any endomorphism of $`𝕋^n`$ is given by a matrix from $`AM(n,)GL(n,)`$. If, in addition, $`detA=\pm 1`$, i.e. if $`A`$ is invertible over $``$, then $`F_A`$ is an automorphism of $`𝕋^n`$ (the group of all such $`A`$ is denoted by $`GL(n,)`$). The map $`F_A`$ preserves Lebesgue (Haar) measure $`\mu `$; it is ergodic with respect to $`\mu `$ if and only if there are no roots of unity among the eigenvalues of $`A`$, as was first pointed out by Halmos (). Furthermore, in this case there are eigenvalues of absolute value greater than one and $`(F_A,\lambda )`$ is an exact endomorphism. If $`F_A`$ is an automorphism it is in fact Bernoulli (). For simplicity we will call such a map $`F_A`$ an ergodic toral endomorphism (respectively, automorphism, if $`A`$ is invertible). If all eigenvalues of $`A`$ have absolute values different from one we will call the endomorphism (automorphism) $`F_A`$ hyperbolic.
When it does not lead to a confusion we will not distinguish between a matrix $`A`$ and corresponding toral endomorphism $`F_A`$.
Let $`\lambda _1,\mathrm{},\lambda _n`$ be the eigenvalues of the matrix $`A`$, listed with their multiplicities. The entropy $`h_\mu (F_A)`$ of $`F_A`$ with respect to Lebesgue measure is equal to
$$\underset{\{i:|\lambda _i|>1\}}{}\mathrm{log}|\lambda _i|.$$
In particular, entropy is determined by the conjugacy class of the matrix $`A`$ over $``$ (or over $``$). Hence all ergodic toral automorphisms which are conjugate over $``$ are measurably conjugate with respect to Lebesgue measure.
Classification, up to a conjugacy over $``$, of matrices in $`SL(n,)`$, which are irreducible and conjugate over $``$ is closely related to the notion of class number of an algebraic number field. A detailed discussion relevant to our purposes appears in Section 4.2. Here we only mention the simplest case $`n=2`$ which is not directly related to rigidity. In this case trace determines conjugacy class over $``$ and, in particular, entropy. However if the class number of the corresponding number field is greater than one there are matrices with the given trace which are not conjugate over $``$. This algebraic distinctiveness is not reflected in the measurable structure: in fact, in the case of equal entropies the classical Adler–Weiss construction of the Markov partition in yields metric isomorphisms which are more concrete and specific than in the general Ornstein isomorphism theory and yet not algebraic.
### 2.2. Higher rank actions
Let $`\alpha `$ be an action by commuting toral automorphisms given by integral matrices $`A_1,\mathrm{},A_d`$. It defines an embedding $`\rho _\alpha :^dGL(n,)`$ by
$$\rho _\alpha ^𝐧=A_1^{n_1}\mathrm{}A_d^{n_d},$$
where $`𝐧=(n_1,\mathrm{},n_d)^d`$, and we have
$$\alpha ^𝐧=F_{\rho _\alpha ^𝐧}.$$
Similarly, we write $`\rho _\alpha :_+^dM(n,)GL(n,)`$ for an action by endomorphisms. Conversely, any embedding $`\rho :^dGL(n,)`$ (respectively, $`\rho :_+^dM(n,)GL(n,))`$ defines an action by automorphisms (respectively, endomorphisms) of $`𝕋^n`$ denoted by $`\alpha _\rho `$.
Sometimes we will not explicitly distinguish between an action and the corresponding embedding, e.g. we may talk about “the centralizer of an action in $`GL(n,)`$” etc.
###### Definitions.
Let $`\alpha `$ and $`\alpha ^{}`$ be two actions of $`^d`$ ($`_+^d`$) by automorphisms (endomorphisms) of $`𝕋^n`$ and $`𝕋^n^{}`$, respectively. The actions $`\alpha `$ and $`\alpha ^{}`$ are measurably (or metrically, or measure–theoretically) isomorphic (or conjugate) if there exists a Lebesgue measure–preserving bijection $`\phi :𝕋^n𝕋^n^{}`$ such that $`\phi \alpha =\alpha ^{}\phi `$.
The actions $`\alpha `$ and $`\alpha ^{}`$ are measurably isomorphic up to a time change if there exist a measure–preserving bijection $`\phi :𝕋^n𝕋^n^{}`$ and a $`CGL(d,)`$ such that $`\phi \alpha C=\alpha ^{}\phi `$.
The action $`\alpha ^{}`$ is a measurable factor of $`\alpha `$ if there exists a Lebesgue measure–preserving transformation $`\phi :𝕋^n𝕋^n^{}`$ such that $`\phi \alpha =\alpha ^{}\phi `$. If, in particular, $`\phi `$ is almost everywhere finite–to–one, then $`\alpha ^{}`$ is called a finite factor or a factor with finite fibres of $`\alpha `$.
Actions $`\alpha `$ and $`\alpha ^{}`$ are weakly measurably isomorphic if each is a measurable factor of the other.
A joining between $`\alpha `$ and $`\alpha ^{}`$ is a measure $`\mu `$ on $`𝕋^n\times 𝕋^n^{}=𝕋^{n+n^{}}`$ invariant under the Cartesian product action $`\alpha \times \alpha ^{}`$ such that its projections into $`𝕋^n`$ and $`𝕋^n^{}`$ are Lebesgue measures. As will be explained in Section 5, conjugacies and factors produce special kinds of joinings.
These measure–theoretic notions have natural algebraic counterparts.
###### Definitions.
The actions $`\alpha `$ and $`\alpha ^{}`$ are algebraically isomorphic (or conjugate) if $`n=n^{}`$ and if there exists a group automorphism $`\phi :𝕋^n𝕋^n`$ such that $`\phi \alpha =\alpha ^{}\phi `$.
The actions $`\alpha `$ and $`\alpha ^{}`$ are algebraically isomorphic up to a time change if there exists an automorphism $`\phi :𝕋^n𝕋^n`$ and $`CGL(d,)`$ such that $`\phi \alpha C=\alpha ^{}\phi `$.
The action $`\alpha ^{}`$ is an algebraic factor of $`\alpha `$ if there exists a surjective homomorphism $`\phi :𝕋^n𝕋^n^{}`$ such that $`\phi \alpha =\alpha ^{}\phi `$.
The actions $`\alpha `$ and $`\alpha ^{}`$ are weakly algebraically isomorphic if each is an algebraic factor of the other. In this case $`n=n^{}`$ and each factor map has finite fibres.
Finally, we call a map $`\phi :𝕋^n𝕋^n^{}`$ affine if there is a surjective continuous group homomorphism $`\psi :𝕋^n𝕋^n^{}`$ and $`x^{}𝕋^n^{}`$ s.t. $`\phi (x)=\psi (x)+x^{}`$ for every $`x𝕋^n`$.
As already mentioned, we intend to show that under certain condition for $`d2`$, measure theoretic properties imply their algebraic counterparts.
We will say that an algebraic factor $`\alpha ^{}`$ of $`\alpha `$ is a rank–one factor if $`\alpha ^{}`$ is an algebraic factor of $`\alpha `$ and $`\alpha ^{}(_+^d)`$ contains a cyclic sub–semigroup of finite index.
The most general situation when certain rigidity phenomena appear is the following :
$`(^{})`$: The action $`\alpha `$ does not possess nontrivial rank–one algebraic factors.
In the case of actions by automorphisms the condition $`(^{})`$ is equivalent to the following condition $`()`$ (cf. ):
$`()`$: The action $`\alpha `$ contains a group, isomorphic to $`^2`$, which consists of ergodic automorphisms.
By Proposition 6.6 in , Condition $`()`$ is equivalent to saying that the restriction of $`\alpha `$ to a subgroup isomorphic to $`^2`$ is mixing.
A Lyapunov exponent for an action $`\alpha `$ of $`^d`$ is a function $`\chi :^d`$ which associates to each $`𝐧^d`$ the logarithm of the absolute value of the eigenvalue for $`\rho _\alpha ^𝐧`$ corresponding to a fixed eigenvector. Any Lyapunov exponent is a linear function; hence it extends uniquely to $`^d`$. The multiplicity of an exponent is defined as the sum of multiplicities of eigenvalues corresponding to this exponent. Let $`\chi _i,i=1,\mathrm{},k`$, be the different Lyapunov exponents and let $`m_i`$ be the multiplicity of $`\chi _i`$. Then the entropy formula for a single toral endomorphism implies that
$$h_\alpha (𝐧)=h_\mu (\rho _\alpha ^𝐧)=\underset{\{i:\chi _i(𝐧)>0\}}{}m_i\chi _i(𝐧).$$
The function $`h_\alpha :^d`$ is called the entropy function of the action $`\alpha `$. It naturally extends to a symmetric, convex piecewise linear function of $`^d`$. Any cone in $`^d`$ where all Lyapunov exponents have constant sign is called a Weyl chamber. The entropy function is linear in any Weyl chamber.
The entropy function is a prime invariant of measurable isomorphism; since entropy does not increase for factors the entropy function is also invariant of a weak measurable isomorphism. Furthermore it changes equivariantly with respect to automorphisms of $`^d`$.
###### Remark.
it is interesting to point out that the convex piecewise linear structure of the entropy function persists in much greater generality, namely for smooth actions on differentiable manifolds with a Borel invariant measure with compact support.
### 2.3. Finite algebraic factors and invariant lattices
Every algebraic action has many algebraic factors with finite fibres. These factors are in one–to–one correspondence with lattices $`\mathrm{\Gamma }^n`$ which contain the standard lattice $`\mathrm{\Gamma }_0=^n`$, and which satisfy that $`\rho _\alpha (\mathrm{\Gamma })\mathrm{\Gamma }`$. The factor–action associated with a particular lattice $`\mathrm{\Gamma }\mathrm{\Gamma }_0`$ is denoted by $`\alpha _\mathrm{\Gamma }`$. Let us point out that in the case of actions by automorphisms such factors are also invertible: if $`\mathrm{\Gamma }\mathrm{\Gamma }_0`$ and $`\rho _\alpha (\mathrm{\Gamma })\mathrm{\Gamma }`$, then $`\rho _\alpha (\mathrm{\Gamma })=\mathrm{\Gamma }`$.
Let $`\mathrm{\Gamma }\mathrm{\Gamma }_0`$ be a lattice. Take any basis in $`\mathrm{\Gamma }`$ and let $`SGL(n,)`$ be the matrix which maps the standard basis in $`\mathrm{\Gamma }_0`$ to this basis. Then obviously the factor–action $`\alpha _\mathrm{\Gamma }`$ is equal to the action $`\alpha _{S\rho _\alpha S^1}`$. In particular, $`\rho _\alpha `$ and $`\rho _{\alpha _\mathrm{\Gamma }}`$ are conjugate over $``$, although not necessarily over $``$. Notice that conjugacy over $``$ is equivalent to conjugacy over $``$ or over $``$.
For any positive integer $`q`$, the lattice $`\frac{1}{q}\mathrm{\Gamma }_0`$ is invariant under any automorphism in $`GL(n,)`$ and gives rise to a factor which is conjugate to the initial action: one can set $`S=\frac{1}{q}\mathrm{Id}`$ and obtains that $`\rho _\alpha =\rho _{\alpha _{\frac{1}{q}\mathrm{\Gamma }_0}}`$. On the other hand one can find, for any lattice $`\mathrm{\Gamma }\mathrm{\Gamma }_0`$, a positive integer $`q`$ such that $`\frac{1}{q}\mathrm{\Gamma }_0\mathrm{\Gamma }`$ (take $`q`$ the least common multiple of denominators of coordinates for a basis of $`\mathrm{\Gamma }`$). Thus $`\alpha _{\frac{1}{q}\mathrm{\Gamma }_0}`$ appears as a factor of $`\alpha _\mathrm{\Gamma }`$. Summarizing, we have the following properties of finite factors.
###### Proposition 2.1.
Let $`\alpha `$ and $`\alpha ^{}`$ be $`^d`$–actions by automorphism of the torus $`𝕋^n`$. The following are equivalent.
1. $`\rho _\alpha `$ and $`\rho _\alpha ^{}`$ are conjugate over $``$;
2. there exists an action $`\alpha ^{\prime \prime }`$ such that both $`\alpha `$ and $`\alpha ^{}`$ are isomorphic to finite algebraic factors of $`\alpha ^{\prime \prime }`$;
3. $`\alpha `$ and $`\alpha ^{}`$ are weakly algebraically isomorphic, i.e. each of them is isomorphic to a finite algebraic factor of the other.
Obviously, weak algebraic isomorphism implies weak measurable isomorphism. For $``$–actions by Bernoulli automorphisms, weak isomorphism implies isomorphism since it preserves entropy, the only isomorphism invariant for Bernoulli maps. In Section 5 we will show that, for actions by toral automorphisms satisfying Condition $`()`$, measurable isomorphism implies algebraic isomorphism. Hence, existence of such actions which are conjugate over $``$ but not over $``$ provides examples of actions by Bernoulli maps which are weakly isomorphic but not isomorphic.
### 2.4. Dual modules
For any action $`\alpha `$ of $`^d`$ by automorphisms of a compact abelian group $`X`$ we denote by $`\widehat{\alpha }`$ the dual action on the discrete group $`\widehat{X}`$ of characters of $`X`$. For an element $`\chi \widehat{X}`$ we denote $`\widehat{X}_{\alpha _,\chi }`$ the subgroup of $`\widehat{X}`$ generated by the orbit $`\widehat{\alpha }\chi `$.
###### Definition.
The action $`\alpha `$ is called cyclic if $`\widehat{X}_{\alpha _,\chi }=\widehat{X}`$ for some $`\chi \widehat{X}`$.
Cyclicity is obviously an invariant of algebraic conjugacy of actions up to a time change.
More generally, the dual group $`\widehat{X}`$ has the structure of a module over the ring $`[u_1^{\pm 1},\mathrm{},u_d^{\pm 1}]`$ of Laurent polynomials in $`d`$ commuting variables. Action by the generators of $`\widehat{\alpha }`$ corresponds to multiplications by independent variables. This module is called the dual module of the action $`\alpha `$ (cf. ). Cyclicity of the action corresponds to the condition that this module has a single generator. The structure of the dual module up to isomorphism is an invariant of algebraic conjugacy of the action up to a time change.
In the case of the torus $`X=𝕋^n`$ which concerns us in this paper one can slightly modify the construction of the dual module to make it more geometric. A $`^d`$-action $`\alpha `$ by automorphisms of the torus $`^n/^n`$ naturally extends to an action on $`^n`$ (this extension coincides with the embedding $`\rho _\alpha `$ if matrices are identified with linear transformations). This action preserves the lattice $`^n`$ and furnishes $`^n`$ with the structure of a module over the ring $`[u_1^{\pm 1},\mathrm{},u_d^{\pm 1}]`$. This module is — in an obvious sense — a *transpose* of the dual module defined above. In particular, the condition of cyclicity of the action does not depend on which of these two definitions of dual module one adopts.
### 2.5. Algebraic and affine centralizers
Let $`\alpha `$ be an action of $`^d`$ by toral automorphisms, and let $`\rho _\alpha (^d)=\{\rho _\alpha ^𝐧:n^d\}`$. The *centralizer* of $`\alpha `$ in the group of automorphisms of $`𝕋^n`$ is denoted by $`Z(\alpha )`$ and is not distinguished from the centralizer of $`\rho _\alpha (^d)`$ in $`GL(n,)`$.
Similarly, the centralizer of $`\alpha `$ in the semigroup of all endomorphisms of $`𝕋^n`$ (identified with the centralizer of $`\rho _\alpha (^d)`$ in the semigroup $`M(n,)GL(n,)`$) is denoted by $`C(\alpha )`$.
The centralizer of $`\alpha `$ in the group of affine automorphisms of $`𝕋^n`$ will be denoted by $`Z_{\mathrm{𝐴𝑓𝑓}}(\alpha )`$.
The centralizer of $`\alpha `$ in the semigroup of surjective affine maps of $`𝕋^n`$ will be denoted by $`C_{\mathrm{𝐴𝑓𝑓}}(\alpha )`$.
## 3. Irreducible actions
### 3.1. Definition
The action $`\alpha `$ on $`𝕋^n`$ is called irreducible if any nontrivial algebraic factor of $`\alpha `$ has finite fibres.
The following characterization of irreducible actions is useful (cf. ).
###### Proposition 3.1.
The following conditions are equivalent:
1. $`\alpha `$ is irreducible;
2. $`\rho _\alpha `$ contains a matrix with characteristic polynomial irreducible over $``$;
3. $`\rho _\alpha `$ does not have a nontrivial invariant rational subspace or, equivalently, any $`\alpha `$–invariant closed subgroup of $`𝕋^n`$ is finite.
###### Corollary 3.2.
Any irreducible action $`\alpha `$ of $`_+^d,d2`$, satisfies condition ($`^{}`$).
###### Proof.
A rank one algebraic factor has to have fibres of positive dimension. Hence the pre–image of the origin under the factor map is a union of finitely many rational tori of positive dimension and by Proposition 3.1 $`\alpha `$ cannot be irreducible. ∎
### 3.2. Uniqueness of cyclic actions
Cyclicity uniquely determines an irreducible action up to algebraic conjugacy within a class of weakly algebraically conjugate actions.
###### Proposition 3.3.
If $`\alpha `$ is an irreducible cyclic action of $`^d,d1`$, on $`𝕋^n`$ and $`\alpha ^{}`$ is another cyclic action such that $`\rho _\alpha `$ and $`\rho _\alpha ^{}`$ are conjugate over $``$, then $`\alpha `$ and $`\alpha ^{}`$ are algebraically isomorphic.
For the proof of Proposition 3.3 we need an elementary lemma.
###### Lemma 3.4.
Let $`\rho :^dGL(n,)`$ be an irreducible embedding. The centralizer of $`\rho `$ in $`GL(n,)`$ acts transitively on $`^n\{0\}`$.
###### Proof.
By diagonalizing $`\rho `$ over $``$ and taking the real form of it, one immediately sees that the centralizer of $`\rho `$ in $`GL(n,)`$ acts transitively on vectors with nonzero projections on all eigenspaces and thus has a single open and dense orbit. Since the centralizer over $``$ is the closure of the centralizer over $``$, the $``$-linear span of the orbit of any integer or rational vector under the centralizer is an invariant rational subspace. Hence any integer point other than the origin belongs to the single open dense orbit of the centralizer of $`\rho `$ in $`GL(n,)`$. This implies the statement of the lemma. ∎
###### Proof of Proposition 3.3.
Choose $`CM(n,)`$ such that $`C\rho _\alpha ^{}C^1=\rho _\alpha `$. Let $`𝐤,𝐥^n`$ be cyclic vectors for $`\rho _\alpha |_^n`$ and $`\rho _\alpha ^{}|_^n`$, respectively.
Now consider the integer vector $`C(𝐥)`$ and find $`DGL(n,)`$ commuting with $`\rho _\alpha `$ such that $`DC(𝐥)=𝐤`$. We have $`DC\rho _\alpha ^{}C^1D^1=\rho _\alpha `$. The conjugacy $`DC`$ maps bijectively the $``$–span of the $`\rho _\alpha ^{}`$–orbit of $`𝐥`$ to $``$–span of the $`\rho _\alpha `$–orbit of $`𝐤`$. By cyclicity both spans coincide with $`^n`$, and hence $`DCGL(n,)`$. ∎
### 3.3. Centralizers of integer matrices and algebraic number fields
There is an intimate connection between irreducible actions on $`𝕋^n`$ and groups of units in number fields of degree $`n`$. Since this connection (in the particular case where the action is Cartan and hence the number field is totally real) plays a central role in the construction of our principal examples (type (ii) and (iii) of the Introduction), we will describe it here in detail even though most of this material is fairly routine from the point of view of algebraic number theory.
Let $`AGL(n,)`$ be a matrix with an irreducible characteristic polynomial $`f`$ and hence distinct eigenvalues. The centralizer of $`A`$ in $`M(n,)`$ can be identified with the ring of all polynomials in $`A`$ with rational coefficients modulo the principal ideal generated by the polynomial $`f(A)`$, and hence with the field $`K=(\lambda )`$, where $`\lambda `$ is an eigenvalue of $`A`$, by the map
(1)
$$\gamma :p(A)p(\lambda )$$
with $`p[x]`$. Notice that if $`B=p(A)`$ is an integer matrix then $`\gamma (B)`$ is an algebraic integer, and if $`BGL(n,)`$ then $`\gamma (B)`$ is an algebraic unit (converse is not necessarily true).
###### Lemma 3.5.
The map $`\gamma `$ in (1) is injective.
###### Proof.
If $`\gamma (p(A))=1`$ for $`p(A)\mathrm{Id}`$, then $`p(A)`$ has $`1`$ as an eigenvalue, and hence has a rational subspace consisting of all invariant vectors. This subspace must be invariant under $`A`$ which contradicts its irreducibility. ∎
Denote by $`𝒪_K`$ the ring of integers in $`K`$, by $`𝒰_K`$ the group of units in $`𝒪_K`$, by $`C(A)`$ the centralizer of $`A`$ in $`M(n,)`$ and by $`Z(A)`$ the centralizer of $`A`$ in the group $`GL(n,)`$.
###### Lemma 3.6.
$`\gamma (C(A))`$ is a ring in $`K`$ such that $`[\lambda ]\gamma (C(A))𝒪_K`$, and $`\gamma (Z(A))=𝒰_K\gamma (C(A))`$.
###### Proof.
$`\gamma (C(A))`$ is a ring because $`C(A)`$ is a ring. As we pointed out above images of integer matrices are algebraic integers and images of matrices with determinant $`\pm 1`$ are algebraic units. Hence $`\gamma (C(A))𝒪_K`$. Finally, for every polynomial $`p`$ with integer coefficients, $`p(A)`$ is an integer matrix, hence $`[\lambda ]\gamma (C(A))`$. ∎
Notice that $`(\lambda )`$ is a finite index subring of $`𝒪_K`$; hence $`\gamma (C(A))`$ has the same property.
###### Remark.
The groups of units in two different rings, say $`𝒪_1𝒪_2`$, may coincide. Examples can be found in the table of totally real cubic fields in .
###### Proposition 3.7.
$`Z(A)`$ is isomorphic to $`^{r_1+r_21}\times F`$ where $`r_1`$ is the number the real embeddings, $`r_2`$ is the number of pairs of complex conjugate embeddings of the field $`K`$ into $``$, and $`F`$ is a finite cyclic group.
###### Proof.
By lemma 3.6, $`Z(A)`$ is isomorphic to the group of units in the order $`𝒪`$, the statement follows from the Dirichlet Unit Theorem (, Ch.2, §3).∎
Now consider an irreducible action $`\alpha `$ of $`^d`$ on $`𝕋^n`$. Denote $`\rho _\alpha (^d)`$ by $`\mathrm{\Gamma }`$, and let $`\lambda `$ be an eigenvalue of a matrix $`A\mathrm{\Gamma }`$ with an irreducible characteristic polynomial. The centralizers of $`\mathrm{\Gamma }`$ in $`M(n,)`$ and $`GL(n,)`$ coincide with $`C(A)`$ and $`Z(A)`$ correspondingly. The field $`K=(\lambda )`$ has degree $`n`$ and we can consider the map $`\gamma `$ as above. By Lemma 3.6 $`\gamma (\mathrm{\Gamma })𝒰_K`$.
For the purposes of purely algebraic considerations in this and the next section it is convenient to consider actions of integer $`n\times n`$ matrices on $`^n`$ rather than on $`^n`$ and correspondingly to think of $`\alpha `$ as an action by automorphisms of the rational torus $`𝕋_{}^n=^n/^n`$.
Let $`v=(v_1,\mathrm{},v_n)`$ be an eigenvector of $`A`$ with eigenvalue $`\lambda `$ whose coordinates belong to $`K`$. Consider the “projection” $`\pi :^nK`$ defined by $`\pi (r_1,\mathrm{}r_n)=_{i=1}^nr_iv_i`$. It is a bijection (, Prop. 8) which conjugates the action of the group $`\mathrm{\Gamma }`$ with the action on $`K`$ given by multiplication by corresponding eigenvalues $`_{i=1}^d\lambda _i^{k_i},k_1,\mathrm{},k_d`$. Here $`A_1,\mathrm{},A_d\mathrm{\Gamma }`$ are the images of the generators of the action $`\alpha `$, and $`A_iv=\lambda _iv,i=1,\mathrm{},d`$. The lattice $`\pi ^nK`$ is a module over the ring $`[\lambda _1,\mathrm{},\lambda _d]`$.
Conversely, any such data, consisting of an algebraic number field $`K=(\lambda )`$ of degree $`n`$, a $`d`$-tuple $`\overline{\lambda }=(\lambda _1,\mathrm{},\lambda _d)`$ of multiplicatively independent units in $`K`$, and a lattice $`K`$ which is a module over $`[\lambda _1,\mathrm{},\lambda _d]`$, determine an $`^d`$-action $`\alpha _{\overline{\lambda },}`$ by automorphisms of $`𝕋^n`$ up to algebraic conjugacy (corresponding to a choice of a basis in the lattice $``$). This action is generated by multiplications by $`\lambda _1,\mathrm{},\lambda _d`$ (which preserve $``$ by assumption). The action $`\alpha _{\overline{\lambda },}`$ diagonalizes over $``$ as follows. Let $`\varphi _1=\mathrm{id},\varphi _2,\mathrm{},\varphi _n`$ be different embeddings of $`K`$ into $``$. The multiplications by $`\lambda _i,i=1,\mathrm{},d`$, are simultaneously conjugate over $``$ to the respective matrices
$$\left(\begin{array}{cccc}\lambda _i& 0& \mathrm{}& 0\\ 0& \varphi _2\left(\lambda _i\right)& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& \mathrm{}& \varphi _n\left(\lambda _i\right)\end{array}\right),i=1,\mathrm{},d.$$
We will assume that the action is irreducible which in many interesting cases can be easily checked.
Thus, all actions $`\alpha _{\overline{\lambda },}`$ with fixed $`\overline{\lambda }`$ are weakly algebraically isomorphic since the corresponding embeddings are conjugate over $``$ (Proposition 2.1). Actions produced with different sets of units in the same field, say $`\overline{\lambda }`$ and $`\overline{\mu }=(\mu _1,\mathrm{},\mu _d)`$, are weakly algebraically isomorphic if and only if there is an element $`g`$ of the Galois group of $`K`$ such that $`\mu _i=g\lambda _i,i=1,\mathrm{},d`$. By Proposition 3.3 there is a unique cyclic action (up to algebraic isomorphism) within any class of weakly algebraically isomorphic actions: it corresponds to setting $`=[\lambda _1,\mathrm{},\lambda _d]`$; we will denote this action by $`\alpha _{\overline{\lambda }}^{\mathrm{min}}`$. Cyclicity of the action $`\alpha _{\overline{\lambda }}^{\mathrm{min}}`$ is obvious since the whole lattice is obtained from its single element $`1`$ by the action of the ring $`[\lambda _1^{\pm 1},\mathrm{},\lambda _d^{\pm 1}]`$.
Let us summarize this discussion.
###### Proposition 3.8.
Any irreducible action $`\alpha `$ of $`^d`$ by automorphisms of $`𝕋^n`$ is algebraically conjugate to an action of the form $`\alpha _{\overline{\lambda },}`$. It is weakly algebraically conjugate to the cyclic action $`\alpha _{\overline{\lambda }}^{\mathrm{min}}`$. The field $`K=[\lambda _1,\mathrm{},\lambda _d]`$ has degree $`n`$, and the vector of units $`\overline{\lambda }=(\lambda _1,\mathrm{},\lambda _d)`$ is defined up to action by an element of the Galois group of $`K:`$.
Apart from the cyclic model $`\alpha _{\overline{\lambda }}^{\mathrm{min}}`$ there is another canonical choice of the lattice $``$, namely the ring of integers $`𝒪_K`$. We will denote the action $`\alpha _{\overline{\lambda },𝒪_K}`$ by $`\alpha _{\overline{\lambda }}^{\mathrm{max}}`$. More generally, one can choose as the lattice $``$ any subring $`𝒪`$ such that $`[\lambda _1,\mathrm{},\lambda _d]𝒪𝒪_K`$.
###### Proposition 3.9.
Assume that $`𝒪[\lambda _1,\mathrm{},\lambda _d]`$. Then the action $`\alpha _{\overline{\lambda },𝒪}`$ is not algebraically isomorphic up to a time change to $`\alpha _{\overline{\lambda }}^{\mathrm{min}}`$. In particular, if $`𝒪_K[\lambda _1,\mathrm{},\lambda _d]`$, then the actions $`\alpha _{\overline{\lambda }}^{\mathrm{max}}`$ and $`\alpha _{\overline{\lambda }}^{\mathrm{min}}`$ are not algebraically isomorphic up to a time change.
###### Proof.
Let us denote the centralizers in $`M(n,)`$ of the actions $`\alpha _{\overline{\lambda },𝒪}`$ and $`\alpha _{\overline{\lambda }}^{\mathrm{min}}`$ by $`C_1`$ and $`C_2`$, respectively. The centralizer $`C_1`$ contains multiplications by all elements of $`𝒪`$. For, if one takes any basis in $`𝒪`$, the multiplication by an element $`\mu 𝒪`$ takes elements of the basis into elements of $`𝒪`$, which are linear combinations with integral coefficients of the basis elements; hence the multiplication is given by an integer matrix. On the other hand any element of each centralizer is a multiplication by an integer in $`K`$ (Lemma 3.6).
Now assume that the multiplication by $`\mu 𝒪_K`$ belongs to $`C_2`$. This means that this multiplication preserves $`[\lambda _1,\mathrm{},\lambda _d]`$; in particular, $`\mu =\mu 1[\lambda _1,\mathrm{},\lambda _d]`$. Thus $`C_2`$ consists of multiplication by elements of $`[\lambda _1,\mathrm{},\lambda _d]`$. An algebraic isomorphism up to a time change has to preserve both the module of polynomials with integer coefficients in the generators of the action and the centralizer of the action in $`M(n,)`$, which is impossible. ∎
The central question which appears in connection with our examples is the classification of weakly algebraically isomorphic Cartan actions up to algebraic isomorphism.
Proposition 3.9 is useful in distinguishing weakly algebraically isomorphic actions when $`𝒪_K[\lambda _1,\mathrm{},\lambda _d]`$. Cyclicity also can serve as a distinguishing invariant.
###### Corollary 3.10.
The action $`\alpha _{\overline{\lambda },𝒪}`$ is cyclic if and only if $`𝒪=[\lambda _1,\mathrm{},\lambda _d]`$.
###### Proof.
The action $`\alpha _{\overline{\lambda }}^{\mathrm{min}}`$ corresponding to the ring $`[\lambda _1,\mathrm{},\lambda _d]`$ is cyclic by definition since the ring coincides with the orbit of $`1`$. By Proposition 3.3, if $`\alpha _{\overline{\lambda },𝒪}`$ were cyclic, it would be algebraically conjugate to $`\alpha _{\overline{\lambda }}^{\mathrm{min}}`$, which, by Proposition 3.9, implies that $`𝒪=[\lambda _1,\mathrm{},\lambda _d]`$. ∎
The property common to all actions of the $`\alpha _{\overline{\lambda },𝒪}`$ is transitivity of the action of the centralizer $`C(\alpha _{\overline{\lambda },𝒪})`$ on the lattice. Similarly to cyclicity this property is obviously an invariant of algebraic conjugacy up to a time change.
###### Proposition 3.11.
Any irreducible action $`\alpha `$ of $`^d`$ by automorphisms of $`𝕋^n`$ whose centralizer $`C(\alpha )`$ in $`M(n,)`$ acts transitively on $`^n`$ is algebraically isomorphic to an action $`\alpha _{\overline{\lambda },𝒪}`$, where $`𝒪𝒪_K`$ is a ring which contains $`[\lambda _1,\mathrm{},\lambda _d]`$.
###### Proof.
By Proposition 3.8 any irreducible action $`\alpha `$ of $`^d`$ by automorphisms of $`𝕋^n`$ is algebraically conjugate to an action of the form $`\alpha _{\overline{\lambda },}`$ for a lattice $`K`$. Let $`C`$ be the centralizer of $`\alpha _{\overline{\lambda },}`$ in the semigroup of linear endomorphisms of $``$. We fix an element $`\beta `$ with $`C(\alpha )\beta =`$ and consider conjugation of the action $`\alpha _{\overline{\lambda },}`$ by multiplication by $`\beta ^1`$; this is simply $`\alpha _{\overline{\lambda },\beta ^1}`$. The centralizer of $`\alpha _{\overline{\lambda },\beta ^1}`$ acts on the element $`1\beta ^1`$ transitively. By Lemma 3.6 the centralizer consists of all multiplications by elements of a certain subring $`𝒪𝒪_K`$ which contains $`[\lambda _1,\mathrm{},\lambda _d]`$. Thus $`1\beta ^1=𝒪`$. ∎
### 3.4. Structure of algebraic and affine centralizers for irreducible actions
By Lemma 3.6, the centralizer $`C(\alpha )`$, as an additive group, is isomorphic to $`^n`$ and has an additional ring structure. In the terminology of Proposition 3.7, the centralizer $`Z(\alpha )`$ for an irreducible action $`\alpha `$ by toral automorphisms is isomorphic to $`^{r_1+r_21}\times F`$.
An irreducible action $`\alpha `$ has maximal rank if $`d=r_1+r_21`$. In this case $`Z(\alpha )`$ is a finite extension of $`\alpha `$.
Notice that any affine map commuting with an action $`\alpha `$ by toral automorphisms preserves the set $`\mathrm{Fix}(\alpha )`$ of fixed points of the action. This set is always a subgroup of the torus and hence, for an irreducible action, always finite. The translation by any element of $`\mathrm{Fix}(\alpha )`$ commutes with $`\alpha `$ and thus belongs to $`Z_{\mathrm{𝐴𝑓𝑓}}(\alpha )`$. Furthermore, the affine centralizers $`Z_{\mathrm{𝐴𝑓𝑓}}(\alpha )`$ and $`C_{\mathrm{𝐴𝑓𝑓}}(\alpha )`$ are generated by these translations and, respectively, $`Z(\alpha )`$ and $`C(\alpha )`$.
###### Remark.
Most of the material of this section extends to general irreducible actions of $`^d`$ by automorphisms of compact connected abelian groups; a group possessing such an action must be a torus or a solenoid (). In the solenoid case, which includes natural extensions of $`^d`$–actions by toral endomorphisms, the algebraic numbers $`\lambda _1,\mathrm{},\lambda _d`$ which appear in the constructions are not in general integers. As we mentioned in the introduction we restrict our algebraic setting here since we are able to exhibit some of the most interesting and striking new phenomena using Cartan actions and certain actions directly derived from them. However, other interesting examples appear for actions on the torus connected with not totally real algebraic number fields, actions on solenoids, and actions on zero-dimensional abelian groups (cf. e.g. ).
One can also extend the setup of this section to certain classes of reducible actions. Since some of these satisfy condition $`()`$ basic rigidity results still hold and a number of further interesting examples can be constructed.
## 4. Cartan actions
### 4.1. Structure of Cartan actions
Of particular interest for our study are abelian groups of ergodic automorphisms of $`𝕋^n`$ of maximal possible rank $`n1`$ (in agreement with the real rank of the Lie group $`SL(n,)`$).
###### Definition.
An action of $`^{n1}`$ on $`𝕋^n`$ for $`n3`$ by ergodic automorphisms is called a Cartan action.
###### Proposition 4.1.
Let $`\alpha `$ be a Cartan action on $`𝕋^n`$.
1. Any element of $`\rho _\alpha `$ other than identity has real eigenvalues and is hyperbolic and thus Bernoulli.
2. $`\alpha `$ is irreducible.
3. The centralizer of $`Z(\alpha )`$ is a finite extension of $`\rho _\alpha (^{n1})`$.
###### Proof.
First, let us point out that it is sufficient to prove the proposition for irreducible actions. For, if $`\alpha `$ is not irreducible, it has a nontrivial irreducible algebraic factor of dimension, say, $`mn1`$. Since every factor of an ergodic automorphism is ergodic, we thus obtain an action of $`^{n1}`$ in $`𝕋^m`$ by ergodic automorphisms. By considering a restriction of this action to a subgroup of rank $`m1`$ which contains an irreducible matrix, we obtain a Cartan action on $`𝕋^m`$. By Statement 3. for irreducible actions, the centralizer of this Cartan action is a finite extension of $`^{m1}`$, and thus cannot contain $`^{n1}`$, a contradiction.
Now assuming that $`\alpha `$ is irreducible, take a matrix $`A\rho _\alpha (^{n1})`$ with irreducible characteristic polynomial $`f`$. Such a matrix exists by Proposition 3.1. It has distinct eigenvalues, say $`\lambda =\lambda _1,\mathrm{},\lambda _n`$. Consider the correspondence $`\gamma `$ defined in (1). By Lemma 3.6 for every $`B\rho _\alpha (^{n1})`$ we have $`\gamma (B)𝒰_K`$, hence the group of units $`𝒰_K`$ in $`K`$ contains a subgroup isomorphic to $`^{n1}`$. By the Dirichlet Unit Theorem the rank of the group of units in $`K`$ is equal to $`r_1+r_21`$, where $`r_1`$ is the number of real embeddings and $`r_2`$ is the number of pairs of complex conjugate embeddings of $`K`$ into $``$. Since $`r_1+2r_2=n`$ we deduce that $`r_2=0`$, so the field $`K`$ is totally real, that is all eigenvalues of $`A`$, and hence of any matrix in $`\rho _\alpha (^{n1})`$, are real. The same argument gives Statement 3, since any element of the centralizer of $`\rho _\alpha (^{n1})`$ in $`GL(n,)`$ corresponds to a unit in $`K`$. Hyperbolicity of matrices in $`\rho _\alpha (^{n1})`$ is proved in the same way as Lemma 3.5. ∎
###### Lemma 4.2.
Let $`A`$ be a hyperbolic matrix in $`SL(n,)`$ with irreducible characteristic polynomial and distinct real eigenvalues. Then every element of the centralizer $`Z(A)`$ other than $`\{\pm 1\}`$ is hyperbolic.
###### Proof.
Assume that $`BZ(A)`$ is not hyperbolic. As $`B`$ is simultaneously diagonalizable with $`A`$ and has real eigenvalues, it has an eigenvalue $`+1`$ or $`1`$. The corresponding eigenspace is rational and $`A`$–invariant. Since $`A`$ is irreducible, this eigenspace has to coincide with the whole space and hence $`B=\pm 1`$. ∎
###### Corollary 4.3.
Cartan actions are exactly the maximal rank irreducible actions corresponding to totally real number fields.
###### Corollary 4.4.
The centralizer $`Z(\alpha )`$ for a Cartan action $`\alpha `$ is isomorphic to $`^{n1}\times \{\pm 1\}`$.
We will call a Cartan action $`\alpha `$ maximal if $`\alpha `$ is an index two subgroup in $`Z(\alpha )`$.
Let us point out that $`Z_{\mathrm{𝐴𝑓𝑓}}(\alpha )`$ is isomorphic $`Z(\alpha )\times \mathrm{Fix}(\alpha )`$. Thus, the factor of $`Z_{\mathrm{𝐴𝑓𝑓}}(\alpha )`$ by the subgroup of finite order elements is always isomorphic to $`^{n1}`$. If $`\alpha `$ is maximal, this factor is identified with $`\alpha `$ itself. In the next Section we will show (Corollary 5.4) that for a Cartan action $`\alpha `$ on $`𝕋^n,n3`$ the isomorphism type of the pair $`(Z_{\mathrm{𝐴𝑓𝑓}}(\alpha ),\alpha )`$ is an invariant of the measurable isomorphism. Thus, in particular, for a maximal Cartan action the order of the group $`\mathrm{Fix}(\alpha )`$ is a measurable invariant.
###### Remark.
An important geometric distinction between Cartan actions and general irreducible actions by hyperbolic automorphisms is the absence of multiple Lyapunov exponents. This greatly simplifies proofs of various rigidity properties both in the differentiable and measurable context.
### 4.2. Algebraically nonisomorphic maximal Cartan actions
In Section 3.3 we described a particular class of irreducible actions $`\alpha _{\overline{\lambda },𝒪}`$ which is characterized by the transitivity of the action of the centralizer $`C(\alpha _{\overline{\lambda },𝒪})`$ on the lattice (Proposition 3.11). In the case $`𝒪_K=[\lambda ]`$ there is only one such action, namely the cyclic one (Corollary 3.10). Now we will analyze this special case for totally real fields in detail and show how information about the class number of the field helps to construct algebraically nonisomorphic maximal Cartan actions. This will in particular provide examples of Cartan actions not isomorphic up to a time change to any action of the form $`\alpha _{\overline{\lambda },𝒪}`$.
It is well–known that for $`n=2`$ there are natural bijections between conjugacy classes of hyperbolic elements in $`SL(2,)`$ of a given trace, ideal classes in the corresponding real quadratic field, and congruence classes of primitive integral indefinite quadratic forms of the corresponding discriminant. This has been used by Sarnak in his proof of the Prime Geodesic Theorem for surfaces of constant negative curvature (see also ). It follows from an old Theorem of Latimer and MacDuffee (see , , and a more modern account in ), that the first bijection persists for $`n>2`$. Let $`A`$ a hyperbolic matrix $`ASL(n,)`$ with irreducible characteristic polynomial $`f`$, and hence distinct real eigenvalues, $`K=(\lambda )`$, where $`\lambda `$ is an eigenvalue of $`A`$, and $`𝒪_K=[\lambda ]`$. To each matrix $`A^{}`$ with the same eigenvalues, we assign the eigenvector $`v=(v_1,\mathrm{},v_n)`$ with eigenvalue $`\lambda `$: $`A^{}v=\lambda v`$ with all its entries in $`𝒪_K`$, which can be always done, and to this eigenvector, an ideal in $`𝒪_K`$ with the $``$–basis $`v_1,\mathrm{},v_n`$. The described map is a bijection between the $`GL(n,)`$–conjugacy classes of matrices in $`SL(n,)`$ which have the same characteristic polynomial $`f`$ and the set of ideal classes in $`𝒪_K`$. Moreover, it allows us to reach conclusions about centralizers as well.
###### Theorem 4.5.
Let $`ASL(n,)`$ be a hyperbolic matrix with irreducible characteristic polynomial $`f`$ and distinct real eigenvalues, $`K=(\lambda )`$ where $`\lambda `$ is an eigenvalue of $`A`$, and $`𝒪_K=[\lambda ]`$. Suppose the number of eigenvalues among $`\lambda _1,\mathrm{},\lambda _n`$ that belong to $`K`$ is equal to $`r`$. If the class number $`h(K)>r`$, then there exists a matrix $`A^{}SL(n,)`$ having the same eigenvalues as $`A`$ whose centralizer $`Z(A^{})`$ is not conjugate in $`GL(n,)`$ to $`Z(A)`$. Furthermore, the number of matrices in $`SL(n,)`$ having the same eigenvalues as $`A`$ with pairwise nonconjugate (in $`GL(n,)`$) centralizers is at least $`[\frac{h(K)}{r}]+1`$, where $`[x]`$ is the largest integer $`<x`$.
###### Proof.
Suppose the matrix $`A`$ corresponds to the ideal class $`I_1`$ with the $``$–basis $`v^{(1)}`$. Then
$$Av^{(1)}=\lambda v^{(1)}.$$
Since $`h(K)>1`$, there exists a matrix $`A_2`$ having the same eigenvalues which corresponds to a different ideal class $`I_2`$ with the basis $`v^{(2)}`$, and we have
$$A_2v^{(2)}=\lambda v^{(2)}.$$
The eigenvectors $`v^{(1)}`$ and $`v^{(2)}`$ are chosen with all their entries in $`𝒪_K`$. Now assume that $`Z(A_2)`$ is conjugate to $`Z(A)`$. Then $`Z(A_2)`$ contains a matrix $`B_2`$ conjugate to $`A`$. Since $`B_2`$ commutes with $`A_2`$ we have $`B_2v^{(2)}=\mu _2v^{(2)}`$, and since $`B_2`$ is conjugate to $`A`$, $`\mu _2`$ is one of the roots of $`f`$. Moreover, since $`B_2SL(n,)`$ and all entries of $`v^{(2)}`$ are in $`K`$, $`\mu _2K`$. Thus $`\mu _2`$ is one of $`r`$ roots of $`f`$ which belongs to $`K`$.
From $`B_2=S^1AS`$ ($`SGL(n,)`$) we deduce that $`\mu _2(Sv^{(2)})=A(Sv^{(2)})`$. Since $`I_1`$ and $`I_2`$ belong to different ideal classes, $`Sv^{(2)}kv^{(1)}`$ for any $`k`$ in the quotient field of $`𝒪_K`$, and since $`\lambda `$ is a simple eigenvalue for $`A`$, we deduce that $`\mu _2\lambda `$, and thus $`\mu _2`$ can take one of the $`r1`$ remaining values.
Now assume that $`A_3`$ corresponds to the third ideal class, i.e
$$A_3v^{(3)}=\lambda v^{(3)},$$
and $`B_3`$ commutes with $`A_3`$ and is conjugate to $`A`$, and hence to $`B_2`$. Then $`B_3v^{(3)}=\mu _3v^{(3)}`$ where $`\mu _3`$ is a root of $`f`$ belonging to the field $`K`$. By the previous considerations, $`\mu _3\lambda `$ and $`\mu _3\mu _2`$. An induction argument shows that if the class number of $`K`$ is greater than $`r`$, there exists a matrix $`A^{}`$ such that no matrix in $`Z(A^{})`$ is conjugate to $`A`$, i.e. $`Z(A^{})`$ and $`Z(A)`$ are not conjugate in $`GL(n,)`$.
Since $`A^{}`$ has the same characteristic polynomial as $`A`$, continuing the same process, we can find not more than $`r`$ matrices representing different ideal classes having centralizers conjugate to $`Z(A^{})`$, and the required estimate follows. ∎
## 5. Measure–theoretic rigidity of conjugacies, centralizers, and factors
### 5.1. Conjugacies
Suppose $`\alpha `$ and $`\alpha ^{}`$ are measurable actions of the same group $`G`$ by measure–preserving transformations of the spaces $`(X,\mu )`$ and $`(Y,\nu )`$, respectively. If $`H:(X,\mu )(Y,\nu )`$ is a metric isomorphism (conjugacy) between the actions then the lift of the measure $`\mu `$ onto the $`\text{graph}HX\times Y`$ coincides with the lift of $`\nu `$ to $`\text{graph}H^1`$. The resulting measure $`\eta `$ is a very special case of a joining of $`\alpha `$ and $`\alpha ^{}`$: it is invariant under the diagonal (product) action $`\alpha \times \alpha ^{}`$ and its projections to $`X`$ and $`Y`$ coincide with $`\mu `$ and $`\nu `$, respectively. Obviously the projections establish metric isomorphism of the action $`\alpha \times \alpha ^{}`$ on $`(X\times Y,\eta )`$ with $`\alpha `$ on $`(X,\mu )`$ and $`\alpha ^{}`$ on $`(Y,\nu )`$ correspondingly.
Similarly, if an automorphism $`H:(X,\mu )(X,\mu )`$ commutes with the action $`\alpha `$, the lift of $`\mu `$ to $`\text{graph}HX\times X`$ is a self-joining of $`\alpha `$, i.e. it is $`\alpha \times \alpha `$–invariant and both of its projections coincide with $`\mu `$. Thus an information about invariant measures of the products of different actions as well as the product of an action with itself may give an information about isomorphisms and centralizers.
The use of this joining construction in order to deduce rigidity of isomorphisms and centralizers from properties of invariant measures of the product was first suggested in this context to the authors by J.-P. Thouvenot.
In both cases the ergodic properties of the joining would be known because of the isomorphism with the original actions. Very similar considerations apply to the actions of semi–groups by noninvertible measure–preserving transformations. We will use the following corollary of the results of .
###### Theorem 5.1.
Let $`\alpha `$ be an action of $`^2`$ by ergodic toral automorphisms and let $`\mu `$ be a weakly mixing $`\alpha `$–invariant measure such that for some $`𝐦^2`$, $`\alpha ^𝐦`$ is a $`K`$-automorphism. Then $`\mu `$ is a translate of Haar measure on an $`\alpha `$–invariant rational subtorus.
###### Proof.
We refer to Corollary 5.2’ from (, “Corrections…”). According to this corollary the measure $`\mu `$ is an extension of a zero entropy measure for an algebraic factor of smaller dimension with Haar conditional measures in the fiber. But since $`\alpha `$ contains a $`K`$-automorphism it does not have non–trivial zero entropy factors. Hence the factor in question is the action on a single point and $`\mu `$ itself is a Haar measure on a rational subtorus. ∎
Conclusion of Theorem 5.1 obviously holds for any action of $`^d,d2`$ which contains a subgroup $`^2`$ satisfying assumptions of Theorem 5.1. Thus we can deduce the following result which is central for our constructions.
###### Theorem 5.2.
Let $`\alpha `$ and $`\alpha ^{}`$ be two actions of $`^d`$ by automorphisms of $`𝕋^n`$ and $`𝕋^n^{}`$ correspondingly and assume that $`\alpha `$ satisfies condition $`()`$. Suppose that $`H:𝕋^n𝕋^n^{}`$ is a measure–preserving isomorphism between $`(\alpha ,\lambda )`$ and $`(\alpha ^{},\lambda )`$, where $`\lambda `$ is Haar measure. Then $`n=n^{}`$ and $`H`$ coincides (mod 0) with an affine automorphism on the torus $`𝕋^n`$, and hence $`\alpha `$ and $`\alpha ^{}`$ are algebraically isomorphic.
###### Proof.
First of all, condition $`()`$ is invariant under metric isomorphism, hence $`\alpha ^{}`$ also satisfies this condition. But ergodicity with respect to Haar measure can also be expressed in terms of the eigenvalues; hence $`\alpha \times \alpha ^{}`$ also satisfies ($``$). Now consider the joining measure $`\eta `$ on $`\text{graph}H𝕋^{n+n^{}}`$. The conditions of Theorem 5.1 are satisfied for the invariant measure $`\eta `$ of the action $`\alpha \times \alpha ^{}`$. Thus $`\eta `$ is a translate of Haar measure on a rational $`\alpha \times \alpha ^{}`$–invariant subtorus $`𝕋^{}𝕋^{n+n^{}}=𝕋^n\times 𝕋^n^{}`$. On the other hand we know that projections of $`𝕋^{}`$ to both $`𝕋^n`$ and $`𝕋^n^{}`$ preserve Haar measure and are one–to–one. The partitions of $`𝕋^{}`$ into pre–images of points for each of the projections are measurable partitions and Haar measures on elements are conditional measures. This implies that both projections are onto, both partitions are partitions into points, and hence $`n=n^{}`$ and $`𝕋^{}=\text{graph}I`$, where $`I:𝕋^n𝕋^n`$ is an affine automorphism which has to coincide $`(\text{mod}\mathrm{\hspace{0.33em}0})`$ with the measure–preserving isomorphism $`H`$. ∎
Since a time change is in a sense a trivial modification of an action we are primarily interested in distinguishing actions up to a time change. The corresponding rigidity criterion follows immediately from Theorem 5.2.
###### Corollary 5.3.
Let $`\alpha `$ and $`\alpha ^{}`$ be two actions of $`^d`$ by automorphisms of $`𝕋^n`$ and $`𝕋^n^{}`$, respectively, and assume that $`\alpha `$ satisfies condition $`()`$. If $`\alpha `$ and $`\alpha ^{}`$ are measurably isomorphic up to a time change then they are algebraically isomorphic up to a time change.
### 5.2. Centralizers
Applying Theorem 5.2 to the case $`\alpha =\alpha ^{}`$ we immediately obtain rigidity of the centralizers.
###### Corollary 5.4.
Let $`\alpha `$ be an action of $`^d`$ by automorphisms of $`𝕋^n`$ satisfying condition $`()`$. Any invertible Lebesgue measure–preserving transformation commuting with $`\alpha `$ coincides (mod 0) with an affine automorphism of $`𝕋^n`$.
Any affine transformation commuting with $`\alpha `$ preserves the finite set of fixed points of the action. Hence the centralizer of $`\alpha `$ in affine automorphisms has a finite index subgroups which consist of automorphisms and which corresponds to the centralizer of $`\rho _\alpha (^d)`$ in $`GL(n,)`$.
Thus, in contrast with the case of a single automorphism, the centralizer of such an action $`\alpha `$ is not more than countable, and can be identified with a finite extension of a certain subgroup of $`GL(n,)`$. As an immediate consequence we obtain the following result.
###### Proposition 5.5.
For any $`d`$ and $`k`$, $`2dk`$, there exists a $`^d`$–action by hyperbolic toral automorphisms such that its centralizer in the group of Lebesgue measure–preserving transformations is isomorphic to $`\{\pm 1\}\times ^k`$.
###### Proof.
Consider a hyperbolic matrix $`ASL(k+1,)`$ with irreducible characteristic polynomial and real eigenvalues such that the origin is the only fixed point of $`F_A`$. Consider a subgroup of $`Z(A)`$ isomorphic to $`^d`$ and containing $`A`$ as one of its generators. This subgroup determines an embedding $`\rho :^dSL(k+1,)`$. Since $`d2`$ and by Proposition 4.2, all matrices in $`\rho (^d)`$ are hyperbolic and hence ergodic, condition $`()`$ is satisfied. Hence by Corollary 5.4, the measure–theoretic centralizer of the action $`\alpha _\rho `$ coincides with its algebraic centralizer, which, in turn, and obviously, coincides with centralizer of the single automorphism $`F_A`$ isomorphic to $`\{\pm 1\}\times ^k`$. ∎
### 5.3. Factors, noninvertible centralizers and weak isomorphism
A small modification of the proof of Theorem 5.2 produces a result about rigidity of factors.
###### Theorem 5.6.
Let $`\alpha `$ and $`\alpha ^{}`$ be two actions of $`^d`$ by automorphisms of $`𝕋^n`$ and $`𝕋^n^{}`$ respectively, and assume that $`\alpha `$ satisfies condition $`()`$. Suppose that $`H:𝕋^n𝕋^n^{}`$ is a Lebesgue measure–preserving transformation such that $`H\alpha =\alpha ^{}H`$. Then $`\alpha ^{}`$ also satisfies $`()`$ and $`H`$ coincides (mod 0) with an epimorphism $`h:𝕋^n𝕋^n^{}`$ followed by translation. In particular, $`\alpha ^{}`$ is an algebraic factor of $`\alpha `$.
###### Proof.
Since $`\alpha ^{}`$ is a measurable factor of $`\alpha `$, every element which is ergodic for $`\alpha `$ is also ergodic for $`\alpha ^{}`$. Hence $`\alpha ^{}`$ also satisfies condition $`()`$. As before consider the product action $`\alpha \times \alpha ^{}`$ which now by the same argument also satisfies $`()`$. Take the $`\alpha \times \alpha ^{}`$ invariant measure $`\eta =(\text{Id}\times H)_{}\lambda `$ on $`\text{graph}H`$. This measure provides a joining of $`\alpha `$ and $`\alpha ^{}`$. Since $`(\alpha \times \alpha ^{},(\text{Id}\times H)_{}\lambda )`$ is isomorphic to $`(\alpha ,\lambda )`$ the conditions of Corollary 5.1 are satisfied and $`\eta `$ is a translate of Haar measure on an invariant rational subtorus $`𝕋^{}`$. Since $`𝕋^{}`$ projects to the first coordinate one-to-one we deduce that $`H`$ is an algebraic epimorphism (mod 0) followed by a translation. ∎
Similarly to the previous section the application of Theorem 5.6 to the case $`\alpha =\alpha ^{}`$ gives a description of the centralizer of $`\alpha `$ in the group of all measure–preserving transformations.
###### Corollary 5.7.
Let $`\alpha `$ be an action of $`^d`$ by automorphisms of $`𝕋^n`$ satisfying condition $`()`$. Any Lebesgue measure–preserving transformation commuting with $`\alpha `$ coincides (mod 0) with an affine map on $`𝕋^n`$.
Now we can obtain the following strengthening of Proposition 2.1 for actions satisfying condition $`()`$ which is one of the central conclusions of this paper.
###### Theorem 5.8.
Let $`\alpha `$ be an action of $`^d`$ by automorphisms of $`𝕋^n`$ satisfying condition $`()`$ and $`\alpha ^{}`$ another $`^d`$-action by toral automorphisms. Then $`(\alpha ,\lambda )`$ is weakly isomorphic to $`(\alpha ^{},\lambda ^{})`$ if and only if $`\rho _\alpha `$ and $`\rho _\alpha ^{}`$ are isomorphic over $``$, i.e. if $`\alpha `$ and $`\alpha ^{}`$ are finite algebraic factors of each other.
###### Proof.
By Theorem 5.6, $`\alpha `$ and $`\alpha ^{}`$ are algebraic factors of each other. This implies that $`\alpha ^{}`$ acts on the torus of the same dimension $`n`$ and hence both algebraic factor–maps have finite fibres. Now the statement follows from Proposition 2.1. ∎
### 5.4. Distinguishing weakly isomorphic actions
Similarly we can translate criteria for algebraic conjugacy of weakly algebraically conjugate actions to the measurable setting.
###### Theorem 5.9.
If $`\alpha `$ is an irreducible cyclic action of $`^d,d2`$, on $`𝕋^n`$ and $`\alpha ^{}`$ is a non–cyclic $`^d`$-action by toral automorphisms. Then $`\alpha `$ and $`\alpha ^{}`$ are not measurably isomorphic up to a time change.
###### Proof.
Since action $`\alpha `$ satisfies condition $`()`$ (Corollary 3.2) we can apply Theorem 5.8 and conclude that we only need to consider the case when $`\rho _\alpha `$ and $`\rho _\alpha ^{}`$ are isomorphic over $``$ up to a time change. But then, by Proposition 3.3, $`\alpha `$ and $`\alpha ^{}`$ are not algebraically isomorphic up to a time change and hence, by Corollary 5.3, they are not measurably isomorphic up to a time change. ∎
Combining Proposition 3.9 and Corollary 5.3 we immediately obtain rigidity for the minimal irreducible models.
###### Corollary 5.10.
Assume that $`𝒪[\lambda _1,\mathrm{},\lambda _d]`$. Then the action $`\alpha _{\overline{\lambda },𝒪}`$ is not measurably isomorphic up to a time change to $`\alpha _{\overline{\lambda }}^{\mathrm{min}}`$. In particular, if $`𝒪_K[\lambda _1,\mathrm{},\lambda _d]`$, then the actions $`\alpha _{\overline{\lambda }}^{\mathrm{max}}`$ and $`\alpha _{\overline{\lambda }}^{\mathrm{min}}`$ are not measurably isomorphic up to a time change.
## 6. Examples
Now we proceed to produce examples of actions for which the entropy data coincide but which are not algebraically isomorphic, and hence by Theorem 5.2 not measure–theoretically isomorphic.
### 6.1. Weakly nonisomorphic actions
In this section we consider actions which are not algebraically isomorphic over $``$ (or, equivalently, over $``$) and hence by Theorem 5.8 are not even weakly isomorphic. The easiest way is as follows.
Example 1a. Start with any action $`\alpha `$ of $`^d,d2`$, by ergodic automorphisms of $`𝕋^n`$. We may double the entropies of all its elements in two different ways: by considering the Cartesian square $`\alpha \times \alpha `$ acting on $`𝕋^{2n}`$, and by taking second powers of all elements: $`\alpha _2^𝐧=\alpha ^{2𝐧}`$ for all $`𝐧^d`$. Obviously $`\alpha \times \alpha `$ is not algebraically isomorphic to $`\alpha _2`$, since, for example, they act on tori of different dimension. Hence by Theorem 5.2 $`(\alpha \times \alpha ,\lambda )`$ is not metrically isomorphic to $`(\alpha _2,\lambda )`$ either.
Now we assume that $`\alpha `$ contains an automorphism $`F_A`$ where $`A`$ is hyperbolic with an irreducible characteristic polynomial and distinct positive real eigenvalues. In this case it is easy to find an invariant distinguishing the two actions, namely, the algebraic type of the centralizer of the action in the group of measure–preserving transformations. By Corollary 5.4, the centralizer of $`\alpha `$ in the group of measure–preserving transformations coincides with the centralizer in the group of affine maps, which is a finite extension of the centralizer in the group of automorphisms. By the Dirichlet Unit Theorem, the centralizer of $`Z(\alpha _2)`$ in the group of automorphisms of the torus is isomorphic to $`\{\pm 1\}\times ^{n1}`$, whereas the centralizer of $`\alpha \times \alpha `$ contains the $`^{2(n1)}`$–action by product transformations $`\alpha ^{𝐧_1}\times \alpha ^{𝐧_2},𝐧_1,𝐧_2^{n1}`$. In fact, the centralizer of $`\alpha \times \alpha `$ can be calculated explicitly:
###### Proposition 6.1.
Let $`\lambda `$ be an eigenvalue of $`A`$. Then $`K=(\lambda )`$ is a totally real algebraic field. If its ring of integers $`𝒪_K`$ is equal to $`[\lambda ]`$ then the centralizer of $`\alpha \times \alpha `$ in $`GL(2n,)`$ is isomorphic to the group $`GL(2,𝒪_K)`$, i.e. the group of $`2\times 2`$ matrices with entries in $`𝒪_K`$ whose determinant is a unit in $`𝒪_K`$.
###### Proof.
First we notice that a matrix in block form $`B=\left(\begin{array}{cc}X& Y\\ Z& T\end{array}\right)`$ with $`X,Y,Z,TM(n,)`$ commutes with $`\left(\begin{array}{cc}A& 0\\ 0& A\end{array}\right)`$ if an only if $`X,Y,Z,T`$ commute with $`A`$ and can thus be identified with elements of $`𝒪_K`$. In this case $`B`$ can be identified with a matrix in $`M(2,𝒪_K)`$. Since $`det\left(\begin{array}{cc}X& Y\\ Z& T\end{array}\right)=det(XTYZ)=\pm 1`$ (cf. ), the norm of the determinant of the $`2\times 2`$ matrix corresponding to $`B`$ is equal $`\pm 1`$. Hence this determinant is a unit in $`𝒪_K`$, and we obtain the desired isomorphism. ∎
It is not difficult to modify Example 1a to obtain weakly nonisomorphic actions with the same entropy on the torus of the same dimension.
Example 1b. For a natural number $`k`$ define the action $`\alpha _k`$ similarly to $`\alpha _2`$: $`\alpha _k^𝐧=\alpha ^{k𝐧}`$ for all $`𝐧^d`$.
The actions $`\alpha _3\times \alpha `$ and $`\alpha _2\times \alpha _2`$ act on $`𝕋^{2n}`$, have the same entropies for all elements and are not isomorphic.
As before, we can see that centralizers of these two actions are not isomorphic. In particular, the centralizer of $`\alpha _3\times \alpha `$ is abelian since it has simple eigenvalues, while the centralizer of $`\alpha _2\times \alpha _2`$ is not.
### 6.2. Cartan actions distinguished by cyclicity or maximality
We give two examples which illustrate the method of Section 3.3. They provide weakly algebraically isomorphic Cartan actions of $`^2`$ on $`𝕋^3`$ which are not algebraically isomorphic even up to a time change (i.e. a linear change of coordinates in $`^2`$) by Proposition 3.9. These examples utilize the existence of number fields $`K=(\lambda )`$ and units $`\overline{\lambda }=(\lambda _1,\lambda _2)`$ in them for which $`𝒪_K[\lambda _1,\lambda _2]`$. In each example one action has a form $`\alpha _{\overline{\lambda }}^{\mathrm{min}}`$ and the other $`\alpha _{\overline{\lambda }}^{\mathrm{max}}`$. Hence by Corollary 5.10 they are not measurably isomorphic up to a time change
In other words, in each example one action, namely, $`\alpha _{\overline{\lambda }}^{\mathrm{min}}`$, is a cyclic Cartan action, and the other is not.
We will aslo show that in these examples the conjugacy type of the pair $`(Z(\alpha ),\alpha )`$ distinguishes weakly isomorphic actions. Let us point out that a noncylic action for example $`\alpha _{\overline{\lambda }}^{\mathrm{max}}`$ may be maximal, for example when fundamental units lie in a proper subring of $`𝒪_K`$. However in our examples centralizers for the cyclic actions will be dirrefern and thus will serve as a distuinguishing invariant.
The information about cubic fields is either taken from or obtained with the help of the computer package Pari-GP. Some calculations were made by Arsen Elkin during the REU program at Penn State in summer of 1999.
We construct two $`^2`$–actions, $`\alpha `$, generated by commuting matrices $`A`$ and $`B`$, and $`\alpha ^{}`$, generated by commuting matrices $`A^{}`$ and $`B^{}`$ in $`GL(3,)`$. These actions are weakly algebraically isomorphic by Proposition 3.8 since they are produced with the same set of units on two different orders, $`[\lambda ]`$ and $`𝒪_K`$, but not algebraically isomorphic by Proposition 3.9. In these examples the action $`\alpha `$ is cyclic by Corollary 3.10 and will be shown to be a maximal Cartan action. Thus $`Z(\alpha )=\alpha \times \{\pm \mathrm{Id}\}`$. The action $`\alpha ^{}`$ is not maximal, specifically, $`Z(\alpha ^{})/\{\pm \mathrm{Id}\}`$ is a nontrivial finite extension of $`\alpha ^{}`$.
Example 2a. Let $`K`$ be a totally real cubic field given by the irreducible polynomial $`f(x)=x^3+3x^26x+1`$, i.e. $`K=(\lambda )`$ where $`\lambda `$ is one of its roots. The discriminant of $`K`$ is equal to $`81`$, hence its Galois group is cyclic, and $`[𝒪_K:[\lambda ]]=3`$. The algebraic integers $`\lambda _1=\lambda `$ and $`\lambda _2=24\lambda \lambda ^2`$ are units with $`f(\lambda _1)=f(\lambda _2)=0`$. The minimal order in $`K`$ containing $`\lambda _1`$ and $`\lambda _2`$ is $`[\lambda _1,\lambda _2]=[\lambda ]`$, and the maximal order is $`𝒪_K`$. A basis in fundamental units is $`ϵ=\frac{\lambda ^2+5\lambda +1}{3}`$ and $`ϵ1`$, hence $`𝒰_K`$ is not contained in $`[\lambda ]`$.
With respect to the basis $`\{1,\lambda ,\lambda ^2\}`$ in $`[\lambda ]`$, multiplications by $`\lambda _1`$ and $`\lambda _2`$ are given by the matrices
$$A=\left(\begin{array}{ccc}0& 1& 0\\ 0& 0& 1\\ 1& 6& 3\end{array}\right),B=\left(\begin{array}{ccc}2& 4& 1\\ 1& 4& 1\\ 1& 5& 1\end{array}\right),$$
respectively (if acting from the right on row–vectors). A direct calculation shows that this action is maximal.
With respect to the basis $`\{\frac{2}{3}+\frac{5}{3}\lambda +\frac{1}{3}\lambda ^2,\frac{1}{3}+\frac{7}{3}\lambda +\frac{2}{3}\lambda ^2\}`$ in $`𝒪_K`$, multiplications by $`\lambda _1`$ and $`\lambda _2`$ are given by the matrices
$$A^{}=\left(\begin{array}{ccc}1& 2& 1\\ 1& 2& 2\\ 2& 5& 2\end{array}\right),B^{}=\left(\begin{array}{ccc}1& 1& 1\\ 1& 2& 1\\ 1& 4& 2\end{array}\right).$$
We have $`A^{}=VAV^1`$, $`B^{}=VBV^1`$ for $`V=\left(\begin{array}{ccc}2& 2& 1\\ 0& 3& 0\\ 1& 4& 2\end{array}\right)`$. Since $`A`$ is a companion matrix of $`f`$, $`\alpha =A,B`$ has a cyclic element in $`^3`$. If $`A^{}`$ also had a cyclic element $`𝐦=(m_1,m_2,m_3)^3`$, then the vectors
$$\begin{array}{c}𝐦=(m_1,m_2,m_3),𝐦A^{}=(m_1m_2+2m_3,2m_12m_2+5m_3,m_1+2m_22m_3)\\ 𝐦\left(A^{}\right)^2=(3m_1+5m_27m_3,7m_1+12m_216m_3,5m_17m_2+12m_3),\end{array}$$
would have to generate $`^3`$ or, equivalently
$$det\left(\begin{array}{ccc}m_1& m_2& m_3\\ m_1m_2+2m_3& 2m_12m_2+5m_3& m_1+2m_22m_3\\ 3m_1+5m_27m_3& 7m_1+12m_216m_3& 5m_17m_2+12m_3\end{array}\right)$$
$$\begin{array}{cc}\hfill =3m_1^3& +18m_1^2m_39m_1m_2^29m_1m_2m_3\hfill \\ & +27m_1m_3^2+3m_2^39m_2m_3^2+3m_3^3=1.\hfill \end{array}$$
This contradiction shows that $`A^{}`$ has no cyclic vector, and since $`B^{}=24A^{}A_{}^{}{}_{}{}^{2}`$ , the action $`\alpha ^{}`$ is not cyclic. In this example both actions $`\alpha `$ and $`\alpha ^{}`$ have a single fixed point $`(0,0,0)`$, hence their linear and affine centralizers coincide, and by Corollary 5.3 $`\alpha `$ and $`\alpha ^{}`$ are not measurably isomorphic up to a time change.
The action $`\alpha ^{}`$ is not maximal beacuse $`Z(\alpha ^{})`$ contains fundamental units.
Example 2b. Let us consider a totally real cubic field $`K`$ given by the irreducible polynomial $`f(x)=x^37x^2+11x1`$. Thus $`K=(\lambda )`$ where $`\lambda `$ is one of its roots. In this field the ring of integers $`𝒪_K`$ has basis $`\{1,\lambda ,\frac{1}{2}\lambda ^2+\frac{1}{2}\}`$ and hence $`[𝒪_K:[\lambda ]]=2`$. The fundamental units in $`𝒪_K`$ are $`\{\frac{1}{2}\lambda ^22\lambda +\frac{1}{2},\lambda 2\}`$. We choose the units $`\lambda =\lambda _1=(\frac{1}{2}\lambda ^22\lambda +\frac{1}{2})^2`$ and $`\lambda _2=\lambda 2`$ which are contained in both orders, $`𝒪_K`$ and $`[\lambda ]`$.
In $`[\lambda ]`$ we consider the basis $`\{1,\lambda ,\lambda ^2\}`$ relative to which the multiplication by $`\lambda _1`$ is represented by the companion matrix $`A=\left(\begin{array}{ccc}0& 1& 0\\ 0& 0& 1\\ 1& 11& 7\end{array}\right)`$ and multiplication by $`\lambda _2`$ is represented by the matrix $`B=\left(\begin{array}{ccc}2& 1& 0\\ 0& 2& 1\\ 1& 11& 5\end{array}\right)`$.
For $`𝒪_K`$ with the basis $`\{1,\lambda ,\frac{1}{2}\lambda ^2+\frac{1}{2}\}`$ multiplications by $`\lambda _1`$ and $`\lambda _2`$ are represented by the matrices $`A^{}=\left(\begin{array}{ccc}0& 1& 0\\ 1& 0& 2\\ 3& 5& 7\end{array}\right)`$ and $`B^{}=\left(\begin{array}{ccc}2& 1& 0\\ 1& 2& 2\\ 3& 5& 5\end{array}\right)`$.
It can be seen directly that $`\alpha `$ and $`\alpha ^{}`$ are not algebraically conjugate up to a time change since $`A^{}`$ is a square of a matrix from $`SL(3,)`$: $`A^{}=\left(\begin{array}{ccc}0& 2& 1\\ 1& 5& 3\\ 2& 9& 6\end{array}\right)^2`$, while $`A`$ is not a square of a matrix in $`GL(3,)`$, which is checked by reducing modulo $`2`$. In this case it is also easily seen that the action $`\alpha ^{}`$ is not cyclic since the corresponding determinant is divisible by 2. The action $`\alpha `$ has $`2`$ fixed points on $`𝕋^3`$: $`(0,0,0)`$ and $`(\frac{1}{2},\frac{1}{2},\frac{1}{2})`$, while the action $`\alpha ^{}`$ has $`4`$ fixed points: $`(0,0,0)`$, $`(\frac{1}{2},\frac{1}{2},\frac{1}{2})`$, $`(\frac{1}{2},\frac{1}{2},0)`$, and $`(0,0,\frac{1}{2})`$. Hence the affine centralizer of $`\alpha `$ is $`Z(\alpha )\times /2`$, and the affine centralizer of $`\alpha ^{}`$ is $`Z(\alpha ^{})\times (/2\times /2)`$.
By Lemma 4.2, the group of elements of finite order in $`Z_{\mathrm{𝐴𝑓𝑓}}(\alpha )`$ is $`/2\times /2`$ and in $`Z_{\mathrm{𝐴𝑓𝑓}}(\alpha ^{})`$ it is $`/2\times /2\times /2`$. The indices of each action in its affine centralizer are $`[Z_{\mathrm{𝐴𝑓𝑓}}(\alpha ):\alpha ]=4`$ and $`[Z_{\mathrm{𝐴𝑓𝑓}}(\alpha ^{}):\alpha ^{}]=16`$.
This gives two alternative arguments that the actions are not measurably isomorphic up to a time change.
### 6.3. Nonisomorphic maximal Cartan actions
We find examples of weakly algebraically isomorphic maximal Cartan actions which are not algebraically isomorphic up to time change. For such an action $`\alpha `$ the structure of the pair $`(Z(\alpha ),\alpha )`$ is always the same: $`Z(\alpha )`$ is isomorphic as a group to $`\alpha \times \{\pm \mathrm{Id}\}`$. The algebraic tool which allows to distinguish the actions is Theorem 4.5.
Example 3a. An example for $`n=3`$ can be obtained from a totally real cubic field with class number $`2`$ and the Galois group $`S_3`$. The smallest discriminant for such a field is $`1957`$ (, Table B4), and it can be represented as $`K=(\lambda )`$ where $`\lambda `$ is a unit in $`K`$ with minimal polynomial $`f(x)=x^32x^28x1`$. In this field the ring of integers $`𝒪_K=[\lambda ]`$ and the fundamental units are $`\lambda _1=\lambda `$ and $`\lambda _2=\lambda +2`$. Two actions are constructed with this set of units (fundamental, hence multiplicatively independent) on two different lattices, $`𝒪_K`$ with the basis $`\{1,\lambda ,\lambda ^2\}`$, representing the principal ideal class, and $``$ with the basis $`\{2,1+\lambda ,1+\lambda ^2\}`$ representing to the second ideal class. Notice that the units $`\lambda _1`$ and $`\lambda _2`$ do not belong to $``$, but $``$ is a $`[\lambda ]`$-module. The first action $`\alpha `$ is generated by the matrices $`A=\left(\begin{array}{ccc}0& 1& 0\\ 0& 0& 1\\ 1& 8& 2\end{array}\right)`$ and $`B=\left(\begin{array}{ccc}2& 1& 0\\ 0& 2& 1\\ 1& 8& 4\end{array}\right)`$ which represent multiplication by $`\lambda _1`$ and $`\lambda _2`$, respectively, on $`𝒪_K`$. The second action $`\alpha ^{}`$ is generated by matrices $`A^{}=\left(\begin{array}{ccc}1& 2& 0\\ 1& 1& 1\\ 5& 9& 2\end{array}\right)`$ and $`B^{}=\left(\begin{array}{ccc}1& 2& 0\\ 1& 3& 1\\ 5& 9& 5\end{array}\right)`$ which represent multiplication by $`\lambda _1`$ and $`\lambda _2`$, respectively, on $``$ in the given basis. By Proposition 3.8 these actions are weakly algebraically isomorphic. By Theorem 4.5 they are not algebraically isomorphic. Since the Galois group is $`S_3`$ there are no nontrivial time changes which produce conjugacy over $``$. Therefore, but Theorem 5.2 the actions are not measurably isomorphic.
It is interesting to point out that for actions $`\alpha `$ and $`\alpha ^{}`$ the affine centralizers $`Z_{\mathrm{𝐴𝑓𝑓}}(\alpha )`$ and $`Z_{\mathrm{𝐴𝑓𝑓}}(\alpha ^{})`$ are not isomorphic as abstract groups. The action $`\alpha `$ has $`2`$ fixed points on $`𝕋^3`$: $`(0,0,0)`$ and $`(\frac{1}{2},\frac{1}{2},\frac{1}{2})`$, while the action $`\alpha ^{}`$ has a single fixed point $`(0,0,0)`$. Hence $`Z_{\mathrm{𝐴𝑓𝑓}}(\alpha )`$ is isomorphic to $`Z(\alpha )\times /2`$, $`Z_{\mathrm{𝐴𝑓𝑓}}(\alpha ^{})`$ is isomorphic to $`Z(\alpha ^{})`$. As abstract groups, $`Z_{\mathrm{𝐴𝑓𝑓}}(\alpha )^2\times /2\times /2`$ and $`Z_{\mathrm{𝐴𝑓𝑓}}(\alpha ^{})^2\times /2`$.
Hence by Corollary 5.4 the measurable centralizers of $`\alpha `$ and $`\alpha ^{}`$ are not conjugate in the group of measure–preserving transformation providing a distinguishing invariant of measurable isomorphism.
Example 3b. This example is obtained from a totally real cubic field with class number $`3`$, Galois group $`S_3`$, and discriminant $`2597`$. It can be represented as $`K=(\lambda )`$ where $`\lambda `$ is a unit in $`K`$ with minimal polynomial $`f(x)=x^32x^28x+1`$. In this field the ring of integers $`𝒪_K=[\lambda ]`$ and the fundamental units are $`\lambda _1=\lambda `$ and $`\lambda _2=\lambda +2`$. Three actions are constructed with this set of units on three different lattices, $`𝒪_K`$ with the basis $`\{1,\lambda ,\lambda ^2\}`$, representing the principal ideal class, $``$ with the basis $`\{2,1+\lambda ,1+\lambda ^2\}`$ representing the second ideal class, and $`^2`$ with the basis $`\{4,3+\lambda ,3+\lambda ^2\}`$ representing the third ideal class.
Multiplications by $`\lambda _1`$ and $`\lambda _2`$ generate the following three weakly algebraically isomorphic actions which are not algebraically isomorphic by Theorem 4.5 even up to a time change, and therefore not measurably isomorphic:
$$A=\left(\begin{array}{ccc}0& 1& 0\\ 0& 0& 1\\ 1& 8& 2\end{array}\right)\text{and}B=\left(\begin{array}{ccc}2& 1& 0\\ 0& 2& 1\\ 1& 8& 4\end{array}\right);$$
$$A^{}=\left(\begin{array}{ccc}1& 2& 0\\ 1& 1& 1\\ 6& 9& 2\end{array}\right)\text{and}B^{}=\left(\begin{array}{ccc}1& 2& 0\\ 1& 3& 1\\ 6& 9& 4\end{array}\right);$$
$$A^{\prime \prime }=\left(\begin{array}{ccc}3& 4& 0\\ 3& 3& 1\\ 10& 11& 2\end{array}\right)\text{and}B^{\prime \prime }=\left(\begin{array}{ccc}1& 4& 0\\ 3& 5& 1\\ 10& 11& 4\end{array}\right).$$
Each action has $`2`$ fixed point in $`𝕋^3`$, $`(0,0,0)`$ and $`(\frac{1}{2},\frac{1}{2},\frac{1}{2})`$. Hence all affine centralizers are isomorphic as abstract groups to $`^2\times /2\times /2`$.
Example 3c Finally we give an example of two nonisomorphic maximal Cartan actions which come from the vector of fundamental units $`\overline{\lambda }=(\lambda _1,\lambda _2)`$ in a totally real cubic field $`K`$ such that $`(\lambda _1,\lambda _2)𝒪_K`$. Thus the whole group of units does not generate the ring $`𝒪_K`$. Both actions $`\alpha _{\overline{\lambda }}^{\mathrm{min}}`$ and $`\alpha _{\overline{\lambda }}^{\mathrm{max}}`$ of the group $`^2`$ are maximal Cartan actions by Lemma 3.6. However by Corollary 3.10 the former is cyclic and the latter is not and hence they are not measurably isomorphic up to a time change by Corollary 5.10.
For a specific example we pick the totally real cubic field $`K=(\alpha )`$ with class number $`1`$ discriminant $`1304`$ given by the polynomial $`x^3x^211x1`$. For this filed we have $`[𝒪_K:(\alpha )]=2`$. Generators in $`𝒪_K`$ can be taken to be $`\{1,\alpha ,\beta =\frac{\alpha ^2+1}{2}\}`$. Fundamental units are $`\lambda _1=\alpha ,\lambda _2=5+14\alpha +10\beta =14\alpha +5\alpha ^2[\alpha ]`$. Thus the whole group of units lies in $`[\lambda ]`$. To construct the generators for two non–isomorphic action $`\alpha _{\overline{\lambda }}^{\mathrm{min}}`$ and $`\alpha _{\overline{\lambda }}^{\mathrm{max}}`$ we write multiplications by $`\lambda _1`$ and $`\lambda _2`$ in bases $`\{1,\alpha ,\alpha ^2\}`$ and $`\{1,\alpha ,\beta \}`$, correspondingly. The resulting matrices are:
$$A=\left(\begin{array}{ccc}0& 1& 0\\ 1& 0& 1\\ 1& 11& 1\end{array}\right)B=\left(\begin{array}{ccc}0& 14& 5\\ 5& 55& 19\\ 19& 214& 74\end{array}\right),$$
$$A^{}=\left(\begin{array}{ccc}0& 1& 0\\ 1& 0& 2\\ 0& 6& 1\end{array}\right)B=\left(\begin{array}{ccc}5& 14& 10\\ 14& 55& 38\\ 30& 114& 79\end{array}\right).$$
The first action has only one fixed point, the origin; the second has four fixed points $`(0,0,0)`$, $`(\frac{1}{2},\frac{1}{2},\frac{1}{2})`$, $`(\frac{1}{2},\frac{1}{2},0)`$, and $`(0,0,\frac{1}{2})`$. Thus we have an example of two maximal Cartan actions of $`^2`$ which have nonisomorphic affine and hence measurable centralizers.
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# Untitled Document
ON THE ADM EQUATIONS FOR GENERAL RELATIVITY
Giampiero EspositoandCosimo Stornaiolo
INFN, Sezione di Napoli
Complesso Universitario di Monte S. Angelo
Via Cintia, Edificio N’
80126 Napoli, Italy
and
Università di Napoli Federico II
Dipartimento di Scienze Fisiche
Complesso Universitario di Monte S. Angelo
Via Cintia, Edificio N’
80126 Napoli, Italy
The Arnowitt–Deser–Misner (ADM) evolution equations for the induced metric and the extrinsic-curvature tensor of the spacelike surfaces which foliate the space-time manifold in canonical general relativity are a first-order system of quasi-linear partial differential equations, supplemented by the constraint equations. Such equations are here mapped into another first-order system. In particular, an evolution equation for the trace of the extrinsic-curvature tensor $`K`$ is obtained whose solution is related to a discrete spectral resolution of a three-dimensional elliptic operator $`𝒫`$ of Laplace type. Interestingly, all nonlinearities of the original equations give rise to the potential term in $`𝒫`$. An example of this construction is given in the case of a closed Friedmann–Lemaitre–Robertson–Walker universe. Eventually, the ADM equations are re-expressed as a coupled first-order system for the induced metric and the trace-free part of $`K`$. Such a system is written in a form which clarifies how a set of first-order differential operators and their inverses, jointly with spectral resolutions of operators of Laplace type, contribute to solving, at least in principle, the original ADM system.
The canonical formulation of general relativity relies on the assumption that space-time $`(M,g)`$ is topologically $`\mathrm{\Sigma }\times 𝐑`$ and can be foliated by a family of spacelike surfaces $`\mathrm{\Sigma }_t`$, all diffeomorphic to the three-manifold $`\mathrm{\Sigma }`$. The space-time metric $`g`$ is then locally cast in the form
$$g=(N^2N_iN^i)dtdt+N_i(dx^idt+dtdx^i)+h_{ij}dx^idx^j,$$
$`(1)`$
where $`N`$ is the lapse function and $`N^i`$ are components of the shift vector of the foliation \[1–3\]. The induced metric $`h_{ij}`$ on $`\mathrm{\Sigma }_t`$ and the associated extrinsic-curvature tensor $`K_{ij}`$ turn out to obey the first-order equations \[1–3\]
$$\frac{h_{ij}}{t}=2NK_{ij}+N_{ij}+N_{ji},$$
$`(2)`$
$$\begin{array}{ccc}& \frac{K_{ij}}{t}=N_{ij}+N\left[{}_{}{}^{(3)}R_{ij}^{}+K_{ij}(\mathrm{tr}K)2K_{im}K_j^m\right]\hfill & \\ & +\left[N^mK_{ijm}+N_i^mK_{jm}+N_j^mK_{im}\right],\hfill & (3)\hfill \end{array}$$
where the stroke denotes covariant differentiation with respect to the induced connection on $`\mathrm{\Sigma }_t`$. Moreover, the constraint equations hold. For Einstein theory in vacuum they read
$$K_{il}^i(\mathrm{tr}K)_l0,$$
$`(4)`$
$${}_{}{}^{(3)}RK_{ij}K^{ij}+(\mathrm{tr}K)^20,$$
$`(5)`$
where $``$ is the weak-equality symbol introduced by Dirac to denote equations which only hold on the constraint surface .
Equation (3) is quasi-linear in that it contains terms quadratic in the extrinsic-curvature tensor, i.e.
$$N\left(K_{ij}(\mathrm{tr}K)2K_{im}K_j^m\right).$$
We are now aiming to obtain from eqs. (2) and (3) another set of quasi-linear first-order equations. For this purpose, we contract Eq. (2) with $`K^{ij}`$ and Eq. (3) with $`h^{ij}`$. This leads to (hereafter $`_i^i`$ is the Laplacian on $`\mathrm{\Sigma }_t`$)
$$K^{ij}\frac{h_{ij}}{t}=2NK_{ij}K^{ij}+2K^{ij}N_{ij},$$
$`(6)`$
$$\begin{array}{ccc}& h^{ij}\frac{K_{ij}}{t}=N+N[{}_{}{}^{(3)}R+(\mathrm{tr}K)^22K_{im}K^{im}]\hfill & \\ & +\left[N^m(\mathrm{tr}K)_m+2N_i^mK_m^i\right].\hfill & (7)\hfill \end{array}$$
It is now possible to obtain an evolution equation for $`(\mathrm{tr}K)`$, because
$$\frac{}{t}(\mathrm{tr}K)=\frac{h^{ij}}{t}K_{ij}+h^{ij}\frac{K_{ij}}{t}.$$
$`(8)`$
The second term on the right-hand side of Eq. (8) is given by Eq. (7), whereas the first term is obtained after using the identity
$$\frac{}{t}(h^{ij}h_{jl})=\frac{}{t}\delta _l^i=0,$$
$`(9)`$
which implies
$$\frac{h^{ij}}{t}=h^{ip}h^{jl}\frac{h_{pl}}{t},$$
$`(10)`$
and hence
$$\frac{h^{ij}}{t}K_{ij}=\frac{h_{ij}}{t}K^{ij}.$$
$`(11)`$
By virtue of Eqs. (6)–(8) and (11), and imposing the constraint equation (5), we get
$$\frac{}{t}(\mathrm{tr}K)=(K_{ij}K^{ij})N+N^m(\mathrm{tr}K)_m,$$
$`(12a)`$
which can also be cast in the form (here $`\widehat{}_m_m`$)
$$(\frac{}{t}N^m\widehat{}_m)(\mathrm{tr}K)=(+K_{ij}K^{ij})N.$$
$`(12b)`$
This is a first non-trivial result because it tells us that, given the operator of Laplace type
$$𝒫+K_{ij}K^{ij},$$
$`(13)`$
the right-hand side of Eq. (12b) is determined by a discrete spectral resolution of $`𝒫`$. By this one means a complete orthonormal set of eigenfunctions $`f_\lambda ^p`$ belonging to the eigenvalue $`\lambda `$, so that, for each fixed value of $`t`$, the lapse can be expanded in the form
$$N(\stackrel{}{x},t)=\underset{\lambda }{}C_\lambda f_\lambda ^p(\stackrel{}{x},t),$$
$`(14)`$
with Fourier coefficients $`C_\lambda `$ given by the scalar product
$$C_\lambda =(f_\lambda ^p,N).$$
$`(15a)`$
This point is simple but non-trivial: for each fixed value of $`t`$, the restriction of the lapse to $`\mathrm{\Sigma }_t`$ becomes a function on $`\mathrm{\Sigma }_t`$ only, and hence the Fourier coefficients in (15a) read
$$C_\lambda =_{\mathrm{\Sigma }_t}(f_\lambda ^p)^{}N\sqrt{h}d^3xC_{\lambda ,t},$$
$`(15b)`$
where the star denotes complex conjugation. In the application to the initial-value problem, one studies the lapse at different values of $`t`$, and hence it is better to write its expansion in the form (14), with the time parameter explicitly included. The integration measure for $`C_\lambda `$, however, remains the invariant integration measure on $`\mathrm{\Sigma }_t`$.
The existence of discrete spectral resolutions of $`𝒫`$ is guaranteed if $`\mathrm{\Sigma }_t`$ is a compact Riemannian manifold without boundary , which is what we assume hereafter. Thus, Eq. (12b) can be re-expressed in the form
$$\left(\frac{}{t}N^m\widehat{}_m\right)\mathrm{tr}K(\stackrel{}{x},t)=\underset{\lambda }{}\lambda C_{\lambda ,t}f_\lambda ^p(\stackrel{}{x},t).$$
$`(16)`$
On denoting by $`L`$ the operator
$$L\frac{}{t}N^m\widehat{}_m,$$
$`(17)`$
and writing $`G(\stackrel{}{x},t;\stackrel{}{y},t^{})`$ for the Green function of $`L`$, Eq. (16) can be therefore solved for $`(\mathrm{tr}K)`$ in the form
$$\mathrm{tr}K(\stackrel{}{x},t)=_{\mathrm{\Sigma }_t^{}\times 𝐑}G(\stackrel{}{x},t;\stackrel{}{y},t^{})\underset{\lambda }{}\lambda C_{\lambda ,t^{}}f_\lambda ^p(\stackrel{}{y},t^{})d\mu (\stackrel{}{y},t^{}),$$
$`(18)`$
where $`d\mu (\stackrel{}{y},t^{})`$ is the invariant integration measure on $`\mathrm{\Sigma }_t^{}\times 𝐑`$. Note that the eigenfunctions $`f_\lambda ^p`$ are different at different times and hence, to compute them at all times, one has to integrate the equations of motion so that the metric and the extrinsic-curvature tensor are known at all times. Thus, Eq. (18) does not reduce the amount of entanglement of the original equations (2) and (3). A relevant particular case is obtained on considering the canonical form of the foliation , for which the shift vector vanishes. The operator $`L`$ reduces then to $`\frac{}{t}`$, and its Green function can be chosen to be of the form
$$G_c(t,t^{})=\frac{1}{2}[\theta (tt^{})\theta (t^{}t)],$$
$`(19)`$
where $`\theta `$ is the step function. Such a Green function equals $`\frac{1}{2}`$ for $`t>t^{}`$ and $`\frac{1}{2}`$ for $`t<t^{}`$, and hence its “jump” at $`t^{}`$ equals $`1`$, as it should be from the general theory. It is sometimes considered in a completely different branch of physics, i.e. the theory of solitons and nonlinear evolution equations.
A cosmological example shows how the right-hand side of Eq. (16) can be worked out explicitly in some cases. We are here concerned with closed Friedmann–Lemaitre–Robertson–Walker models, for which the three-manifold $`\mathrm{\Sigma }`$ reduces to a three-sphere of radius $`a`$. It is indeed well known that the space of functions on the three-sphere can be decomposed by using the invariant subspaces corresponding to irreducible representations of $`O(4)`$, the orthogonal group in four dimensions. These invariant subspaces are spanned by the hyperspherical harmonics , i.e. generalizations to the three-sphere of the familiar spherical harmonics. The scalar hyperspherical harmonics $`Q^{(n)}(\chi ,\theta ,\phi )`$ form a basis spanning the invariant subspace labeled by the integer $`n1`$ corresponding to the $`n`$-th scalar representation of $`O(4)`$. The label $`(n)`$ of $`Q^{(n)}`$ refers both to the order $`n`$ and to other labels denoting the different elements spanning the subspace. The number of elements $`Q^{(n)}`$ is determined by the dimension of the corresponding $`O(4)`$ representation. On using the general formula for $`O(2k)`$ representations, the dimension of the irreducible scalar representations is found to be $`d_Q(n)=n^2`$. The Laplacian on a unit three-sphere when acting on $`Q^{(n)}`$ gives
$$Q^{(n)}=(n^21)Q^{(n)}n=1,2,\mathrm{}.$$
$`(20)`$
Any arbitrary function $`F`$ on the three-sphere can be expanded in terms of these hyperspherical harmonics as they form a complete set:
$$F=\underset{(n)}{}q_{(n)}Q^{(n)},$$
$`(21)`$
where the sum $`(n)`$ runs over $`n=1`$ to $`\mathrm{}`$ and over the $`n^2`$ elements in each invariant subspace, and the $`q_{(n)}`$ are constant coefficients. The lapse function on the three-sphere is therefore expanded according to (cf. (14))
$$N=\underset{(n)}{}c_{(n)}Q^{(n)},$$
$`(22)`$
and bearing in mind that $`K_{ij}K^{ij}=3`$ on a unit three-sphere we find, on a three-sphere of radius $`a`$,
$$𝒫N=\underset{(n)}{}c_{(n)}\lambda _nQ^{(n)},$$
$`(23)`$
where (here $`a=a(t)`$)
$$\lambda _n=\frac{(n^2+2)}{a^2}n=1,2,\mathrm{}.$$
$`(24)`$
We have therefore found that all nonlinearities of the original equations give rise to the potential term in the operator $`𝒫`$ defined in Eq. (13), which is a positive-definite operator of Laplace type (unlike the Laplacian, which is bounded from below but has also a zero eigenvalue when $`n=1`$).
The interest in the evolution equation for the trace of the extrinsic-curvature tensor is motivated by a careful analysis of the variational problem in general relativity. More precisely, the “cosmological action” is sometimes considered, which is the form of the action functional $`I`$ appropriate for the case when $`(\mathrm{tr}K)`$ and the conformal three-metric $`\stackrel{~}{h}_{ij}h^{\frac{1}{3}}h_{ij}`$ are fixed on the boundary. One then finds in $`c=1`$ units
$$I=\frac{1}{16\pi G}_M{}_{}{}^{(4)}R\sqrt{g}d^4x+\frac{1}{24\pi G}_M(\mathrm{tr}K)\sqrt{h}d^3x,$$
$`(25)`$
and this formula finds important applications also to the quantum cosmology of a closed universe, giving rise to the $`K`$-representation of the wave function of the universe .
The program we have outlined in our letter shows an intriguing link between the Cauchy problem in general relativity on the one hand \[10–13\], and the spectral theory of elliptic operators of Laplace type on the other hand , which is obtained by exploiting the nonlinear form of the ADM equations and the nonpolynomial form of the Hamiltonian constraint (5). On denoting by $`\sigma _{ij}`$ the trace-free part of the extrinsic-curvature tensor:
$$\sigma _{ij}K_{ij}\frac{1}{3}h_{ij}(\mathrm{tr}K)$$
$`(26)`$
we therefore obtain a set of equations, equivalent to (2) and (3), in the form (here $`L^1`$ denotes the inverse of the operator $`L`$ defined in (17))
$$\mathrm{tr}K=L^1𝒫N,$$
$`(27)`$
$$\frac{h_{ij}}{t}=\frac{2}{3}Nh_{ij}(\mathrm{tr}K)2N\sigma _{ij}+2N_{(ij)},$$
$`(28)`$
$$\frac{\sigma _{ij}}{t}=\mathrm{\Omega }_{ij}N\frac{2}{3}N_{(ij)}(\mathrm{tr}K)+N^m\widehat{}_m\sigma _{ij}+2N_{(i}^m[\sigma _{j)m}+\frac{1}{3}h_{j)m}(\mathrm{tr}K)],$$
$`(29)`$
where the operator of Laplace type $`\mathrm{\Omega }_{ij}`$ is given by
$$\begin{array}{ccc}\hfill \mathrm{\Omega }_{ij}& \widehat{}_j\widehat{}_i\frac{1}{3}h_{ij}𝒫+\frac{1}{3}\sigma _{ij}(\mathrm{tr}K)\hfill & \\ & +\frac{1}{3}h_{ij}(\mathrm{tr}K)^2+{}_{}{}^{(3)}R_{ij}^{}2\sigma _{im}\sigma _j^m.\hfill & (30)\hfill \end{array}$$
It now appears desirable to understand to which extent the differential-operator viewpoint plays a role in solving the coupled system (27)–(30). For this purpose, we consider a relevant particular case, i.e. the canonical foliation studied in Ref. , for which the shift vector vanishes. The operator $`L`$ reduces then to $`\frac{}{t}`$ and, on defining the operators
$$Q\frac{}{t}+\frac{2}{3}N\mathrm{tr}K=L+\frac{2}{3}NL^1𝒫N,$$
$`(31)`$
$$S\frac{}{t}\frac{N}{3}\mathrm{tr}K=L\frac{N}{3}L^1𝒫N,$$
$`(32)`$
$$A_{ij}\widehat{}_j\widehat{}_i+{}_{}{}^{(3)}R_{ij}^{},$$
$`(33)`$
one finds the equations
$$\left[S+\frac{2}{3}(L^1𝒫N)^2Q^1N\right]\sigma _{ij}\frac{2}{3}(Q^1N\sigma _{ij})𝒫N+2N\sigma _{im}\sigma _j^m=A_{ij}N,$$
$`(34)`$
$$h_{ij}=2Q^1N\sigma _{ij},$$
$`(35)`$
supplemented, of course, by Eq. (27). In other words, once Eq. (34) is solved for the trace-free part of the extrinsic-curvature tensor of $`\mathrm{\Sigma }_t`$, the induced metric on $`\mathrm{\Sigma }_t`$ is obtained from Eq. (35). This form of the coupled first-order ADM system of equations has possibly the merit of stressing the need to invert the operators $`L,Q`$ and $`S`$ to find a solution for given initial conditions. Approximate solutions, to the desired accuracy, will correspond to the construction of approximate inverse operators $`L^1,Q^1`$ and $`S^1`$. For this purpose, it can be useful to re-express Eq. (34) in the form
$$\begin{array}{ccc}& \left[I+\frac{2}{3}S^1(L^1𝒫N)^2Q^1N\right]\sigma _{ij}\frac{2}{3}S^1\left[(Q^1N\sigma _{ij})𝒫N\right]\hfill & \\ & S^1(A_{ij}N)=2S^1\left(N\sigma _{im}\sigma _j^m\right).\hfill & (34^{})\hfill \end{array}$$
The action of the inverse operator $`L^1`$ on $`𝒫N`$ is already given by Eq. (18), but the inverses $`Q^1`$ and $`S^1`$ make it necessary to develop an algorithm for their exact or approximate evaluation. More precisely, one has
$$Q=L\left(I+\frac{2}{3}L^1NL^1𝒫N\right),$$
$`(31^{})`$
$$S=L\left(I\frac{1}{3}L^1NL^1𝒫N\right),$$
$`(32^{})`$
and hence
$$Q^1=\left(I+\frac{2}{3}L^1NL^1𝒫N\right)^1L^1,$$
$`(36)`$
$$S^1=\left(I\frac{1}{3}L^1NL^1𝒫N\right)^1L^1,$$
$`(37)`$
by virtue of the operator identity $`(AB)^1=B^1A^1`$. From this point of view, two levels of approximation seem to emerge: (i) The degree of accuracy in the evaluation of $`Q^1`$ and $`S^1`$. It is clear from Eqs. (36) and (37) that this reduces to finding
$$\left(I+\rho L^1NL^1𝒫N\right)^1$$
where $`\rho =\frac{2}{3}`$ or $`\frac{1}{3}`$. A series expansion will be, in general, only of formal value. However, on using the norms defined in section 7.4 of Ref. , which relies on well known properties of Sobolev spaces, if the operator
$$T_\rho \rho L^1NL^1𝒫N$$
$`(38)`$
is a bounded operator on a Banach space $`X`$ with norm $`T_\rho <1`$, the inverse of $`I+T\rho `$ exists and the series $`_{n=0}^{\mathrm{}}(1)^n(T_\rho )^n`$ converges uniformly to $`(I+T_\rho )^1`$ with respect to the norm on the set of bounded maps from $`X`$ into $`X`$, and one can write
$$(I+T_\rho )^1=\underset{n=0}{\overset{\mathrm{}}{}}(1)^n(T_\rho )^n.$$
$`(39)`$
Equation (39) can be therefore used if $`L^1`$ is such that $`T_\rho `$ is a bounded operator with $`T_\rho <1`$. (ii) The way in which the non-linear term $`\sigma _{im}\sigma _j^m`$ is dealt with. For example, one may regard the right-hand side of Eq. (34’) as the “known term”, and hence consider the equation
$$\begin{array}{ccc}& \sigma _{ij}dx^idx^j\frac{2}{3}F^1S^1\left[(Q^1N\sigma _{ij})𝒫N\right]dx^idx^j\hfill & \\ & F^1S^1(A_{ij}N)dx^idx^j=2F^1S^1\left(N\sigma _{im}\sigma _j^m\right)dx^idx^j,\hfill & (40)\hfill \end{array}$$
where $`F^1`$ is the inverse of the operator
$$FI+\frac{2}{3}S^1(L^1𝒫N)^2Q^1N.$$
$`(41)`$
Equation (40) is an integral equation for $`\sigma _{ij}`$, for which a perturbation approach might be useful if one takes
$$2F^1S^1\left(N\sigma _{im}\sigma _j^m\right)dx^idx^j$$
as the known term mentioned before. Note that we have resorted to the use of the tensor product $`dx^idx^j`$ because in the resulting equation (40) one has a well defined integration of a “symmetric two-form” over the space-time manifold, here taken to be diffeomorphic to $`\mathrm{\Sigma }_t\times 𝐑`$. Strictly, also Eq. (35) should be written as
$$h_{ij}dx^idx^j=2(Q^1N\sigma _{ij})dx^idx^j,$$
$`(35^{})`$
and in all such equations only the symmetric part of the tensor product survives, because both $`h_{ij}`$ and $`\sigma _{ij}`$ are symmetric rank-two tensor fields.
Note once more that it would be wrong to regard Eq. (27) as independent of the solution of Eqs. (28) and (29), because $`L^1𝒫N`$ is only known at all times when the time evolution of $`h_{ij}`$ and $`K_{ij}`$ has been determined for given initial conditions. Our equations (35’)–(37) and (40), although not obviously more powerful than previous schemes, prepare the ground for an operator approach to the ADM equations for general relativity, and hence might contribute to the understanding of structural properties of general relativity. In particular, our way of writing the coupled system of non-linear evolution equations might lead to a better understanding of the interplay between elliptic operators on the spacelike surfaces $`\mathrm{\Sigma }_t`$ and hyperbolic equations on the space-time manifold. As far as we can see, this expectation is supported by the closed Friedmann–Lemaitre–Robertson–Walker model discussed before (where $`\sigma _{ij}`$ vanishes), and further examples of cosmological interest might be found, e.g. Bianchi IX models describing a closed but anisotropic universe. For open universes or asymptotically flat space-times, however, we are not aware of theorems that make it possible to expand the lapse as in (14). In such cases, the counterpart of the elliptic theory on $`\mathrm{\Sigma }_t`$ advocated in our paper is therefore another open problem. Acknowledgments. We are grateful to Gabriele Gionti and Giuseppe Marmo for several enlightening conversations, and to Luca Lusanna for very useful comments.
REFERENCES
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4. P. A. M. Dirac, Lectures on Quantum Mechanics, Belfer Graduate School (Yeshiva University, New York, 1964).
5. P. B. Gilkey, Invariance Theory, the Heat Equation and the Atiyah–Singer Index Theorem (CRC Press, Boca Raton, 1995).
6. D. Christodoulou and S. Klainerman, The Global Nonlinear Stability of the Minkowski Space (Princeton University Press, Princeton, 1993).
7. E. M. Lifshitz and I. M. Khalatnikov, Adv. Phys. 12, 185 (1963).
8. J. W. York, Found. Phys. 16, 249 (1986).
9. J. B. Hartle and S. W. Hawking, Phys. Rev. D 28, 2960 (1983).
10. D. Christodoulou, Class. Quantum Grav. 16, A23 (1999).
11. A. Anderson and J. W. York, Phys. Rev. Lett. 82, 4384 (1999).
12. A. Anderson, Y. Choquet-Bruhat and J. W. York, “Einstein’s Equations and Equivalent Hyperbolic Dynamical Systems” (GR-QC 9907099).
13. H. Friedrich and A. Rendall, in Einstein’s Field Equations and their Physical Interpretation, edited by B. G. Schmidt (Springer–Verlag, Berlin, 2000; GR-QC 0002074).
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# Decompositions of Reflexive Modules This work is supported by the project No. G-545-173.06/97 of the German-Israeli Foundation for Scientific Research & Development AMS subject classification: primary: 13C05, 13C10, 13C13, 20K15,20K25,20K30 secondary: 03E05, 03E35 Key words and phrases: almost free modules, reflexive modules, duality theory, modules with particular monomorphism GbSh 716 in Shelah’s list of publications
## 1 Introduction
Let $`R`$ be a countable principal ideal domain with $`10`$ and $`S`$ a multiplicative closed subset of $`R\{0\}`$ containing $`1`$. If $`S=\{s_n:n\omega \},s_0=1`$ and $`q_n=\underset{i<n}{}s_i`$, then
$`q_{n+1}=q_ns_n\text{ and }{\displaystyle \underset{n\omega }{}}q_nR=0.`$ (1.1)
The condition (1.1) requires that $`R`$ is not a field, but then we may choose $`S=R\{0\}`$ or any other ‘classical example’. The right ideals $`q_nR`$ form a basis of the S-topology on $`R`$ which is Hausdorff by (1.1). The S-adic completion of $`R`$ under the $`S`$-topology is the ring $`\widehat{R}`$ which has the size $`2^\mathrm{}_0`$, see Göbel, May for properties on $`\widehat{R}`$. Similar consideration carry over to (free) $`R`$-modules which we will use in the next sections. If $`F`$ is a free $`R`$-module then the $`S`$-topology is generated by $`Fq_n(n\omega )`$ and if $`\widehat{F}`$ denotes the S-adic completion of $`F`$ then $`F\widehat{F}`$ and $`F`$ is pure and dense in $`\widehat{F}`$. Recall that $`F_{}\widehat{F}`$ is pure (w.r.t. $`S`$) if and only if
$$\widehat{F}sFFs\text{ for all }sS.$$
Also $`F`$ is dense in $`\widehat{F}`$ if and only if $`\widehat{F}/F`$ is $`S`$-divisible (in the obvious sense). Note that an element $`e`$ of an $`R`$-module $`G`$ is pure, we write $`e_{}G`$ if $`eR_{}G`$. If $`XG`$, then $`X`$ denotes the submodule generated by $`X`$ and $`X_{}`$ denotes the submodule purely generated by $`X`$. The following observation is well-known and can be looked up in .
###### Observation 1.1
If $`0r_nR`$ infinitely often ($`n\omega `$), then we can find $`ϵ_n\{0,1\}`$ such that
$$\underset{n\omega }{}r_nq_nϵ_n\widehat{R}R.$$
If $`G`$ is any $`R`$-module then $`G^{}=\mathrm{Hom}(G,R)`$ denotes its dual module, and $`G`$ is a dual module if $`GD^{}`$ for some $`R`$-module $`D`$. Particular dual modules are reflexive modules $`D`$ introduced by Bass \[1, p. 476\]. We need the evaluation map
$$\sigma =\sigma _D:DD^{}(d\sigma (d))$$
where $`\sigma (d)D^{}`$ is defined by evaluation
$$\sigma (d):D^{}R(\phi \phi (d)).$$
The module $`D`$ is reflexive if the evaluation map $`\sigma _D`$ is an isomorphism. Obviously $`D(D^{})^{}`$ is a dual module if $`D`$ is reflexive. However, there are very many $`\mathrm{}_1`$-free modules $`G`$ with $`G^{}=0`$, see . An $`R`$-module is $`\mathrm{}_1`$-free if all countable submodules are free. Using the special continuum hypothesis CH the existence of many reflexive modules (Section 3) will follow from considerations of free modules with bilinear form in Section 2. Classical examples are due to Specker and Łos and it follows from Specker’s theorem that reflexive modules must be $`\mathrm{}_1`$-free, see Fuchs . Hence it is not too surprising that we will deal with free $`R`$-modules first in Section 2. The bilinear form is needed to control their duals when passing from $`\mathrm{}_0`$ to size $`\mathrm{}_1`$ in Section 3.
In order to find reflexive groups $`G`$ of cardinality $`\mathrm{}_1`$ with $`G\cong ̸RG`$ we must discard all possible monomorphisms
$`\phi :GG\text{ with }G\phi eR=G.`$ (1.2)
This will be established with the help of an $`\mathrm{}_1`$-filtration
$$G=\underset{\alpha \omega _1}{}G_\alpha \text{ of countable, pure, free submodules }G_\alpha $$
such that $`G_\alpha `$ is a summand with $`R`$-free complement of any $`G_\beta `$ for $`\alpha <\beta `$ if $`\alpha `$ does not belong to a fixed stationary subset $`E`$ of $`\omega _1`$. Given $`\phi `$ as in (1.2) by a back and forth argument there is a cub $`C\omega _1`$ such that $`G_\alpha \phi G_\alpha `$ and $`eG_\alpha `$ for all $`\alpha C`$.
If $`\alpha EC`$ then
$`G_\alpha \phi eR=G_\alpha .`$ (1.3)
And if $`\phi G_\alpha `$ is predicted by a function $`\phi _\alpha `$ as under the assumption of the diamond principle $`\mathrm{}_\mathrm{}_1`$, then we construct $`G_{\alpha +1}`$ by a Step-Lemma from $`G_\alpha `$ and $`\phi _\alpha `$ such that $`\phi _\alpha `$ does not extend to $`G_{\alpha +1}`$. Moreover $`G_{\alpha +1}`$ is the $`S`$-adic closure of $`G_\alpha `$ in $`G`$ by the summand property mentioned above. Hence $`\phi `$ coincides with $`\phi _\alpha `$ on $`G_\alpha `$ and maps $`G_{\alpha +1}`$ into itself, a contradiction. However assuming CH only weaker prediction principles like weak diamond $`\mathrm{\Phi }_\mathrm{}_1`$ are available, see Devlin and Shelah or Eklof and Mekler \[8, p. 143, Lemma 1.7\]. If we discard $`\phi _\alpha `$ in the construction as before, then $`\phi G_\alpha `$ might extend because it does not entirely agree with $`\phi _\alpha `$. In this case $`\phi `$ will show up at some $`\alpha <\beta EC`$ and we have a new chance to discard $`\phi G_{\alpha +1}`$ at level $`\beta `$ \- provided we know $`\phi G_{\alpha +1}`$ when constructing $`G_{\beta +1}`$.
This time we need a stronger algebraic algebraic Step-Lemma: Note that $`G_{\alpha +1}F_\beta =G_\beta `$. Now we have a partial map $`\phi _\beta :=\phi G_{\alpha +1}`$ with domain $`\mathrm{Dom}\phi _\beta =G_{\alpha +1}`$ a summand of $`G_\beta `$ and the splitting property (1.3) for the new map $`\phi _\beta `$:
$`(G_{\alpha +1}\phi _\beta eR)F_\beta =G_\beta .`$ (1.4)
In order to proceed we must discard these partial maps and indeed we are able to prove a generalized Step-Lemma taking care of partial maps (1.4) in Section 2. Moreover, a counting argument also shows that a list of such partial maps
$$\{\phi _\beta :\beta E\}$$
exists which predict any given map $`\phi :GG`$ such that the following holds.
###### Lemma 1.2
Let $`G=_{\alpha \omega _1}G_\alpha `$ be an $`\mathrm{}_1`$-filtration of $`G`$ and $`E\omega _1`$ be a stationary subset of $`\omega _1`$. Then there is a list of predicting partial maps
$$\{\phi _\alpha :G_\alpha G_\alpha :\alpha E\}$$
such that for any countable subset $`A`$ of $`G`$ there is an ordinal $`\beta E`$ (in fact an unbounded set of such ordinals) with
$$\phi A=\phi _\beta (G_\beta A).$$
For a suitable $`A=G_{\alpha +1}`$ and $`\phi A`$ we choose $`\beta E`$ by Lemma 1.2 with $`\phi G_{\alpha +1}=\phi _\beta G_{\alpha +1}`$ and $`G_\beta =G_{\alpha +1}F_\beta `$. Then we construct $`G_{\beta +1}`$ from $`G_\beta `$ and the given $`\phi _\beta G_{\alpha +1}`$ such that $`\phi _\beta G_{\alpha +1}`$ does not extend to $`G_{\beta +1}`$. This contradiction will provide the
###### Main Theorem 1.3
(ZFC + CH) If $`R`$ is a countable domain but not a field, then there is a family of $`2^\mathrm{}_1`$ pair-wise non-isomorphic reflexive $`R`$-modules $`G`$ of cardinality $`\mathrm{}_1`$ such that $`G\cong ̸RG`$.
We also would like to draw attention to a slight modification of the proof of the Main Theorem 1.3. In addition we may assume that $`G`$ in the Main Theorem 1.3 is essentially indecomposable, that any decomposition into two summands has one summand free of finite rank. This follows from a split realization result Corollary 4.3 with $`\mathrm{End}GA\mathrm{Fin}(G)`$ where $`A`$ is any $`R`$-algebra which is free of countable rank. Recall that $`\mathrm{Fin}(G)`$ is the ideal of all endomorphisms of $`G`$ of finite rank.
The formal proof for the prediction Lemma 1.2 is not complicated and uses repeatedly often the weak diamond prediction $`\mathrm{\Phi }_\mathrm{}_1`$. It is also clear from what we said that the underlying module theory is not essential for proving Lemma 1.2 and that it should be possible to replace modules by many other categories like non-commutative groups, fields or Boolean algebras. In order to cover all these possibilities, the prediction principle is formulated in terms of model theory and will appear in this setting in a forth coming book by Shelah \[14, Chapter IX, Claim 1.5\].
We close this introduction with some historical remarks. Using a theorem of Łos̀ (see Fuchs ) on slender groups, the first ‘large’ reflexive abelian groups are free groups or (cartesian) products of $``$ \- assuming for a moment that all cardinals under consideration are $`<\mathrm{}_m`$, the first measurable cardinal. Also the members of the class of groups generated by $``$ and taking direct sums and products alternatively are reflexive, called Reid groups. Using a generalized ‘Chase Lemma’, which controls homomorphisms from products of modules into direct sums of modules, Dugas and Zimmermann-Huisgen showed that the class of Reid groups is ‘really large’. Nevertheless there are more reflexive groups - Eda and Otha applied their ‘theory of continuous functions on $`0`$-dimensional topological spaces’ to find reflexive groups not Reid-groups. As a by-product we also get dual groups which are not reflexive, see also . All these groups $`G`$ have the property that they are either free of finite rank or
$`GG.`$ (1.5)
As indicated in the abstract, we applied $`\mathrm{}_\mathrm{}_1`$ to find examples $`G`$ of size $`\mathrm{}_1`$, where (1.5) is violated. The obvious question to replace $`\mathrm{}_\mathrm{}_1`$ by CH was the main goal of this paper. The question whether the Main Theorem 1.3 holds in any model of ZFC remains open. On the other hand we are able to show that the conclusion of the Main Theorem 1.3 also follows in models of ZFC and Martin’s axiom MA, see . Hence CH is not necessary to derive the existence of these reflexive modules.
## 2 Free modules with bilinear form and partial dual maps
###### Definition 2.1
Let $`(\mathrm{\Phi },𝔉_0,𝔉_1)`$ be a triple of a bilinear map $`\mathrm{\Phi }:F_0F_1R`$ for some countable, free $`R`$-modules $`F_i`$ of infinite rank, $`𝔉_iF_i^{}(i=0,1)`$ families of dual maps subject to the following conditions
1. $`\mathrm{\Phi }`$ is not degenerated. This is to say if $`\mathrm{\Phi }(e,)F_1^{}`$ or $`\mathrm{\Phi }(,f)F_0^{}`$ is the trivial map then $`e=0`$ or $`f=0`$, respectively.
2. $`\mathrm{\Phi }`$ preserves purity, that is $`\mathrm{\Phi }(e,)_{}F_1^{}`$ if $`e_{}F_0`$ and dually $`\mathrm{\Phi }(,f)_{}F_0^{}`$ if $`f_{}F_1.`$
3. $`𝔉_i`$ is a countable, non-empty family of homomorphisms $`\phi :F_\phi R`$ such that $`\mathrm{Dom}\phi =F_\phi _{}F_i`$. The set $`\mathrm{Dom}𝔉_i=\{F_\phi :\phi 𝔉_i\}`$ is well-ordered by inclusion, for $`i\{0,1\}.`$
4. For any $`0xF_1`$ and any finite subset $`E𝔉_0`$ we have $`\mathrm{ker}E\mathrm{ker}\mathrm{\Phi }(,x)`$, and dually for any $`0yF_0`$ and any finite subset $`E𝔉_1`$, we have $`\mathrm{ker}E\mathrm{ker}\mathrm{\Phi }(y,).`$
Here, and in the subsequent parts we use the following
###### Notation 2.2
1. $`𝔉`$ is the collection of all triples $`(\mathrm{\Phi },𝔉_0,𝔉_1)`$ as in Definition 2.1.
2. If $`E\mathrm{Hom}(G,H)`$ then $`\mathrm{ker}E=\underset{\phi E}{}\mathrm{ker}\phi .`$
3. Similarly $`\mathrm{ker}\mathrm{\Phi }(,E)=\underset{eE}{}\mathrm{ker}\mathrm{\Phi }(,e)`$ and dually.
Next we define a partial order on $`𝔉`$.
###### Definition 2.3
1. $`(\mathrm{\Phi },𝔉_0,𝔉_1)(\mathrm{\Phi }^{},𝔉_0^{},𝔉_1^{})`$
2. $`\{\begin{array}{cc}\text{(a) }\mathrm{\Phi }\mathrm{\Phi }^{}\text{ and }\mathrm{Dom}\mathrm{\Phi }_{}\mathrm{Dom}\mathrm{\Phi }^{}\text{ is a pure submodule.}\hfill & \\ \text{(b) }\text{ If }\phi 𝔉_i\text{ then there is a unique }\phi ^{}𝔉_i^{}\text{ such that }\phi \phi ^{}\text{ and }\hfill & \\ \mathrm{Dom}\phi _{}\mathrm{Dom}\phi ^{}.\hfill & \end{array}`$
We will construct the reflexive modules of size $`\mathrm{}_1`$ by using an order preserving continuous map and let
$`p:{}_{}{}^{\omega _1>}2𝔉(\eta p_\eta )`$ (2.1)
from the tree $`𝐓`$ of all branches of length $`\mathrm{lg}(\eta )=\alpha `$
$$\eta :\alpha 2=\{0,1\}\text{for all}\alpha <\omega _1.$$
The order on $`𝐓`$ is defined naturally by extensions, i.e. if $`\eta ,\eta ^{}`$ T, then $`\eta \eta ^{}`$ if and only if $`\eta \eta ^{}`$ as maps. Hence $`\mathrm{lg}(\eta )\mathrm{lg}(\eta ^{})`$ and $`\eta ^{}\mathrm{lg}(\eta )=\eta `$, and we will require that $`p_\eta p_\eta ^{}`$ by the ordering of $`𝔉`$ as defined in Definition 2.3. If $`\eta {}_{}{}^{\omega _1}2`$, then $`p_{\eta \alpha }(\alpha \omega _1)`$ is linearly ordered and the triple
$`p_\eta ={\displaystyle \underset{\alpha \omega _1}{}}p_{\eta \alpha }`$ (2.2)
is well-defined. In details we have $`\mathrm{\Phi }_\eta :F_{0\eta }F_{1\eta }R`$ is defined by continuity such that $`\mathrm{\Phi }_\eta =\underset{\alpha \omega _1}{}\mathrm{\Phi }_{\eta \alpha }`$ and $`F_{i\eta }=\underset{\alpha \omega _1}{}F_{i\eta \alpha }`$. The bilinear forms $`\mathrm{\Phi }_\eta :F_{0\eta }F_{1\eta }R`$ will be our candidates for modules $`G`$ as in the Main Theorem 1.3. First we will show that $`𝔉\mathrm{}`$ and the arguments will be refined for Lemma 2.5.
###### Lemma 2.4
The partially ordered set $`𝔉`$ is non-empty.
Proof. Choose $`F_0=F_1=\underset{n\omega }{}e_nR`$ and extend
$$\mathrm{\Phi }(e_i,e_j)=\delta _{ij}=\{\begin{array}{ccc}\hfill 0& \hfill if& \hfill ij\\ \hfill 1& \hfill if& \hfill i=j\end{array}$$
linearly to get a bilinear map $`\mathrm{\Phi }:F_0F_1R.`$ The map $`\mathrm{\Phi }`$ satisfies Definition 2.1 for $`𝔉_0=𝔉_1=\mathrm{}`$. Next we want find $`𝔉_0=\{\phi _0\}`$ and $`𝔉_1=\{\phi _1\}`$. If $`\phi _i(i2)`$ satisfy Definition 2.1 $`(iv)`$ and $`\mathrm{Dom}\phi _i=F_i`$ then obviously
$$(\mathrm{\Phi },𝔉_0,𝔉_1)𝔉.$$
We will work for $`\phi =\phi _0:F_0R`$ and enumerate $`F_1\{0\}=\{x_i:i\omega \}`$. If $`\phi _i=\mathrm{\Phi }(,x_i)`$ then $`\phi _i0`$ by $`(i)`$ and for $`(iv)`$ we must show that
$`\mathrm{ker}\phi \mathrm{ker}\phi _i\text{for all}i\omega `$ (2.3)
Write $`F_0=\underset{i\omega }{}L_i`$ such that each $`L_i=e_iRe_i^{}R`$ is free of rank 2 and let $`L_i^{}=\underset{ji}{}L_j`$ be a complement of $`L_i`$ . If $`L_i\phi _i=0`$ then we use $`\phi _i0`$ to find some $`0y_{}L_i^{}`$ such that $`y\phi _i0.`$ Choose a new complement of $`L_i^{}`$ and rename it $`L_i=(e_i+y)Re_i^{}R`$.
Hence $`L_i\phi _i0`$ and there is a pure element $`y_iL_i`$ with $`y_i\phi _i0`$. We found an independent family $`\{y_i:i\omega \}`$ with $`F_0=\underset{i\omega }{}y_iRC`$ for some $`0CF_0`$. Choose $`\phi \mathrm{Hom}(F_0,R)`$ such that $`y_i\phi =0`$ for all $`i\omega `$ and $`\phi C0`$. Hence
$$y_i\mathrm{ker}\phi \mathrm{ker}\phi _i\text{for all}i\omega $$
and (2.3) holds. Hence Definition 2.1 holds and $`𝔉_1`$ can be chosen dually. $`\mathrm{}`$
The crucial step in proving the next result is again verification of Definition 2.1 $`(iv)`$, this time for $`\mathrm{\Phi }^{}(,x)`$. The proof is similar to the last one, hence we can be less explicit and just refine the old arguments.
###### Lemma 2.5
Let $`p=(\mathrm{\Phi },𝔉_0,𝔉_1)𝔉`$ be with $`\mathrm{Dom}\mathrm{\Phi }=F_0F_1`$ and $`\phi L^{}`$ for some pure submodule $`L`$ of finite rank in $`F_0`$. Then we find $`pp^{}=(\mathrm{\Phi }^{},𝔉_0,𝔉_1)𝔉`$ with $`\mathrm{Dom}\mathrm{\Phi }^{}=F_0F_1^{},F_1^{}=F_1xR`$ and $`\phi \mathrm{\Phi }^{}(,x).`$
Proof. First we want to extend $`\phi `$ to $`\phi ^{}:F_0R`$ such that
$`\mathrm{ker}E\mathrm{ker}\phi ^{}\text{ for all finite }E𝔉_0.`$ (2.4)
Enumerate all $`\mathrm{ker}E`$ for $`E𝔉_0`$ finite by $`\{K_i:0i\omega \}`$. Hence $`K_i=\underset{\phi E_i}{}\mathrm{ker}\phi `$ is a pure submodule of $`F_0`$ with $`F_0/K_i`$ free of rank $`|E_i|`$. As in the proof of (2.5) we can choose inductively $`\underset{i\omega }{}L_i=F_0`$ and $`0y_i_{}L_iK_i`$ for $`i>0`$ and $`L_0=L.`$ For some $`CF_0`$ we have $`F_0=\underset{i\omega }{}y_iRLC`$. Now we extend $`\phi L^{}`$ such that $`y_i\phi ^{}=1`$ and $`\phi C=0.`$ Hence
$$y_iK_i\mathrm{ker}\phi ^{}\text{ for all }0i\omega $$
and (2.4) holds. Finally we extend $`\mathrm{\Phi }`$ to $`\mathrm{\Phi }^{}:F_0(F_1xR)R`$ by taking $`\mathrm{\Phi }^{}(,x)=\phi ^{}`$. Condition (2.4) carries over to
$$\mathrm{ker}E\mathrm{ker}\mathrm{\Phi }(,x)\text{ for all finite }E𝔉_0.$$
From this it is immediate that $`\mathrm{ker}E\mathrm{ker}\mathrm{\Phi }(,y)`$ for all $`0yF_1^{}`$ and finite sets $`E𝔉_0`$. Hence Definition 2.1$`(iv)`$ holds, $`p(\mathrm{\Phi }^{},𝔉_0,𝔉_1)𝔉`$ and $`\phi \mathrm{\Phi }(,x)`$. $`\mathrm{}`$
Next we move elements from $`\mathrm{Hom}(F_\phi ,R)`$ for some $`\phi 𝔉_i`$ to $`𝔉_i`$.
###### Lemma 2.6
Let $`p=(\mathrm{\Phi },𝔉_0,𝔉_1)𝔉`$ be with $`\mathrm{Dom}\mathrm{\Phi }=F_0F_1`$. If $`\psi \mathrm{Hom}(F_\phi ,R)`$ for some $`\phi 𝔉_0`$ such that
$$\mathrm{ker}\psi \mathrm{ker}E\mathrm{ker}\mathrm{\Phi }(,x)\text{for all finite}E𝔉_0\text{and}xF_1,$$
then
$$pp^{}=(\mathrm{\Phi },𝔉_0^{},𝔉_1)𝔉\text{ for }𝔉_0^{}=𝔉_0\{\psi \}$$
Proof. We only have to check Definition 2.1$`(iv)`$ for finite subsets of $`𝔉_0^{}=𝔉_0\{\psi \}.`$ But this follows by hypothesis on $`\psi `$.
The proof of the following observation is obvious.
###### Observation 2.7
If $`p_n=(\mathrm{\Phi }_n,𝔉_{0n},𝔉_{1n})(n\omega )`$ is an ascending chain of elements $`p_n𝔉`$ and elements in $`𝔉_i`$ are unions of extensions in $`𝔉_{in}𝔉_{in+1}(n\omega )`$, then $`p=(\underset{n\omega }{}\mathrm{\Phi }_n,𝔉_0,𝔉_1)𝔉.`$
###### Definition 2.8
If $`(\mathrm{\Phi },𝔉_0,𝔉_1)𝔉`$ with $`\mathrm{Dom}\mathrm{\Phi }=F_0F_1`$ then $`\phi F_0^{}`$ is essential for $`\mathrm{\Phi }`$ if for any finite rank summand $`L`$ of $`F_0`$ and any finite subset $`E`$ of $`F_1`$ there is $`gF_0L`$ with $`g\phi 0`$ and $`\mathrm{\Phi }(g,e)=0`$ for all $`eE`$.
The notion ‘$`\phi F_1^{}`$ is essential for $`\mathrm{\Phi }`$ ’ is dual.
If $`g\phi =0=\mathrm{\Phi }(g,e)`$ for $`eE`$ and some $`E𝔉_1`$ then $`g\mathrm{\Phi }(,e):eEF_0^{}`$ by induction on $`|E|`$. Hence $`\phi F_0^{}`$ is essential for $`\mathrm{\Phi }`$ is equivalent to say that $`\phi `$ is not in $`\mathrm{\Phi }(,F_1)`$ modulo summands of finite rank. This leads to the following
###### Observation 2.9
Let $`(\mathrm{\Phi },𝔉_0,𝔉_1)𝔉`$ with $`\mathrm{Dom}\mathrm{\Phi }=F_0F_1`$. If $`EF_1`$ is a finite subset and $`\phi F_0^{}`$ inessential for $`\mathrm{\Phi }`$ with $`\mathrm{ker}\phi \mathrm{ker}E`$ i.e.
$$(\mathrm{\Phi }(x,e)=0\text{ for all }eE)x\phi =0,$$
then there is $`e_0EF_1`$ such that $`\phi =\mathrm{\Phi }(,e_0)`$.
Proof. By induction on $`|E|`$.
###### First Killing-Lemma 2.10
Suppose $`\phi F_0^{}`$ is essential for $`\mathrm{\Phi }`$ with $`\mathrm{Dom}\mathrm{\Phi }=F_0F_1`$ and $`p=(\mathrm{\Phi },𝔉_0,𝔉_1)𝔉`$. Then we find $`pp^{}𝔉`$ with $`p^{}=(\mathrm{\Phi }^{},𝔉_0^{},𝔉_1^{}),\mathrm{Dom}\mathrm{\Phi }^{}=F_0^{}F_1`$ and $`F_0^{}=F_0,y_{}\widehat{F}_0`$ for some $`y\widehat{F}_0`$ such that $`\phi `$ does not extend to $`\phi ^{}:F_0^{}R`$.
Remark A dual lemma holds for $`\phi F_1^{}`$.
Proof. Let $`F_0=\underset{i\omega }{}e_iR`$ and $`F_1=\underset{n\omega }{}f_nR`$. First we apply that $`\phi `$ is essential. It is easy to find inductively elements $`g_nF_0F^n`$ with $`F^n=e_i,g_i:i<n_{}F_0`$ such that the following holds
1. $`\mathrm{\Phi }(g_n,f_i)=0`$ for all $`i<n.`$
2. $`\underset{i<n}{}g_iR`$ is a direct summand - also $`\underset{i\omega }{}g_iR`$ is a summand of $`G`$.
3. $`g_n\phi 0`$ for all $`n\omega `$.
Decompose $`\omega `$ into a disjoint union of infinite subsets $`S_i(i\omega )`$ and let $`\{E_i:0i\omega \}`$ be an enumeration of all finite subsets of $`𝔉_1`$ and write $`K_i=\mathrm{ker}E_i`$ for all $`i>0`$. In order to get $`K_i\mathrm{ker}\mathrm{\Phi }(x,)`$ for all $`i>0,xF_0^{}F`$, we choose for each $`nS_i`$ an element $`k_nK_i`$ such that $`\mathrm{\Phi }(g_n,k_n)0`$. This is possible as $`F_1/K_i`$ is free of finite rank $`|E_i|`$, hence $`K_i`$ is ‘quite large’.
Now we use $`g_n\phi 0(n\omega )`$ and apply Observation 1.1 (for $`S_0`$) and choose a sequence $`ϵ_n\{0,1\}(n\omega )`$ and suitable $`q_nS`$ as in (1.1) such that
$`r={\displaystyle \underset{n\omega }{}}(g_n\phi )q_nϵ_n\widehat{R}R`$ (2.5)
$`\text{If }nS_i,\text{ then }{\displaystyle \underset{j=0}{\overset{n}{}}}\mathrm{\Phi }(g_j,k_i)q_jϵ_j0modq_{n+1}`$ (2.6)
If $`s\omega `$ and $`ϵ_s=1`$, then let
$$y_s=\underset{sn\omega }{}(q_s)^1q_ng_nϵ_n\widehat{F}_0$$
and consider the $`R`$-module $`F_0^{}=F_0,y_0_{}\widehat{F}_0`$, which can be generated by
$$F_0^{}=F_0,y_sR:s\omega ,ϵ_s=1.$$
Note that $`F_0^{}`$ is a countable $`R`$-module. It is easy to see that $`F_0^{}`$ is free - either apply Pontryagin’s theorem (Fuchs \[9, p. 93\]) or determine a free basis. The bilinear form $`\mathrm{\Phi }:F_0F_1R`$ extends uniquely to $`\mathrm{\Phi }^{}:F_0^{}F_1\widehat{R}`$ by continuity and density. We want to show that
$`(\mathrm{\Phi }^{},𝔉_0,𝔉_1)𝔉.`$ (2.7)
First we claim that Im $`\mathrm{\Phi }^{}R`$. By (i) we have
$$\mathrm{\Phi }^{}(y_0,f_j)=\mathrm{\Phi }^{}(\underset{n\omega }{}ϵ_nq_ng_n,f_j)=\underset{n\omega }{}\mathrm{\Phi }(g_n,f_j)q_nϵ_n=\underset{n<j}{}\mathrm{\Phi }(g_n,f_j)q_nϵ_nR$$
and $`\mathrm{Im}\mathrm{\Phi }^{}R`$ follows.
We also must check $`(iv)`$ from Definition 2.1 for the new elements $`yF_0^{}F_0`$. It is enough to consider
$$\mathrm{ker}E\mathrm{ker}\mathrm{\Phi }(y_s,)\text{ for all }s\omega \text{ with }ϵ_s=1.$$
By definitions and enumerations this is equivalent to say that
$$K_i\mathrm{ker}\mathrm{\Phi }(y_s,)\text{for each}i>0.$$
If $`nS_i`$ and $`n>s`$ by (2.6) we have that
$$\mathrm{\Phi }(y_s,k_i)=\underset{j\omega }{}\mathrm{\Phi }(g_j,k_i)q_jϵ_j\underset{j=0}{\overset{n}{}}\mathrm{\Phi }(g_j,k_i)q_jϵ_j0modq_{n+1}$$
Hence $`k_i\mathrm{ker}\mathrm{\Phi }(y_s,)`$ but $`k_iK_i`$ and (2.7) follows. Finally we must show that $`\phi F_0^{}`$ does not extend to $`(F_0^{})^{}`$. By continuity $`\phi :F_0R`$ extends uniquely to $`\phi ^{}:F_0^{}\widehat{R}`$. However
$$y_0\phi =(\underset{n\omega }{}g_nq_nϵ_n)\phi ^{}=\underset{n\omega }{}(g_n\phi )q_nϵ_n=r\widehat{R}R$$
by (2.5), hence $`\phi `$ does not extend to $`F_0^{}R`$. $`\mathrm{}`$
###### Second Killing-Lemma 2.11
Let $`p=(\mathrm{\Phi },𝔉_0,𝔉_1)𝔉`$ and $`\eta :F_\phi F_\phi `$ be some monomorphism with $`F_\phi =x_0RF_\phi \eta `$ for some $`\phi 𝔉_0`$ with $`x_0F_\phi =\mathrm{ker}\phi `$. Then there is $`pp^{}=(\mathrm{\Phi }^{},𝔉_0^{},𝔉_1^{})𝔉`$ with $`\mathrm{Dom}\mathrm{\Phi }^{}=F_0^{}F_1^{},\phi \phi ^{}𝔉_0^{}`$ such that $`\eta `$ does not extend to a monomorphism
$$\eta ^{\prime \prime }:F_\phi ^{}F_\phi ^{}_{}F_0^{\prime \prime }\text{ with }F_\phi ^{}^{\prime \prime }=x_0RF_\phi ^{}^{\prime \prime }\eta ^{\prime \prime }$$
where $`p^{}p^{\prime \prime }:=(\mathrm{\Phi }^{\prime \prime },𝔉_0^{\prime \prime },𝔉_1^{\prime \prime })`$ and $`\mathrm{Dom}\mathrm{\Phi }^{\prime \prime }=F{}_{}{}^{\prime \prime }{}_{0}{}^{}F_1^{\prime \prime }`$.
Proof. In order to satisfy Definition 2.1 $`(iv)`$ for the new $`p^{}𝔉`$ we argue similar to the First Killing Lemma 2.10. Let $`\omega =\underset{i\omega }{}S_i`$ be a decomposition into infinite subsets $`S_i`$ and $`\{K_i:0i\omega \}`$ be all kernels $`K_i=\mathrm{ker}E_i`$ for an enumeration of finite subsets $`E_i`$ of $`𝔉_1`$. The set $`S_0`$ will be used for killing $`\eta `$ and the $`S_i(i>0)`$ are in charge of $`K_i`$ and $`(iv)`$ above. Extending $`\mathrm{\Phi }\mathrm{\Phi }^{}`$ we must ensure $`\mathrm{Im}\mathrm{\Phi }^{}R`$. Hence we construct an increasing sequence $`s_n\omega (n\omega )`$ and pose more conditions on $`s_n`$ later on. If $`F^{}=\mathrm{Dom}\phi `$ then $`F^{}=x_0RF^{}\eta `$ and if $`x_i=x_0\eta ^i`$ we get
$`F^{}={\displaystyle \underset{i<n}{}}x_iRF^{}\eta ^n\text{ for all }n\omega .`$ (2.8)
Note that $`\underset{i<n}{}x_iR`$ is pure of finite rank, hence $`F_0=\underset{i<n}{}x_iRC_n`$ for some $`C_nF_0`$. Now we construct $`T_n=\underset{s_ni<s_{n+1}}{}e_iR`$ from $`F_0=\underset{i\omega }{}e_iR`$ and refine an argument from Göbel, Shelah . Obviously $`F_0=\underset{n\omega }{}T_n`$ and let $`e^n\mathrm{Hom}(F_0,R)(n\omega )`$ be defined by
$$e_ie^n=\delta _{i,n}=\{\begin{array}{ccc}\hfill 0& \hfill if& \hfill in\\ \hfill 1& \hfill if& \hfill i=n\end{array}$$
for all $`i\omega `$. If $`i<n`$ then $`\pi _i:F_0R`$ denotes the projection modulo $`\underset{ij<n}{}x_jRC_n`$, moreover let $`F_1=\underset{j\omega }{}f_jR`$ be as before. We now seek for elements $`w_nF_0(n\omega )`$ subject to the following four conditions
1. $`0w_nT_n`$ and $`w_n\eta T_n`$.
2. $`\mathrm{\Phi }(w_n,f_n)=w_ne^k=w_n\eta e^k=0`$ for all $`k<s_n`$
3. If $`nS_0,`$ then let $`\pi _n^{}`$ be the projection $`\pi _i`$ with $`i`$ maximal such that
$$w_n\pi _i=0w_n\eta \pi _i,s_ni<s_{n+1}$$
4. If $`nS_i,i>0`$, there is $`y_nK_i`$ such that $`\mathrm{\Phi }(w_n,y_n)0.`$
Suppose $`s_0,\mathrm{},s_n,w_0\mathrm{},w_{n1}`$ are constructed and we want to choose $`w_n,s_{n+1}`$. Then pick $`s_n<s_{n+1}`$ such that
$$\{x_{s_n},\mathrm{},x_{4s_n+1}\}\underset{is_{n+1}}{}e_iR.$$
We want to choose $`w_n=\underset{i=s_n}{\overset{4s_n}{}}x_ia_i^n`$ for some $`a_i^nR`$. If $`w_ne^k=0`$ for $`k<s_n`$, then $`w_nT_n`$. Moreover
$$w_n\eta =\underset{i=s_n}{\overset{4s_n}{}}x_ia_i^n\eta =\underset{i=s_n}{\overset{4s_n}{}}x_{i+1}a_i^n=\underset{i=s_n+1}{\overset{4s_n+1}{}}x_ia_{i1}^n$$
and $`x_n\eta T`$ because $`x_{4s_n+1}\underset{is_{n+1}}{}e_iR`$ as well.
Hence $`(i)`$ follows provided $`w_n0`$ is generated by those $`x_i^{}s`$. The conditions $`(ii)`$ can be viewed as a system of $`3s_n`$ homogeneous linear equations in $`4s_n+1s_n=3s_n+1`$ unknowns $`a_i^nR`$. We find a non-trivial solution $`w_n0`$ by linear algebra. Hence $`(i)`$ and $`(ii)`$ hold. Condition $`(iii)`$ follows by hypothesis on $`𝔉_1`$ for $`E_i`$ and $`K_i=\mathrm{ker}E_i`$. Condition $`(iii)`$ finally follows by the action of $`\eta `$ on $`w_n0`$ and the maximality of $`i`$ with $`s_n<i4s_n`$ and $`w_n\pi _i=0`$ for $`\pi _n^{}=\pi _i`$. Hence $`(i),\mathrm{},(iv)`$ follow.
As in the proof of the First Killing Lemma 2.10 inductively we choose a strictly increasing sequence $`m_j\omega (j\omega )`$. If $`m_j`$ is defined up to $`jn`$ we must choose $`m_{n+1}`$ large enough such that
$$\underset{jn}{}q_{m_j}\mathrm{\Phi }(w_j,y_n)0modq_{m_{n+1}}.$$
This needs inductively the hypothesis that $`\underset{jn}{}q_{m_j}\mathrm{\Phi }(w_j,y_n)0`$. If $`n+1S_i`$, then $`\mathrm{\Phi }(w_{n+1},y_{n+1})0`$ by $`(iv)`$ and we may assume
$$\underset{jn+1}{}p_j^m\mathrm{\Phi }(w_j,y_{n+1})0$$
hence the inductive hypothesis follows and we can proceed. By Observation 1.1 and $`(iii)`$ we also find $`ϵ_j\{0,1\}(jS_0)`$ such that
$`{\displaystyle \underset{jS_0}{}}(w_j\eta )\pi _j^{}q_{m_j}ϵ_j\widehat{R}R.`$ (2.9)
Now we are ready to extend $`\mathrm{\Phi }`$. Choose new elements
$$z_k=\underset{jk}{}w_jq_{m_j}(q_{m_k})^1\widehat{F}_0$$
for all $`k\omega `$. Hence the submodule $`F_0^{}=F_0,z_0_{}\widehat{F}_0`$ purely generated by adding $`z=z_0`$ is generated by
$$F_0^{}=F_0,z_k:k\omega .$$
Again we see that $`F_0^{}`$ is a countable, free $`R`$-module. The map $`\mathrm{\Phi }:F_0F_1R`$ by continuity extends uniquely to
$$\mathrm{\Phi }^{}:F_0^{}F_1\widehat{R}.$$
Recall that $`F_0^{}/F_0`$ is S-divisible, hence $`F_0`$ is S-dense in $`F_0^{}`$ in the S-adic topology. First we must show that $`\mathrm{Im}\mathrm{\Phi }^{}R`$. We apply $`(ii)`$ and continuity to see that
$$\mathrm{\Phi }^{}(z,f_k)=\mathrm{\Phi }^{}(\underset{n\omega }{}q_{m_n}w_n,f_k)=\underset{n\omega }{}q_{m_n}\mathrm{\Phi }(w_n,f_k)=\underset{nk}{}q_{m_n}\mathrm{\Phi }(w_n,f_k)R,$$
hence $`\mathrm{\Phi }^{}:F_0^{}F_1R`$.
We also must show Definition 2.1 $`(iv)`$ for the new elements $`z_tF_0^{}`$. We have $`K_i=\mathrm{ker}E_i`$ and $`S_i`$ is unbounded. Hence we find $`nS_i,n>t`$ and $`y_nK_i`$ such that
$$\mathrm{\Phi }(w_n,y_n)0.$$
We apply $`\mathrm{\Phi }^{}`$ to $`(z_t,y_n)`$ and get
$$\mathrm{\Phi }^{}(z_t,y_n)\underset{tjn}{}q_{m_j}(q_{m_t}^1\mathrm{\Phi }(w_j,y_n)0modq_{m_{n+1}}q_{n_t}^1$$
hence $`\mathrm{\Phi }(z_t,y_n)0`$ and $`y_nK_i`$. This is equivalent to say that $`\mathrm{ker}E_i\mathrm{ker}\mathrm{\Phi }(z_t,)`$ and Definition 2.1 $`(iv)`$ follows. Hence $`(\mathrm{\Phi }^{},𝔉_0,𝔉_1)𝔉`$ with $`\mathrm{Dom}\mathrm{\Phi }^{}=F_0^{}F_1.`$
Next we extend $`\mathrm{\Phi }^{}`$ under the name $`\mathrm{\Phi }^{}`$. Let $`F_1^{}=F_1fR`$ be a free rank-1 extension. We want $`\mathrm{\Phi }^{}:F_0^{}F_1^{}R`$ and must define $`\mathrm{\Phi }^{}(,f):F_0^{}R`$. Put $`\mathrm{\Phi }^{}(,f)T_n=ϵ_n\pi _n^{}T_n`$ if $`nS_0`$ and $`\mathrm{\Phi }^{}(,f)T_n=0`$ otherwise. By linear extension and $`F_0=\underset{n\omega }{}T_n`$ the map $`\mathrm{\Phi }^{}(,f):F_0R`$ is well-defined. It extends further by continuity to
$$\mathrm{\Phi }^{}:F_0^{}F_1^{}\widehat{R}.$$
Again we must show that $`\mathrm{Im}\mathrm{\Phi }^{}R`$. Note that $`\mathrm{\Phi }^{}(w_n,f)=w_nϵ_n\pi _n^{}=0`$ for $`nS_0`$ from (iii), and $`\mathrm{\Phi }^{}(w_n,f)=0`$ for $`n\omega S_0`$ by the above, hence
$$\mathrm{\Phi }^{}(z,f)=\underset{n\omega }{}q_{m_n}\mathrm{\Phi }^{}(w_n,f)=0$$
and $`\mathrm{Im}\mathrm{\Phi }^{}R`$ follows.
We also must check condition $`(iv)`$ of Definition 2.1 for the new element $`fF_1^{}`$. Recall $`K_i(i\omega )`$ is a list of all kernels $`\mathrm{ker}E_i`$ for finite sets in $`𝔉_0`$. Also enumerate all $`T_n`$’s with $`nS_0`$ and $`r_n=1`$ as, say $`T^i=T_{n_i}(i\omega )`$. Choosing the ranks of the $`T^i^{}s`$ large enough we can find (as before) $`y_i^{}K_iT^i`$ with $`y_i^{}\pi _{n_i}^{}0`$. Then
$$y_i^{}\mathrm{ker}E_i\text{ but }\mathrm{\Phi }^{}(y_i^{},f)=y_i^{}\pi _{n_i}^{}0,$$
hence $`y_i^{}\mathrm{ker}\mathrm{\Phi }(,f)`$ as desired for $`(iv)`$ above.
Finally we must get rid of $`\eta `$ by showing that there is no extension $`\eta ^{\prime \prime }\eta `$ as stated in the Lemma. Otherwise we have $`\mathrm{\Phi }^{}\mathrm{\Phi }^{\prime \prime }:F_0^{\prime \prime }RR`$ and
$$x_0RF^{\prime \prime }\eta ^{\prime \prime }=F^{\prime \prime }$$
for $`\mathrm{Dom}\eta ^{\prime \prime }=F^{\prime \prime }_{}F_0^{\prime \prime }`$ with $`z\mathrm{Dom}\eta ^{}F^{\prime \prime }`$. Hence $`r=\mathrm{\Phi }^{\prime \prime }(z\eta ^{\prime \prime },f)R`$. On the other hand
$$r=\mathrm{\Phi }^{\prime \prime }(z\eta ^{\prime \prime },f)=\mathrm{\Phi }^{\prime \prime }(\underset{j\omega }{}q_{m_j}(w_j\eta ^{}),f)=\underset{j\omega }{}q_{m_j}\mathrm{\Phi }^{}(w_j\eta ^{},f)=\underset{jS_0}{}q_{m_j}ϵ_j(w_j\eta \pi _j^{})\widehat{R}R$$
is a contradiction. The lemma follows. $`\mathrm{}`$
## 3 Construction of reflexive modules assuming CH
Let $`\{S_0,\mathrm{},S_5\}`$ be a decomposition of the set of all limit ordinals in $`\omega _1`$ into stationary sets. Using CH we can enumerate sets of cardinality $`2^\mathrm{}_0`$ by any of these stationary sets of size $`\mathrm{}_1`$.
Let
$`𝐓={}_{}{}^{\omega _{1>}}2=\{\eta _\alpha :\alpha \omega _1\}`$ (3.1)
be an enumeration of the tree with countable branches such that
$`\eta _\alpha \eta _\beta \alpha <\beta `$ (3.2)
We are working in the ‘universe’ $`\omega _1`$ and let
$`\mathrm{Map}(\omega _1,R)=\{\phi _i:iS_0\}=\{\phi _i:iS_1\}=\{\phi _i:iS_2\}=\{\phi _i:iS_3\}`$ (3.3)
be four lists of all partial functions $`\phi :\delta R`$ for $`\delta \omega _1`$ any limit ordinal with $`\mathrm{}_1`$ repetitions $`\phi _i`$ for each $`\phi `$. Similarly let
$$\mathrm{Map}(\omega _1,\omega _1)=\{\mu _i:iS_4\}=\{\mu _i:iS_5\}$$
be two lists of all partial maps $`\mu :XX\omega _1`$ with countable domains $`X`$ given by the prediction Lemma 1.2 for $`E\{S_4,S_5\}`$.
Inductively we want to construct an order preserving continuous map
$$p:{}_{}{}^{\omega _1>}2=𝐓𝔉,(\eta p_\eta )$$
subject to certain conditions $`(a),\mathrm{},(d)`$ dictated by the proof of the main theorem and stated below. Some preliminary words are in order. Using the enumeration (3.1) we may write
$`p_{\eta _\alpha }=p_\alpha =(\mathrm{\Phi }_\alpha ,𝔉_{0\alpha },𝔉_{1\alpha })\text{ and }\mathrm{Dom}\mathrm{\Phi }_\alpha =F_{0\alpha }F_{1\alpha }`$ (3.4)
If $`p`$ is order preserving then $`(\eta _\alpha \eta _\beta p_\alpha p_\beta )`$ and continuity applies. If $`\eta {}_{}{}^{\alpha }2`$ has length $`\mathrm{lg}(\eta )=\alpha `$ and $`\alpha <\omega _1`$ is a limit ordinal, then $`\underset{\beta <\alpha }{}\eta \beta =\eta `$ and $`\eta \beta \eta `$, hence
$`\eta {}_{}{}^{\alpha }2p_\eta ={\displaystyle \underset{\beta <\alpha }{}}p_{\eta \beta }`$ (3.5)
If all $`p_{\eta \beta }𝔉`$ then $`p_\eta 𝔉`$ by Observation 2.7. Hence we have no problem in defining $`p`$ at limit stages by continuity. It remains to consider inductive steps for $`p`$. If $`\eta ,\eta ^{}{}_{}{}^{\omega _1>}2`$ then $`\gamma =\mathrm{br}(\eta ,\eta ^{})`$ denotes the branching point of $`\eta ,\eta ^{}`$, this is to say that
1. $`\gamma <\mathrm{min}\{\mathrm{lg}(\eta ),\mathrm{lg}(\eta ^{})\}`$
2. $`\eta \gamma =\eta ^{}\gamma `$
3. $`\eta (\gamma )\eta ^{}(\gamma ).`$
If $`\eta `$ is a branch of length $`\alpha `$ and $`i\{0,1\}`$ then $`\eta ^{}=\eta \{i\}`$ is a branch of length $`\alpha +1`$ with $`\eta ^{}\alpha =\eta `$ and $`\eta ^{}(\alpha )=i.`$ Now we continue defining $`p:{}_{}{}^{\omega _1>}2𝔉`$. Generally we require
(a) If $`\phi 𝔉_{i\alpha }`$ then $`\mathrm{Dom}\phi =F_{i\beta }`$ for some $`\eta _\beta \eta _\alpha .`$
(b) The set $`\{\mathrm{Dom}\phi :\phi 𝔉_{i\alpha }\}`$ is a well ordered set of pure submodule of $`F_{i\alpha }`$.
(c) If $`\gamma =\mathrm{br}(\eta _\alpha ,\eta _\beta )`$ and $`\eta _\alpha \gamma =\eta _\beta \gamma =\eta _ϵ`$ for some $`ϵ<\omega _1`$, and $`yF_{1\beta }F_{1ϵ}`$ then $`\mathrm{\Phi }^\beta (y,)F_{0ϵ}𝔉_{0\alpha }.`$ Dually, if $`yF_{0\beta }F_{1\beta }`$ then $`\mathrm{\Phi }^\beta (y,)F_{1ϵ}𝔉_{1\alpha }.`$
(d) If $`\nu {}_{}{}^{\gamma }2`$ is a branch of length $`\gamma ,i\{0,1\}`$ and $`\eta =\nu \{i\}`$ then we want to define $`p_\eta `$ depending on $`\gamma `$.
1. For $`\nu =\mathrm{}`$ choose $`p_{\{i\}}`$ by Lemma 2.5.
2. If $`\gamma S_0`$ then we want to enlarge $`\mathrm{\Phi }_\nu `$ to make the evaluation map injective: If $`\phi _\gamma F_\phi `$ for some $`\phi 𝔉_{0\nu }`$ is a partial $`R`$-homomorphism $`F_\phi R`$ with $`\mathrm{Dom}(\phi _\gamma F_\phi )`$ a pure $`R`$-submodule of finite rank of $`F_\phi `$, then we apply Lemma 2.6 such that $`p_\nu p_\eta `$ with $`F_{1\eta }=F_{1\nu }x_\eta R,F_{0\eta }=F_{0\nu }`$, $`\mathrm{\Phi }_\nu \mathrm{\Phi }_\eta `$ and $`\phi _\gamma \mathrm{\Phi }_\eta (,x_\eta ).`$ If $`\phi _\gamma `$ does not satisfy the requirements, then we choose any $`\mathrm{\Phi }_\nu \mathrm{\Phi }_\eta `$.
3. If $`\gamma S_1`$ then we argue as in $`(ii)`$ but dually. A dual version of Lemma 2.6 provides $`\mathrm{\Phi }_\nu \mathrm{\Phi }_\eta `$ and $`\phi _\gamma \mathrm{\Phi }_\eta (x_\eta ,)`$ if $`\phi _\gamma `$ meets the requirements.
4. If $`\gamma S_2`$ then we want to kill bad dual maps to make the evaluation map surjective. If $`\phi _\gamma F_{0\nu }`$ is an $`R`$-homomorphism $`F_{0\nu }R`$ which is essential for $`\mathrm{\Phi }_\nu `$ then we apply the First Killing Lemma 2.10 to find $`p_\gamma p_\eta `$ such that $`F_{0\eta }=F_{0\nu },y_{}\widehat{F}_{\nu 0}`$ for some $`y\widehat{F}_{\nu 0}`$ and $`\phi _\gamma F_{0\nu }`$ does not extend to $`F_{0\eta }R.`$
5. If $`\gamma S_3`$ and $`\phi _\gamma F_{1\nu }`$ is an essential $`R`$-homomorphism for $`\mathrm{\Phi }_\nu `$ then we argue as in $`(iv)`$ but dually.
6. If $`\gamma S_4`$ then we want to get rid of potential monomorphisms $`\eta `$ of the final module $`G`$ with $`G\eta xR=G.`$ If $`\phi 𝔉_{0\nu },\mu =\phi _\gamma F_\phi `$ and $`\mu :F_\phi F_\phi `$ is an $`R`$-monomorphism such that $`F_\phi =x_\phi RF_\phi \mu `$, then we apply the Second Killing Lemma 2.11 to find $`p_\nu p_\eta `$ such that $`\mu `$ does not extend to a monomorphism $`\mu ^{}`$ of an extension $`F_0^{}`$ of $`F_{0\eta }`$ with
$$p_\eta p^{}=(\mathrm{\Phi }^{},𝔉_0^{},𝔉_1^{}),\mathrm{Dom}\mathrm{\Phi }^{}=F_0^{}F_1^{}\text{ and }F_0^{}=x_\phi RF_0^{}\mu ^{}.$$
7. If $`\gamma S_5`$, then we argue dually for some partial monomorphism $`\mu `$ with domain and range some $`F_\phi _{}F_{1\nu }`$.
8. If $`\gamma \omega _1`$ is not a limit ordinal, then we are free to choose any ‘trivial’ extension $`p_\nu p_\eta `$.
## 4 Proof of the Main Theorem
Recall from Section 3 that we are given an order preserving continuous map
$$p:\text{ }\text{}𝔉.$$
By continuity we may extend $`p`$ to
$$p_\eta =(\mathrm{\Phi }_\eta ,𝔉_{0\eta },𝔉_{1\eta })\text{ for all }\eta {}_{}{}^{\omega _1}2.$$
It is immediate that
$`F_{i\eta }\text{ is }\mathrm{}_1\text{-free and }\mathrm{\Phi }_\eta :F_{0\eta }F_{1\eta }R\text{ is a bilinear form,}`$ (4.1)
where $`i\{0,1\},\eta {}_{}{}^{\omega _1}2.`$ Moreover, $`\mathrm{\Phi }_\eta `$ preserves purity and is not degenerated in the sense of Definition 2.1. $`\mathrm{\Phi }_\eta :F_{0\eta }F_{1\eta }R`$ is our candidate for a reflexive modules, expressed in an unusual way.
To see that $`\mathrm{\Phi }_\eta `$ preserves purity we must show that each pure element $`e_{}F_{0\eta }`$ induces $`\mathrm{\Phi }_\eta (e,)_{}F_{1\eta }^{}`$ (and dually). We may restrict to the first case. If $`e_{}F_{0\eta }`$ then $`e_{}F_{0\eta \alpha }`$ for any $`\alpha \omega _1`$ large enough, and $`\mathrm{\Phi }_{\eta \alpha }(e,)_{}F_{1\eta \alpha }^{}`$ by $`(\mathrm{\Phi }_{\eta \alpha },𝔉_{0\eta \alpha },𝔉_{1\eta \alpha })𝔉`$ and Definition 2.1 $`(ii)`$. It follows that $`\underset{\alpha \omega _1}{}\mathrm{\Phi }_{\eta \alpha }(e,)_{}F_{1\eta }^{}.`$
To see that $`\mathrm{\Phi }_\eta `$ is not degenerated we consider $`0eF_{0\eta }`$. Hence $`ee^{}R_{}F_{0\eta \alpha }`$ for a pure element $`e^{}`$ and any $`\alpha \omega _1`$ large enough. The partial homomorphism $`\phi :e^{}RR`$ defined by $`e^{}\phi =1`$ has a number $`iS_0`$ in the list and $`\phi =\phi _i`$ . By construction there is $`yF_{1\eta \gamma }`$ with
$$\phi \mathrm{\Phi }_{\eta \gamma }(,y).$$
Hence $`0e\phi =\mathrm{\Phi }_{\eta \gamma }(e,y)=\mathrm{\Phi }_\eta (e,y)`$ and $`\mathrm{\Phi }_\eta `$ is not degenerated. $`\mathrm{}`$
###### Definition 4.1
We will say that $`(\mathrm{\Phi },F_0,F_1)`$ with $`\mathrm{Dom}\mathrm{\Phi }=F_0F_1`$ is fully represented if $`\mathrm{\Phi }_\eta (F_{0\eta },)=F_{1\eta }^{}`$ and $`\mathrm{\Phi }_\eta (,F_{1\eta })=F_{0\eta }^{}`$.
We claim that
$`(\mathrm{\Phi }_\eta ,F_{0\eta },F_{1\eta })\text{ is fully represented for almost all }\eta {}_{}{}^{\omega _1}2`$ (4.2)
There are at most $`\eta W{}_{}{}^{\omega _1}2`$ exceptions with $`|W|<2^\mathrm{}_1`$.
Suppose for contradiction that $`|W|=2^\mathrm{}_1`$ and
$`\phi _\eta F_{1\eta }^{}\mathrm{\Phi }_\eta (F_{0\eta },)\text{for all}\eta W`$ (4.3)
By a pigeon hole argument there are $`\eta ,\eta ^{}W`$ with $`\mathrm{br}(\eta ,\eta ^{})=\alpha `$ and
(a) $`\phi _\eta F_{1\alpha }=\phi _\eta ^{}F_{1\alpha }`$
(b) $`\phi _\eta :F_{1\eta }R`$, and $`\phi _\eta ^{}:F_{1\eta ^{}}R`$ are not represented.
Let $`\psi =\phi _\eta F_{1\alpha }=\phi _\eta ^{}F_{1\alpha }:F_{1\alpha }R`$ and recall that there is some $`\gamma S_3`$ with $`\psi =\phi _\gamma `$ and $`(d)(v)`$ of the construction applies. Hence $`\psi `$ is inessential for $`\mathrm{\Phi }_\gamma `$. There is a finite set $`EF_{0\gamma }`$ such that
$$\mathrm{\Phi }_\gamma (e,x)=0\text{for all}eE,x\mathrm{Dom}\psi =F_{1\alpha }\text{ implies }x\psi =0.$$
By Observation 2.9 we have some $`eEF_{0\gamma }`$ such that $`\psi =\mathrm{\Phi }_\gamma (,e)F_{1\alpha }`$. The same argument applies for $`\phi _\eta ^{}`$ and there are $`\gamma ^{}S_3`$ and $`e^{}F_{0\gamma ^{}}`$ such that $`\gamma <\gamma ^{}`$ and $`\psi =\mathrm{\Phi }_\gamma ^{}(,e^{})F_{1\alpha }`$. Hence $`\psi =\mathrm{\Phi }_\gamma ^{}(,e^{})F_\alpha =\mathrm{\Phi }_\gamma (,e)F_\alpha `$.
Finally we apply $`(c)`$ of the construction to get $`\psi 𝔉_{0\gamma ^{}}`$. Now Definition 2.1 $`(iv)`$ applies and $`\mathrm{ker}\psi \mathrm{ker}\mathrm{\Phi }_\gamma ^{}(,e^{})`$ is a contradiction because $`\psi =\mathrm{\Phi }_\gamma ^{}(,e^{})`$. The claim (4.2) follows. $`\mathrm{}`$
Note that we did not use $`(vi)`$ so far. Hence without the Second Killing Lemma 2.11 we are able to derive reflexivity of the modules $`G_\eta `$, which we will do next.
We will use the following notations.
$$\text{Let }I^{}=\{\eta {}_{}{}^{\omega _1}2\text{ such that }p_\eta \text{ is not fully represented} and }I={}_{}{}^{\omega _1}2I^{}$$
From (4.2) we see that $`|I^{}|<2^\mathrm{}_1`$, hence $`|I|=2^\mathrm{}_1`$ . If $`\eta I`$ and $`i\{0,1\}`$ then we also fix the evaluation map
$$\sigma _{i\eta }:F_{i\eta }F_{i\eta }^{}$$
and claim that
$`\sigma _{i\eta }:F_{i\eta }F_{i\eta }^{}\text{ is injective for all }\eta I,i\{0,1\}.`$ (4.4)
Proof. We consider $`\sigma =\sigma _{0\eta }`$ and apply that $`\mathrm{\Phi }_\eta `$ is not degenerated. If $`0xF_{0\eta }`$ there is $`yF_{1\eta }`$ such that $`\mathrm{\Phi }_\eta (x,y)0`$. Hence $`\phi :=\mathrm{\Phi }_\eta (,y)F_{0\eta }^{}`$ and $`x\phi =\mathrm{\Phi }_\eta (x,y)0`$, thus $`x\sigma 0`$ and $`\sigma `$ is injective. The case $`\sigma _{1\eta }`$ is similar. $`\mathrm{}`$
Next we show that
$`\sigma _{i\eta }:F_{i\eta }F_{i\eta }^{}\text{ is surjective for all }\eta I,i\{0,1\}.`$ (4.5)
Proof. First note that
$`\mathrm{\Phi }_\eta ^{}:F_{0\eta }F_{1\eta }^{}(x\mathrm{\Phi }_\eta (x,))\text{ is bijective}`$ (4.6)
and
$`{}_{}{}^{}\mathrm{\Phi }_{\eta }^{}:F_{1\eta }F_{0\eta }^{}(y\mathrm{\Phi }_\eta (,y))\text{ is bijective}`$ (4.7)
because $`\mathrm{\Phi }_\eta `$ is not degenerated and $`\eta I`$. Hence we can identify $`F_{1\eta }`$ and $`F_{0\eta }^{}`$ by $`{}_{}{}^{}\mathrm{\Phi }_{\eta }^{}`$ and $`F_{0\eta }^{}=\mathrm{Im}(^{}\mathrm{\Phi }_\eta )=\mathrm{\Phi }_\eta (,F_{1\eta })`$. Moreover
$$F_{0\eta }^{}=(F_{0\eta }^{})^{}=(F_{1\eta })^{}=\mathrm{Im}\mathrm{\Phi }_\eta ^{}=\mathrm{\Phi }_\eta (F_{0\eta },)$$
and for any $`\phi F_{0\eta }^{}`$ we find $`fF_{0\eta }`$ with $`\phi =\mathrm{\Phi }_\eta (f,)`$. We consider the case $`\sigma =\sigma _{0\eta }`$, and get
$$\mathrm{\Phi }_\eta (,x)\sigma (f)=\mathrm{\Phi }_\eta (f,x)=\mathrm{\Phi }_\eta (,x)\mathrm{\Phi }_\eta (f,)$$
for all $`xF_{1\eta }`$ and $`\mathrm{\Phi }_\eta (,x)`$ runs through all of $`F_{0\eta }^{}`$. We derive $`\sigma (f)=\mathrm{\Phi }_\eta (f,)=\phi `$ and $`\sigma `$ is surjective. The case $`\sigma _{1\eta }`$ is similar. $`\mathrm{}`$
We have an immediate corollary from (4.4) and (4.5).
###### Corollary 4.2
If $`\eta I`$ and $`i\{0,1\}`$ then $`F_{i\eta }`$ is a reflexive $`R`$-module. Moreover $`I`$ is a subset of $`{}_{}{}^{\omega _1}2`$ of cardinality $`2^\mathrm{}_1`$.
For the Proof of the Main Theorem 1.3 we finally must show that
$`F_{i\eta }\cong ̸RF_{i\eta }\text{ for any }\eta I,i\{0,1\}.`$ (4.8)
We consider $`F_{i\eta }`$ with some monomorphism $`\xi :F_{i\eta }F_{i\eta }`$ such that $`F_{i\eta }=F_{i\eta }\xi xR`$ for any $`\eta I,i\{0,1\}`$. By a back and forth argument there is an $`\alpha <\omega _1`$ such that $`F_{i\eta \alpha }=F_{i\eta \alpha }\xi xR`$ and we take $`\psi =\xi F_{i\eta \alpha }`$ into consideration. There is some $`\gamma S_4`$ such that $`\psi =\phi _\gamma `$ and $`\phi _\gamma `$ is discarded by the construction. Hence $`\xi `$ does not exist. $`\mathrm{}`$
We would like to add a modification of our main result which can be shown using one more stationary subset $`S_6`$ after introducing inessential endomorphisms for our category of reflexive modules. Let $`\mathrm{Fin}(G)`$ be the ideal of all endomorphisms
$$\{\sigma \mathrm{End}(G):G\sigma \text{ has finite rank}\}$$
for some torsion-free $`R`$-module $`G`$. Then we can find a ‘Killing-Lemma’ in Dugas, Göbel , see also , which ‘takes care’ of all endomorphisms which are not in $`\mathrm{Fin}(F_{i\eta })`$ with $`i,\eta `$ as above. Hence we can strengthen our Main Theorem 1.3 and get with slight modification from known results the following
###### Corollary 4.3
(ZFC + CH) Let $`R`$ is a countable domain but not a field and $`A`$ be a countable $`R`$-algebra with free additive structure $`A_R`$. Then there is a family of $`2^\mathrm{}_1`$ pair-wise non-isomorphic reflexive $`R`$-modules $`G`$ of cardinality $`\mathrm{}_1`$ such that $`G\cong ̸RG`$ and $`\mathrm{End}(G)=A\mathrm{Fin}(G)`$ a split extension.
In particular $`\mathrm{End}(G)/\mathrm{Fin}(G)A`$ and if $`A`$ has only trivial idempotents like $`A=R`$ then $`G`$ in the Corollary 4.3 is reflexive, essentially indecomposable of size $`\mathrm{}_1`$ and does not decompose into $`GRG`$.
Rüdiger Göbel
Fachbereich 6, Mathematik und Informatik
Universität Essen, 45117 Essen, Germany
e–mail: R.Goebel@Uni-Essen.De
and
Saharon Shelah
Department of Mathematics
Hebrew University, Jerusalem, Israel
and Rutgers University, Newbrunswick, NJ, U.S.A
e-mail: Shelah@math.huji.ae.il
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# 1 Introduction
## 1 Introduction
In this paper we study the effects of nonuniversality of gaugino masses on dark matter in SUGRA models and in string and D brane models. Nonuniversal gaugino masses arise quite naturally in supergravity and string unified theories. Thus, in N=1 supergravity the kinetic energy and the mass terms for the gauge fields and the gauginos are given by
$`e^1=`$ $`{\displaystyle \frac{1}{4}}\mathrm{}\left[f_{\alpha \beta }F_{\mu \nu }^\alpha F^{\beta \mu \nu }\right]+{\displaystyle \frac{1}{4}}i\mathrm{}\left[f_{\alpha \beta }F_{\mu \nu }^\alpha \stackrel{~}{F}^{\beta \mu \nu }\right]+{\displaystyle \frac{1}{2}}\mathrm{}\left[f_{\alpha \beta }\left({\displaystyle \frac{1}{2}}\overline{\lambda }^\alpha D/\lambda ^\beta \right)\right]`$ (1)
$`{\displaystyle \frac{1}{8}}i\mathrm{}\left[f_{\alpha \beta }e^1D_\mu (e\overline{\lambda }^\alpha \gamma ^\mu \gamma _5\lambda ^\beta )\right]+{\displaystyle \frac{1}{4}}\overline{e}^{G/2}G^a(G^1)_a^b(f_{\alpha \beta }^{}/z^b\lambda ^\alpha \lambda ^\beta )+\mathrm{h}.\mathrm{c}.`$
Here $`\lambda ^\alpha `$ are the gaugino fields, $`G=\mathrm{ln}[\kappa ^6WW^{}]\kappa ^2d`$, where $`W`$ is the superpotential, $`d(z,z^{})`$ is the Kahler potential where $`z^a`$ are the complex scalar fields, and $`\kappa =(8\pi G_N)^{\frac{1}{2}}=0.41\times 10^{18}`$ GeV<sup>-1</sup> where $`G_N`$ is Newton’s constant. The gauge kinetic energy function $`f_{\alpha \beta }`$ in general has a non-trivial field dependence involving fields which transform as a singlet or a non-singlet irreducible representation of the underlying gauge group. After the spontaneous breaking of the gauge symmetry at the unification scale to the Standard Model gauge group $`SU(2)_L\times U(1)\times SU(3)_C`$, one needs to carry out a rescaling of the gauge kinetic energy in the sector of the gauge group that is preserved. This rescaling generates a splitting of the $`SU(2)_L\times U(1)\times SU(3)_C`$ gauge couplings at the unification scale $`M_X`$ and one has
$$\alpha _i^1(M_X)=\alpha _X^1(M_X)+\underset{r}{}c_{0r}n_i^r$$
(2)
where $`n_i^r`$ characterize the Higgs vacuum structure of the irreducible representation r and $`c_{0r}`$ parametrize its relative strength. After rescaling the gaugino mass matrix takes the form
$$m_{\alpha \beta }=\frac{1}{4}\overline{e}^{G/2}G^a(G^1)_a^b(f_{\alpha \gamma }^{}/z^b)f_{\gamma \beta }^1$$
(3)
Here one finds that the contribution to nonuniversality of the gaugino masses is controlled not only by the nature of the GUT sector but also by the nature of the hidden sector. Because of this the splitting of the gaugino masses at $`M_X`$ is characterized by a set of parameters $`c_r`$ different from $`c_{0r}`$. Thus after the breaking of the unified gauge group we parametrize the gaugino masses at $`M_X`$ by $`\stackrel{~}{m}_i(0)`$ where
$$\stackrel{~}{m}_i(0)=m_{\frac{1}{2}}(1+\underset{r}{}c_rn_i^r)$$
(4)
The effect of $`c_{0r}`$ on the gauge coupling unification has been discussed in the previous literature and we do not discuss it here. In the analysis of this paper we assume unification of gauge couplings at $`M_X`$ and the refinement of including $`c_{0r}`$ correction will not have any significant effect on our analysis.
In SU(5) the gauge kinetic energy function $`f_{\alpha \beta }`$ transforms as the symmetric product of $`\mathrm{𝟐𝟒}\times \mathrm{𝟐𝟒}`$ and contains the following representations
$$(\mathrm{𝟐𝟒}\times \mathrm{𝟐𝟒})_{symm}=\mathrm{𝟏}+\mathrm{𝟐𝟒}+\mathrm{𝟕𝟓}+\mathrm{𝟐𝟎𝟎}$$
(5)
The singlet leads to universality of the gaugino masses while the non-singlet terms will generate nonuniversality. We consider models where we have a linear combination of the singlet and a non-singlet representation, i.e., $`\mathrm{𝟏}+\mathrm{𝟐𝟒},\mathrm{𝟏}+\mathrm{𝟕𝟓}`$ or $`\mathrm{𝟏}+\mathrm{𝟐𝟎𝟎}`$. The quantities $`n_i^r`$ for the representations $`\mathrm{𝟏},\mathrm{𝟐𝟒},\mathrm{𝟕𝟓},\mathrm{𝟐𝟎𝟎}`$ are listed in Table1. Some phenomenological aspects of nonuniversality of gaugino masses have recently been discussed. We focus here on their effects on event rates in the direct detection of dark matter (see Ref. for previous work on the effects of gaugino mass nonuniversality on dark matter).
Techniques for computing the event rate in the scattering of neutralinos off nuclear targets has been discussed by many authors. We follow here the procedures discussed in Ref.. In our analysis we impose the $`bs+\gamma `$ constraint and the bounds on SUSY particles from the Tevatron and LEP. Specifically we take for the lower limits $`m_{\chi ^+}>94`$ GeV, $`m_\chi >33`$ GeV, $`m_{\stackrel{~}{\tau }_R}>71`$ GeV, $`m_{\stackrel{~}{t}}>87`$ GeV, $`m_{\stackrel{~}{g}}>190`$ GeV, $`m_h>113.5`$ GeV and for the $`bs+\gamma `$ branching ratio we take a $`2\sigma `$ range around the current experiment, i.e., we take $`2\times 10^4<B(bs+\gamma )<4.5\times 10^4`$. The quantity that constrains theory is $`\mathrm{\Omega }_\chi h^2`$ where $`\mathrm{\Omega }_\chi =\rho _\chi /\rho _c`$, where $`\rho _\chi `$ is the neutralino relic density and $`\rho _c=3H_0^2/8\pi G_N`$ is the critical matter density, and h is the value of the Hubble parameter $`H_0`$ in units of 100 km/sMpc. The current limit on h from the Hubble Space Telescope is $`h=0.71\pm 0.03\pm 0.07`$ and recent analyses of $`\mathrm{\Omega }_m`$ give $`\mathrm{\Omega }_m=0.3\pm 0.08`$. Assuming $`\mathrm{\Omega }_B0.05`$, one gets
$$\mathrm{\Omega }_\chi h^2=0.126\pm 0.043$$
(6)
In this analysis we make a a somewhat liberal choice for the error corridor on $`\mathrm{\Omega }_\chi h^2`$, i.e., we choose $`0.02\mathrm{\Omega }_\chi h^20.3`$. The choice of a more restricted corridor does not affect the general conclusions arrived at in this analysis. In the theoretical computation of the relic density we use
$`\mathrm{\Omega }_\chi h^22.48\times 10^{11}\left({\displaystyle \frac{T_\chi }{T_\gamma }}\right)^3\left({\displaystyle \frac{T_\gamma }{2.73}}\right)^3{\displaystyle \frac{N_f^{1/2}}{J(x_f)}}`$
$`J(x_f)={\displaystyle _0^{x_f}}𝑑x\sigma \upsilon (x)GeV^2`$ (7)
where $`(\frac{T_\chi }{T_\gamma })^3`$ is the reheating factor, $`N_f`$ is the number of degrees of freedom at the freeze-out temperature $`T_f`$ and $`x_f=kT_f/m_{\stackrel{~}{\chi }}`$. In determining $`J(x_f)`$ we use the method developed in Ref.. The role of $`J`$ in the context of nonuniversalities will be elucidated in Sec.5. A number of effects on neutralino dark matter have already been studied. These include the effects of nonuniversality of the scalar masses at the unification scale, effects of variations of the WIMP velocity, effects of rotation of the galaxy, effects of CP violation with EDM constraints, and effects of coannihilation. The focus of this analysis is on the effects of nonuniversality of the gaugino masses on dark matter. In the present analysis we do not include the effects of coannihilation. These effects become important when the NLSP mass $`m_i`$ lies close to the LSP mass, i.e., $`\mathrm{\Delta }_i=(m_i/m_\chi 1)<0.1`$. We have identified several regions of the parameter space where coannihilations involving $`\stackrel{~}{\tau }_1`$, $`\stackrel{~}{e}_R`$, the light chargino $`\chi _1^+`$, or the next to the lightest neutralino $`\chi _2^0`$ occur. However, as we stated above we do not consider coannihilation in this paper and thus eliminate such regions of the parameter space by imposing the constraint $`\mathrm{\Delta }_i>0.1`$. An analysis in this region requires a separate treatment and will be reported elsewhere. Recent analyses have pointed to the uncertainties in the quark masses and in the quark content of the nucleon that enter in analyses of dark matter. We give in this paper an independent analysis of the errors in the quark densities and compute their effects on dark matter.
The outline of the rest of the paper is as follows: In Sec.2 we discuss the basic formulae used to compute the neutralino-proton cross-section. An analysis of errors in the quark densities that enter in the scalar $`\sigma _{\chi p}`$ cross-section is also given. In Sec.3 we first give an analysis of $`\sigma _{\chi p}(scalar)`$ for the universal SUGRA case and analyze the effect of errors on the quark densities of the proton on it. We then discuss three nonuniversal scenarios where we consider admixtures of the singlet with the 24 plet, the 75 plet and the 200 plet representations for the gauge kinetic energy function. In Sec.4 we extend the analysis of $`\sigma _{\chi p}(scalar)`$ to the case of the O-II string model and a brane model based on 9 branes and $`5_i`$ branes. In Sec.5 we discuss the origin of the enlargement of the allowed LSP domain consistent with the relic density constraints due to the presence of nonuniversalities. Conclusions are given in Sec.6. In Appendix A we give an analytic solution of the sparticle masses using the one loop renormalization group equations including the effect of gaugino mass nonuniversalities. Using results of Appendix A we compute in Appendix B the effects of nonuniversalities on the $`\mu `$ parameter. The analytic results of Appendices A and B provide a deeper understanding of the gaugino-mass nonuniversality effects discussed in Secs. 3, 4 and 5.
## 2 Neutralino-proton cross-section
For heavy target nuclei such as germanium the neutralino-nucleus scattering cross-section is dominated by the scalar part of the neutralino-quark interaction and it is the quantity $`\sigma _{\chi p}(scalar)`$ on which constraints have been exhibited in the recent experimental works. For this reason we focus in this paper on the analysis of $`\sigma _{\chi p}(scalar)`$. The basic interaction governing the $`\chi p`$ scattering is the effective four-fermi interaction given by
$`_{eff}=\overline{\chi }\gamma _\mu \gamma _5\chi \overline{q}\gamma ^\mu (AP_L+BP_R)q+C\overline{\chi }\chi m_q\overline{q}q+D\overline{\chi }\gamma _5\chi m_q\overline{q}\gamma _5q+E\overline{\chi }i\gamma _5\chi m_q\overline{q}q`$
$`+F\overline{\chi }\chi m_q\overline{q}i\gamma _5q`$ (8)
where the interaction relevant to our analysis is parametrized by $`C`$. The $`\chi p`$ cross-section arising from scalar interactions is given by
$$\sigma _{\chi p}(scalar)=\frac{4\mu _r^2}{\pi }(\underset{i=u,d,s}{}f_i^pC_i+\frac{2}{27}(1\underset{i=u,d,s}{}f_i^p)\underset{a=c,b,t}{}C_a)^2$$
(9)
Here $`\mu _r`$ is the reduced mass, $`f_i^p`$ (i=u,d,s quarks) are defined by
$$m_pf_i^p=<p|m_{qi}\overline{q}_iq_i|p>$$
(10)
and C is given by
$$C=C_{h^0}+C_{H^0}+C_{\stackrel{~}{f}}$$
(11)
where $`C_{h^0},C_{H^0}`$ are the contributions from the s-channel $`h^0`$ and $`H^0`$ exchanges and $`C_{\stackrel{~}{f}}`$ is the contribution from the t-channel sfermion exchange. They are given by
$`C_{h^0}(u,d)=(+){\displaystyle \frac{g^2}{4M_WM_{h^0}^2}}{\displaystyle \frac{\mathrm{cos}\alpha (sin\alpha )}{\mathrm{sin}\beta (cos\beta )}}Re\sigma `$ (12)
$`C_{H^0}(u,d)={\displaystyle \frac{g^2}{4M_WM_{H^0}^2}}{\displaystyle \frac{\mathrm{sin}\alpha (cos\alpha )}{\mathrm{sin}\beta (cos\beta )}}Re\rho `$ (13)
$$C_{\stackrel{~}{f}}(u,d)=\frac{1}{4m_q}\frac{1}{M_{\stackrel{~}{q1}}^2M_\chi ^2}Re[C_{qL}C_{qR}^{}]\frac{1}{4m_q}\frac{1}{M_{\stackrel{~}{q2}}^2M_\chi ^2}Re[C_{qL}^{^{}}C_{qR}^{{}_{}{}^{}}]$$
(14)
Here (u,d) refer to the quark flavor, $`\alpha `$ is the Higgs mixing angle, and $`C_{qL},C_{qL}^{}`$ etc. are as defined in Ref., and $`\sigma `$ and $`\rho `$ are defined by
$$\sigma =X_{40}^{}(X_{20}^{}\mathrm{tan}\theta _WX_{10}^{})\mathrm{cos}\alpha +X_{30}^{}(X_{20}^{}\mathrm{tan}\theta _WX_{10}^{})\mathrm{sin}\alpha $$
(15)
$$\rho =X_{40}^{}(X_{20}^{}\mathrm{tan}\theta _WX_{10}^{})\mathrm{sin}\alpha +X_{30}^{}(X_{20}^{}\mathrm{tan}\theta _WX_{10}^{})\mathrm{cos}\alpha $$
(16)
where $`X_{n0}`$ are the components of the LSP
$$\chi =X_{10}^{}\stackrel{~}{B}+X_{20}^{}\stackrel{~}{W}_3+X_{30}^{}\stackrel{~}{H}_1+X_{40}^{}\stackrel{~}{H}_2$$
(17)
We discuss now the amount of uncertainty connected with the determination of $`f_i^p`$. The quantities that are used as inputs are $`\sigma _{\pi N}`$, x, and $`\xi `$ defined by
$$<p|2^1(m_u+m_d)(\overline{u}u+\overline{d}d|p>=\sigma _{\pi N},$$
(18)
$$x=\frac{\sigma _0}{\sigma _{\pi N}}=\frac{<p|\overline{u}u+\overline{d}d2\overline{s}s|p>}{<p|\overline{u}u+\overline{d}d|p>}$$
(19)
and
$$\xi =\frac{<p|\overline{u}u\overline{d}d|p>}{<p|\overline{u}u+\overline{d}d|p>}.$$
(20)
We can determine $`f_i^p`$ in terms of these and find
$`f_u^p={\displaystyle \frac{m_u}{m_u+m_d}}(1+\xi ){\displaystyle \frac{\sigma _{\pi N}}{m_p}}`$
$`f_d^p={\displaystyle \frac{m_d}{m_u+m_d}}(1\xi ){\displaystyle \frac{\sigma _{\pi N}}{m_p}}`$
$`f_s^p={\displaystyle \frac{m_s}{m_u+m_d}}(1x){\displaystyle \frac{\sigma _{\pi N}}{m_p}}`$ (21)
A similar analysis holds for the neutralino-neutron scattering and one can determine $`f_i^n`$ in terms of $`\xi ,x,\sigma _{\pi N}`$ as follows
$`f_u^n={\displaystyle \frac{m_u}{m_u+m_d}}(1\xi ){\displaystyle \frac{\sigma _{\pi N}}{m_p}}`$
$`f_d^n={\displaystyle \frac{m_d}{m_u+m_d}}(1+\xi ){\displaystyle \frac{\sigma _{\pi N}}{m_p}}`$
$`f_s^n={\displaystyle \frac{m_s}{m_u+m_d}}(1x){\displaystyle \frac{\sigma _{\pi N}}{m_p}}`$ (22)
We note in passing that from Eqs.21 and 22 one has the relation
$$f_u^pf_d^p=f_u^nf_d^n$$
(23)
which holds independent of the details of the input parameters. We discuss now the numerical evaluation of $`f_i^p`$ and $`f_i^n`$. The various determinations of $`\sigma _{\pi N}`$, $`\sigma _0`$ and x using analyses of Ref. are summarized in Table2. For $`\sigma _{\pi N}`$ one finds an average value of $`48\pm 9`$ MeV, and for $`\sigma _0`$ an average value of $`35.5\pm 6`$ MeV. These give $`x=0.74\pm 0.25`$. Further, there are two independent lattice gauge determinations of $`y=1x`$ which we list in Table2. In recording the result of $`x`$ for Ref. we have reduced the y value by $`35\%`$ as discussed in Ref.. The average of these lattice gauge calculations gives $`x=0.61\pm 0.08`$. Taking the average yet again of this $`x`$ and of $`\overline{\sigma }_0`$/$`\overline{\sigma }_{\pi N}`$ we get the average $`\overline{x}`$ listed in Table2. In addition to the above we need to determine the symmetry breaking parameter $`\xi `$. Here as in the work of Ref. we use the analysis of Ref. on baryon mass splittings to obtain
$$\xi =\frac{(\mathrm{\Xi }^{}+\mathrm{\Xi }^0\mathrm{\Sigma }^+\mathrm{\Sigma }^{})x}{\mathrm{\Xi }^{}+\mathrm{\Xi }^0+\mathrm{\Sigma }^++\mathrm{\Sigma }^{}2m_p2m_n}$$
(24)
where $`x`$ is as defined by Eq.(19). Numerically one finds $`\xi =0.196x`$ and on using Table 2 we find
$$\xi =0.132\pm 0.035$$
(25)
In addition to the above one needs the ratios of the quark masses for which we use
$$\frac{m_u}{m_d}=0.553\pm 0.043,\frac{m_s}{m_d}=18.9\pm 0.8$$
(26)
On using Eqs.(21),(25),(26), and Table 2 we find
$`f_u^p=0.021\pm 0.004`$
$`f_d^p=0.029\pm 0.006`$
$`f_s^p=0.21\pm 0.12`$ (27)
Similarly for $`f_i^n`$ we find
$`f_u^n=0.016\pm 0.003`$
$`f_d^n=0.037\pm 0.007`$
$`f_s^n=0.21\pm 0.12`$ (28)
For the more general case of neutralino-Nucleus ($`\chi N`$) scattering one has
$`\sigma _{\chi N}(scalar)={\displaystyle \frac{4m_r^2}{\pi }}(Z{\displaystyle \underset{i=u,d,s}{}}f_i^pC_i+{\displaystyle \frac{2}{27}}Z(1{\displaystyle \underset{i=u,d,s}{}}f_i^p){\displaystyle \underset{a=c,b,t}{}}C_a`$
$`+(AZ){\displaystyle \underset{i=u,d,s}{}}f_i^nC_i+{\displaystyle \frac{2}{27}}(AZ)(1{\displaystyle \underset{i=u,d,s}{}}f_i^n){\displaystyle \underset{a=c,b,t}{}}C_a)^2`$ (29)
Using Eqs.21 and 22 we write the above in the form
$`\sigma _{\chi N}(scalar)={\displaystyle \frac{4m_r^2A}{\pi }}(\widehat{f}{\displaystyle \frac{m_uC_u+m_dC_d}{m_u+m_d}}+\widehat{f}\xi \mathrm{\Delta }{\displaystyle \frac{m_uC_um_dC_d}{m_u+m_d}}+fC_s`$
$`+{\displaystyle \frac{2}{27}}(1f\widehat{f}\widehat{f}\xi \mathrm{\Delta }{\displaystyle \frac{m_um_d}{m_u+m_d}}){\displaystyle \underset{a=c,b,t}{}}C_a)^2`$ (30)
where $`\mathrm{\Delta }=(2ZA)/A`$, $`f=f_s`$, $`\widehat{f}=\sigma _{\pi N}/m_p`$ and numerically on using Table 2 we have
$$\widehat{f}=0.05\pm 0.01$$
(31)
We note that while the $`\chi p`$ and $`\chi n`$ cross-sections depend on $`\xi `$, the $`\xi `$ dependent term has a cancellation in ($`\chi N`$) cross-section because of the $`(2ZA)`$ factor. Further, if the target nucleus has $`A=2Z`$, i.e., $`\mathrm{\Delta }=0`$, then the $`\xi `$ dependent term will drop out of the $`\chi N`$ cross-section. Because of the above it is experimentally better to plot the $`\chi N`$ cross-section rather than the $`\chi p`$ cross-section as is currently the practice.
## 3 Dark matter in GUT models with gaugino mass nonuniversality
As discussed at the beginning of this section we consider here models where the nonuniversalities arise from admixtures of the singlet with the 24 plet, the 75 plet and the 200 plet representations. However, we begin first by exhibiting the result for the universal SUGRA case. The soft SUSY breaking sector of the theory, under the assumption that SUSY breaking is communicated from the hidden to the visible sector by gravitational interactions, is parametrized in this case by the universal scalar mass $`m_0`$, the universal gaugino mass $`m_{\frac{1}{2}}`$, the universal trilinear coupling $`A_0`$ all taken at the GUT scale, and $`\mathrm{tan}\beta =<H_2>/<H_1>`$ where $`H_2`$ gives mass to the up quark and $`H_1`$ gives mass to the down quark. Throughout this analysis we assume that the Higgs mixing parameter $`\mu `$ (which appears in the superpotential as $`\mu H_1H_2`$) is determined via the electro-weak symmetry breaking constraint. The range of the parameters are limited by a naturalness constraint. We mean this to imply that $`m_0,m_{\stackrel{~}{g}}1`$ TeV, where $`m_{\stackrel{~}{g}}`$ is the gluino mass, $`\mathrm{tan}\beta 25`$, and $`A_0`$, or equivalently $`A_t`$, the value of $`A_0`$ at the electro-weak scale in the top channel, is limited by the electro-weak symmetry breaking constraint. For the analysis here we choose $`\mu >0`$ while for the other $`\mu `$ sign the allowed parameter space for dark matter is strongly limited due to the $`bs+\gamma `$ constraint. The results for $`\mu <0`$ look qualitatively different in that the cross-sections are significantly smaller.
In Fig.1 we plot the maximum and the minimum of $`\sigma _{\chi p}(scalar)`$ as a function of $`m_\chi `$ where the parameters are allowed to vary over the naturalness range discussed above. The analysis is done for three sets of $`f_i^p`$ values corresponding to the corridor given by Eq.(27). They correspond to (I)$`f_u^p=0.025,f_d^p=0.035,f_s^p=0.33`$, (II)$`f_u^p=0.021,f_d^p=0.029,f_s^p=0.21`$, (III)$`f_u^p=0.017,f_d^p=0.023,f_s^p=0.09`$. From Fig.1 we see that the different sets can lead to a variation in $`\sigma _{\chi p}(scalar)`$ up to a factor of about 5. For the rest of the analysis in this paper we use set (II).
Next we consider the $`\mathrm{𝟏}+\mathrm{𝟐𝟒}`$ model. We find that in this case $`\sigma _{\chi p}(scalar)`$ typically decreases for positive values of $`c_{24}`$ and increases with negative values of $`c_{24}`$. This behavior arises primarily from the dependence of the gaugino-Higgsino components $`X_{n0}`$ of the LSP on the gaugino mass nonuniversality. Thus $`X_{n0}`$ are sensitive to the gaugino nonuniversality through their dependence on $`\stackrel{~}{m}_1`$, $`\stackrel{~}{m}_2`$ and $`\mu `$. In Fig.2 we exhibit the dependence of $`X_{n0}`$ on $`C_{24}`$ for some typical input values. The quantity $`\sigma _{\chi p}(scalar)`$ depends on the direct product of the gaugino and Higgsino components of $`\chi `$. Specifically, $`\sigma _{\chi p}(scalar)`$ vanishes if $`\chi `$ is a pure Bino. From Fig.2 we see that negative values of $`c_{24}`$ increase the Higgsino components and hence increase the neutralino-quark scattering and lead to an enhancement of $`\sigma _{\chi p}(scalar)`$ while the opposite situation is realized for positive values of $`c_{24}`$. This is what is found in Fig.3 where we give a plot of the maximum and the minimum of $`\sigma _{\chi p}(scalar)`$ for the cases $`c_{24}=0.08`$, $`c_{24}=0`$ and $`c_{24}=0.1`$ when the soft SUSY breaking parameters are varied over the assumed naturalness range as in Fig.1. The analysis shows that for $`c_{24}=0.1`$ an enhancement of $`\sigma _{\chi p}(scalar)`$ by as much as a factor of 5 can occur as a result of the gaugino mass nonuniversality and the allowed range of $`m_\chi `$ is also increased in this case beyond the range allowed in the universal SUGRA case.
In Fig.4 we give an analysis of the maximum and the minimum of $`\sigma _{\chi p}(scalar)`$ for the $`\mathrm{𝟏}+\mathrm{𝟕𝟓}`$ case for three different values of $`c_{75}`$, i.e., $`c_{75}=0.06`$, $`c_{75}=0`$ and $`c_{75}=0.04`$. As for the $`\mathrm{𝟏}+\mathrm{𝟐𝟒}`$ case, $`\sigma _{\chi p}(scalar)`$ typically increases for negative values of $`c_{\mathrm{𝟕𝟓}}`$ and decreases for positive values of $`c_{\mathrm{𝟕𝟓}}`$. Again this can be understood by analysing the gaugino-Higgsino components of the LSP as a function of $`c_{\mathrm{𝟕𝟓}}`$. Thus here as in the $`\mathrm{𝟏}+\mathrm{𝟐𝟒}`$ case one finds that the Higgsino components of the LSP increase as $`c_{\mathrm{𝟕𝟓}}`$ decreases and decrease as $`c_{\mathrm{𝟕𝟓}}`$ increases. Because of this there is an enhancement of $`\sigma _{\chi p}(scalar)`$ for $`c_{75}<0`$. A comparison of $`c_{75}=0`$ and $`c_{75}=0.06`$ cases in Fig.4 shows that an enhancement of $`\sigma _{\chi p}(scalar)`$ up to a factor of 5 or more occurs in this case. As in the case of $`\mathrm{𝟏}+\mathrm{𝟐𝟒}`$ here also one finds that the allowed range of $`m_\chi `$ consistent with the constraints is extended beyond the values allowed in the universal SUGRA case.
In Fig.5 we give an analysis of the maximum and the minimum of the $`\sigma _{\chi p}(scalar)`$ for the $`\mathrm{𝟏}+\mathrm{𝟐𝟎𝟎}`$ case when $`c_{200}=0.08`$, $`c_{200}=0`$ and $`c_{200}=0.1`$. In this case the dependence of $`\sigma _{\chi p}(scalar)`$ on $`c_{200}`$ is opposite to that one has in the previous cases, i.e., the $`\mathrm{𝟏}+\mathrm{𝟐𝟒}`$ case and the $`\mathrm{𝟏}+\mathrm{𝟕𝟓}`$ case. Here it is for positive values of $`c_{200}`$ that $`\sigma _{\chi p}(scalar)`$ increases and it is for negative values of $`c_{200}`$ that $`\sigma _{\chi p}(scalar)`$ decreases. The origin of this reversal lies in $`n_i^r`$ (see Table 1) and can be understood from the dependence of $`X_{n0}`$ on $`c_{200}`$. The above implies that the Bino component of $`\chi `$ decreases and the Higgsino components increase for $`c_{200}>0`$ while the opposite situation occurs when $`c_{200}<0`$. This dependence implies that $`\sigma _{\chi p}(scalar)`$ should increase for $`c_{200}>0`$ and decrease for $`c_{200}<0`$ which is what is observed in Fig.5. In this case one finds that $`\sigma _{\chi p}(scalar)`$ can increase up to a factor of 10 or more because of nonuniversality. The analysis also shows that as in the case of $`\mathrm{𝟏}+\mathrm{𝟐𝟒}`$ and $`\mathrm{𝟏}+\mathrm{𝟕𝟓}`$ the allowed range of the LSP mass is extended beyond what is allowed in the universal SUGRA case.
One can gain a deeper understanding of the dependence of $`X_{n0}`$ on $`c_r`$ and hence a deeper understanding of the dependence of $`\sigma _{\chi p}(scalar)`$ on $`c_r`$ from studying the dependence of the Higgs mixing parameter $`\mu `$ on $`c_r`$. To appreciate why $`\mu `$ is such an important parameter in this discussion it is useful to look at large $`\mu `$, i.e., the case $`\mu ^2>>M_Z^2`$. In this limit one finds that the LSP eigenvector is given by
$`X_{10}1(M_Z^2/2\mu ^2)\mathrm{sin}^2\theta _W,X_{20}(M_Z^2/2(\stackrel{~}{m}_2\stackrel{~}{m}_1)\mu )\mathrm{sin}2\theta _W\mathrm{sin}2\beta ,`$
$`X_{30}(M_Z/\mu )sin\theta _Wsin\beta ,X_{40}(M_Z/\mu )sin\theta _W\mathrm{cos}\beta ,`$ (32)
Eq.32 shows that in the large $`\mu `$ limit $`\chi `$ is mostly a Bino and the corrections to the pure Bino limit are proportional to $`(M_Z^2/\mu ^2)`$ while the Higgsino components are proportional to $`(M_Z/\mu )`$. As already pointed out the $`\sigma _{\chi p}(scalar)`$ depends on the interference of the gaugino and Higgsino components of the LSP, i.e., $`X_{i0}\times X_{\alpha 0}`$ (i=1,2; $`\alpha =3,4`$). Clearly then as $`|\mu |`$ increases we go more deeply into the pure Bino region reducing $`\sigma _{\chi p}(scalar)`$. Likewise as $`|\mu |`$ decreases $`\chi `$ develops larger Higgsino components $`X_{\alpha 0}`$($`\alpha =3,4`$), even though it is still dominantly a Bino, $`\sigma _{\chi p}(scalar)`$ decreases. Thus to gain an insight on the effect of $`c_r`$ on $`X_{n0}`$ and hence on the effect of $`c_r`$ on $`\sigma _{\chi p}(scalar)`$ we need to understand how $`c_r`$ affects $`\mu `$. We address this topic below.
In SUGRA models $`\mu ^2`$ is determined via the breaking of the electro-weak symmetry and thus depends on the gaugino mass nonuniversality through the Higgs mass parameters (see Eqs.39 and 43 in Appendix A and Eq.50 in Appendix B). We can understand the effect of nonuniversality on $`\mu `$ analytically by expanding $`\mu `$ for the nonuniversal case around the universal value using $`c_r`$ as an expansion parameter
$$\stackrel{~}{\mu }^2=\mu _0^2+\underset{r}{}\frac{\mu ^2}{c_r}c_r+O(c_r^2)$$
(33)
where $`\mu _0`$ is the value of $`\mu `$ for the universal case. Using the analysis of Appendix B the pattern of breaking in the Higgs structure of $`\mathrm{𝟐𝟒},\mathrm{𝟕𝟓},\mathrm{𝟐𝟎𝟎}`$ shows that
$$\frac{\mu _{24}^2}{c_{24}}>0,\frac{\mu _{75}^2}{c_{75}}>0,\frac{\mu _{200}^2}{c_{200}}<0$$
(34)
Thus in the neighborhood of $`c_r=0`$ a negative value of $`c_{24}`$ gives a smaller value of $`|\mu |`$ leading to larger Higgsino components in $`\chi `$ and hence a larger $`\sigma _{\chi p}(scalar)`$ as is observed in the numerical analysis. A similar situation holds for the nonuniversality effects from $`c_{75}`$. However, for the nonuniversality effects from $`c_{200}`$ an opposite situation holds because of the opposite sign of the derivative term as given by Eq.34. More generally one finds the same behavior in a larger $`c_r`$ domain, i.e., $`\mu _{24}^2<\mu _0^2`$ for $`c_{24}<0`$, $`\mu _{75}^2<\mu _0^2`$ for $`c_{75}<0`$, and $`\mu _{200}^2<\mu _0^2`$ for $`c_{200}>0`$ and observations similar to those valid for small $`c_r`$ also apply here. These results imply that gaugino nonuniversality which makes $`\mu `$ small produces a deviation of the LSP from the approximate Bino limit in a direction which leads to a larger value of $`\sigma _{\chi p}(scalar)`$.
Nonuniversality of the gaugino masses also has implications for naturalness. One convenient definition of naturalness is via the fine tuning parameter $`\mathrm{\Phi }`$ introduced in Ref. defined by $`\mathrm{\Phi }=\frac{1}{4}+\frac{\mu ^2}{M_Z^2}`$. Using this definition one can easily compare values of fine tuning for the universal and nonuniversal cases. One finds $`\mathrm{\Phi }_{24}<\mathrm{\Phi }_0`$ $`(c_{24}<0)`$, $`\mathrm{\Phi }_{75}<\mathrm{\Phi }_0(c_{75}<0)`$, $`\mathrm{\Phi }_{200}<\mathrm{\Phi }_0(c_{200}>0)`$, where $`\mathrm{\Phi }_{24}`$ is $`\mathrm{\Phi }`$ for the case $`\mathrm{𝟏}+\mathrm{𝟐𝟒}`$ etc, and $`\mathrm{\Phi }_0`$ is $`\mathrm{\Phi }`$ for the universal case. Since the correction to $`\mu ^2`$ is negative for the case of larger Higgsino components one finds that the deviation from the approximate Bino limit is in the direction of a smaller value of $`\mathrm{\Phi }`$ and towards the direction of greater naturalness relative to the universal case. Thus a smaller $`\mu `$ leads to a larger $`\sigma _{\chi p}`$ and a larger detection rate and also makes the model more natural by making $`\mathrm{\Phi }`$ small. In this sense the more natural the SUSY model the larger is the detection rate.
## 4 Dark matter in string/brane models
One of the main hurdles in the analysis of SUSY phenomenology based on string models is that there is as yet not a full understanding of the breaking of supersymmetry here. However, there do exist efficient ways to parametrize SUSY breaking and one such parametrization is
$`F^S=\sqrt{3}m_{\frac{3}{2}}(S+S^{})\mathrm{sin}\theta e^{i\gamma _S},F^i=\sqrt{3}m_{\frac{3}{2}}(T+T^{})\mathrm{cos}\theta \mathrm{\Theta }_ie^{i\gamma _i}`$ (35)
where $`F^S`$ is the dilation VEV, $`F^i`$ are the moduli VEVs, $`\theta `$ ($`\mathrm{\Theta }_i`$) parametrizes the Goldstino direction in the S ($`T_i`$) field space and $`\gamma _S`$ ($`\gamma _i`$) is the phase. The $`\mathrm{\Theta }_i`$ obey the constraint $`\mathrm{\Theta }_1^2+\mathrm{\Theta }_2^2+\mathrm{\Theta }_3^2=1`$. We begin with an example of the O-II string model with the soft SUSY breaking sector parametrized by
$`\stackrel{~}{m}_i=\sqrt{3}m_{\frac{3}{2}}(\mathrm{sin}\theta e^{i\alpha _S}\gamma _iϵ\mathrm{cos}\theta e^{i\alpha _T})`$
$`m_0^2=ϵ^{}(\delta _{GS})m_{\frac{3}{2}}^2,A_0=\sqrt{3}m_{\frac{3}{2}}\mathrm{sin}\theta e^{i\alpha _S}`$ (36)
Here $`\gamma _1=\frac{33}{5}+\delta _{GS},\gamma _2=1+\delta _{GS},\gamma _3=3+\delta _{GS}`$ where $`\delta _{GS}`$ is the Greene-Schwarz parameter which is fixed by the constraint of anomaly cancellation in a given orbifold model. Further, as in the GUT analyses we treat $`\mu `$ as an independent parameter. The phenomenology of this model has been discussed in Ref. and the EDM constraints in Refs.. However, in the analysis below we do not impose the accelerator constraints and as in the analysis for GUT models we do not assume CP violation and thus set the CP phases to zero.
The stucture of the soft SUSY parameters for the model discussed above shows that the nonuniversality of the gaugino masses is controlled by several parameters in this case: $`\delta _{GS}`$, $`\theta `$ and $`ϵ`$ which play a role similar to the role played by the parameters $`c_r`$ in the case of the GUT models. The presence of several parameters leads to many different possibilities for generating gaugino mass nonuniversality. In addition, there is a new feature in this string model, not present in GUT models, in that the universal scalar $`(mass)^2`$ at the unification scale, i.e. $`m_0^2`$, is dependent on $`\delta _{GS}`$ which therefore correlates the universal scalar mass with the gaugino mass nonuniversality. An analysis of the maximum and the minimum curves for $`\sigma _{\chi p}(scalar)`$ as a function of $`m_\chi `$ is given in Fig.6 when the parameters in the model are varied over a range of allowed values. A comparison with Fig.1 shows that $`\sigma _{\chi p}(scalar)`$ can be larger than for the universal SUGRA case by a factor of as much as 10. One also finds that the range of the neutralino mass extends to about 575 GeV significantly beyond what one finds in the universal SUGRA case even without inclusion of the coannihilation effects.
We discuss next dark matter for a class of D brane models. Over the recent past there has been considerable interest in the study of Type IIB orientifolds and their compactifications. We consider here models with compactifications on a six torus of the type $`T^6=T^2\times T^2\times T^2`$. In models of this type one has a set of 9 branes, $`7_i`$ (i=1,2,3) branes, $`5_i`$ branes and 3 branes. This set is further constrained by the requirement of N=1 supersymmetry which requires that on has either 9 branes and $`5_i`$ branes, or $`7_i`$ branes and 3 branes. Model building on branes allows an additional flexibility in that one can associate different parts of the Standard Model gauge group with different branes. In Ref. a brane model using two five branes $`5_1`$ and $`5_2`$ was investigated while in Ref. models using 9 brane and $`5_i`$ brane were investigated. We pursue here the implications of the latter possibility. In one of the models of Ref. the Standard Model gauge group is associated with the branes in the following way: the $`SU(3)_C\times U(1)_Y`$ is associated with the 9 brane while $`SU(2)_L`$ is associated with the $`5_1`$ brane. Further, it is assumed that the $`SU(2)_R`$ singlets are associated with the 9 brane while the $`SU(2)_L`$ doublets are associated with the intersection of 9 brane and $`5_1`$ brane. The soft SUSY breaking sector of this $`95_1`$ brane model is then given as follows: The gaugino masses $`\stackrel{~}{m}_i`$ (i=1,2,3) corresponding to the gauge group $`SU(3)`$, $`SU(2)`$ and $`U(1)`$ are parametrized by
$`\stackrel{~}{m}_1=\sqrt{3}m_{\frac{3}{2}}\mathrm{sin}\theta e^{i\gamma _S}=\stackrel{~}{m}_3=A_0,\stackrel{~}{m}_2=\sqrt{3}m_{\frac{3}{2}}\mathrm{cos}\theta \mathrm{\Theta }_1e^{i\gamma _1}`$ (37)
while the $`SU(2)_L`$ singlets are parametrized by $`m_9`$ and $`SU(2)_L`$ doublets are parametrized by $`m_{95_1}`$ where
$$m_9^2=m_{\frac{3}{2}}^2(13\mathrm{cos}^2\theta \mathrm{\Theta }_1^2),m_{95_1}^2=m_{\frac{3}{2}}^2(1\frac{3}{2}\mathrm{cos}^2\theta (1\mathrm{\Theta }_1^2)).$$
(38)
Here the $`\theta `$ and $`\mathrm{\Theta }_1`$ are the directions of the Goldstino in the dilaton and the moduli VEV space as discussed earlier. To avoid generating tachyons in the theory one needs to impose the constraint $`cos^2\theta \mathrm{\Theta }_1^2<1`$. A second version of this model was also discussed in Ref.. Here one associates $`SU(3)\times U(1)_Y`$ with the $`5_1`$ brane while the $`SU(2)_L`$ with the 9 brane, and assumes that the $`SU(2)_R`$ singlets were associated with the $`5_1`$ brane while the $`SU(2)_L`$ doublets are associated with the intersection of 9 brane and $`5_1`$ brane as before. The soft SUSY breaking sector of this model can be gotten from the model discussed above by the interchange $`\mathrm{cos}\theta \mathrm{\Theta }_1\mathrm{sin}\theta `$. In the analysis of this paper we focus on the version of the model given by Eqs.37 and 38. As usual we assume that the parameter $`\mu `$ is free and we determine it via the constraint of the radiative breaking of the electro-weak symmetry. Again as in other cases we have considered we set the CP phases to zero. Eqs.37 and 38 show that one has nonuniversality at the unification scale both in the scalar sector as well as in the gaugino sector. The limit of universal scalar mass corresponds to $`\mathrm{\Theta }_1=1/\sqrt{3}`$. Our focus in this paper is on the nonuniversality in the gaugino sector, and so for the numerical analysis we set $`\mathrm{\Theta }_1=1/\sqrt{3}`$. In this case the numerical analysis shows that the allowed neutralino mass range extends up to about 650 GeV.
## 5 Gaugino mass nonuniversality and LSP mass range
In SUGRA models with universal boundary conditions at the GUT scale, the allowed range of the LSP typically does not exceed 200 GeV and with imposition of additional constraints it is often significantly less. As discussed in Secs.3 and 4 in the presence of gaugino mass nonuniversality one finds that the allowed LSP range is increased. For the $`\mathrm{𝟏}+\mathrm{𝟐𝟒}`$ case the allowed LSP range extends to about $`220`$ GeV for the case $`c_{24}=0.1`$. For the $`\mathrm{𝟏}+\mathrm{𝟕𝟓}`$ case one finds that $`c_{75}=0.06`$ gives an LSP range which extends to 250 GeV, while for the $`\mathrm{𝟏}+\mathrm{𝟐𝟎𝟎}`$ case with $`c_{200}=0.1`$ the allowed LSP range extends to 375 GeV. A similar situation holds for the string/brane models. Here one finds that the allowed LSP range can extend up to 600 GeV. These extended regions arise even without inclusion of coannihilation effects which is known to extend the allowed regions also up to about 700 GeV. In the regions of the parameters space considered the effects of coannihilations would infact be negligible since we are considering only those configurations for which $`\mathrm{\Delta }_i>0.1`$. The specific mechanism by which the neutralino mass range is extended is also different than in the case of coannihilation. Thus for the case of coannihilation the increase in the allowed LSP range occurs as a consequence of a coupled channel effect while in the the case of nonuniversalities the extension of the allowed region of the LSP arises due to a significant increase in the value of J for certain ranges of nonuniversalities. Thus one may expand J as the sum over the final state channels in the $`\chi \chi `$ annihilation so that $`J`$=$`J(\stackrel{~}{f}\stackrel{~}{f})+J(WW)+J(ZZ)`$\+ $`J(Zh^0)`$+$`J(ZH^0)+etc`$. The region of the nonuniversality parameter space which leads to an enhancement of $`\sigma _{\chi p}`$ can also lead to an enhancement of the cross-sections for the $`\chi \chi `$ annihilation into the final states $`W^+W^{},ZZ,Zh,ZH`$ etc and hence to an increase in J which leads to a decrease in the relic density down to permissible limits consistent with constraints. Thus regions of the parameter space which would otherwise be excluded are now included when the gaugino mass nonuniversalities are included.
## 6 Conclusion
In this paper we have analyzed the effects of the nonuniversality of the gaugino masses on neutralino dark matter in SUGRA, string and D brane models under the constraint of R parity invariance. It is found that nonuniversality effects can enhance the $`\chi p`$ cross-section for scalar interactions by as much as a factor of 10. We also carried out an analysis of the uncertainties in the numerical determination of $`f_i^p`$ (i=u,d,s) and find that with the current state of uncertainties the $`\chi p`$ cross-section cannot be pinned down to better than a factor of about 5. Our analysis of the gaugino mass nonuniversality also exhibits another important phenomena, i.e., that the allowed range of the neutralino mass can be extended up to about 600 GeV even without inclusion of the coannihilation effects. The effects of coannihilation were not considered and inclusion of these effects may further increase the allowed neutralino mass range. The extended LSP mass range should be of interest in the experimental searches for dark matter in the current dark matter detectors and in the design of new dark matter detectors which are at the planning stage.
Acknowledgements
This research was supported in part by NSF grant PHY-9901057.
Appendix A: Effects of gaugino mass nonuniversality on sparticle masses
In this appendix we give analytic solutions to the one loop renormalization group equations including the effect of gaugino mass nonuniversality. The analytic formulae for these with universal boundary conditions were given in Ref. and for nonuniversality in the scalar mass sector in Ref.. Here we limit ourselves to the gaugino mass nonuniversality. These formulae are found useful in gaining an analytic understanding of the nonuniversality effects. The one loop RG formulae are given in several papers (see, e.g. Ref.) and we do not reproduce them here. Rather we discuss the solutions under the boundary conditions where $`m_0`$ is the universal scalar mass, $`\stackrel{~}{m}_i(0)`$ (i=1,2,3) are the gaugino masses for the gauge group sectors $`U(1),SU(2),SU(3)`$ and $`A_0`$ is the universal trilinear coupling all taken at the unification scale. The Higgs mass parameters, the trilinear couplings, and the squark and slepton masses at the electro-weak scale are all sensitive to the effect of gaugino mass nonuniversality. The simplest case is that of the mass parameter $`m_{H_1}^2`$ for the $`H_1`$ Higgs which couples to the down quark. Here one finds
$$m_{H_1}^2=m_0^2+\stackrel{~}{\alpha }_G(\frac{3}{2}\stackrel{~}{f}_2(t)+\frac{3}{10}\stackrel{~}{f}_1(t))m_{\frac{1}{2}}^2$$
(39)
where
$$\stackrel{~}{f}_i(t)=\frac{1}{\beta _i}(1\frac{1}{(1+\beta _it)^2})(\frac{\stackrel{~}{\alpha }_i(0)}{\stackrel{~}{\alpha }_G})(\frac{\stackrel{~}{m}_i(0)}{m_{\frac{1}{2}}})^2$$
(40)
Here $`t=ln(M_X^2/Q^2)`$, $`\beta _i=(b_i/4\pi )\stackrel{~}{\alpha }_i(0)`$ where $`b_i=(33/5,1,3)`$ for $`U(1)`$, $`SU(2)`$, and $`SU(3)`$, and $`\stackrel{~}{\alpha }_i(0)=\alpha _i/4\pi `$. The $`\stackrel{~}{f}_i(t)`$ contain the nonuniversality effects. The evolution of the mass parameter for the Higgs $`H_2`$ involves the evolution of the trilinear coupling in the stop channel at the electro-weak scale and this coupling is given by
$$A_t(t)=\frac{A_0}{1+6Y_0F}+m_{\frac{1}{2}}(\stackrel{~}{H}_2\frac{6Y_0\stackrel{~}{H}_3}{1+6Y_0F})$$
(41)
where
$`\stackrel{~}{H}_2=\stackrel{~}{\alpha }_G({\displaystyle \frac{16}{3}}\stackrel{~}{h}_3+3\stackrel{~}{h}_2+{\displaystyle \frac{13}{15}}\stackrel{~}{h}_1),\stackrel{~}{H}_3={\displaystyle _0^t}E(t)\stackrel{~}{H}_2(t)`$
$`\stackrel{~}{h}_i={\displaystyle \frac{t}{1+\beta _it}}({\displaystyle \frac{\stackrel{~}{\alpha }_i(0)}{\stackrel{~}{\alpha }_G}})({\displaystyle \frac{\stackrel{~}{m}_i(0)}{m_{\frac{1}{2}}}})`$ (42)
where $`Y_0`$ is the top Yukawa coupling at the GUT scale and the functions E and F are as defined in Ref.. We note that the first term on the right hand side of Eq.(40) which arises purely from the top Yukawa coupling evolution is unaffected by nonuniversality while the second term in affected through the modification of $`\stackrel{~}{h}_i`$. For the mass parameter $`m_{H_2}`$ for Higgs $`H_2`$ that couples with the top one finds
$$m_{H_2}^2=m_{\frac{1}{2}}^2\stackrel{~}{e}(t)+A_0m_0m_{\frac{1}{2}}\stackrel{~}{f}(t)+m_0^2(h(t)k(t)A_0^2)$$
(43)
where the functions $`h(t)`$ and $`k(t)`$ are unaffected by nonuniversality and are as given in Ref. while the functions $`\stackrel{~}{e}`$ and $`\stackrel{~}{f}`$ are modified due to nonuniversality. $`\stackrel{~}{e}`$ and $`\stackrel{~}{f}`$ are given by
$$\stackrel{~}{e}(t)=\frac{3}{2}[\frac{\stackrel{~}{G}_1+Y_0\stackrel{~}{G}_2}{D(t)}+\frac{(\stackrel{~}{H}_2+6Y_0\stackrel{~}{H}_4)^2}{3D(t)^2}+\stackrel{~}{H}_8]$$
(44)
and
$$\stackrel{~}{f}(t)=\frac{6Y_0\stackrel{~}{H}_3(t)}{D(t)^2};D(t)=(1+6Y_0F(t))$$
(45)
Here the various tilde functions containing the nonuniversality are defined below
$`\stackrel{~}{G}_1(t)=\stackrel{~}{F}_2(t){\displaystyle \frac{1}{3}}(\stackrel{~}{H}_2)^2`$
$`\stackrel{~}{G}_2(t)=6\stackrel{~}{F}_3(t)\stackrel{~}{F}_4(t)4\stackrel{~}{H}_2(t)\stackrel{~}{H}_4(t)+2F(\stackrel{~}{H}_2)^22\stackrel{~}{H}_6(t)`$
$`\stackrel{~}{F}_3(t)=F(t)\stackrel{~}{F}_2(t){\displaystyle _0^t}𝑑t^{}E(t^{})\stackrel{~}{F}_2(t^{}),\stackrel{~}{F}_4(t)={\displaystyle _0^t}𝑑t^{}E(t^{})\stackrel{~}{H}_5(t^{})`$
$`\stackrel{~}{H}_4(t)=F(t)\stackrel{~}{H}_2(t)\stackrel{~}{H}_3(t),\stackrel{~}{H}_5(t)=\stackrel{~}{\alpha }_G({\displaystyle \frac{16}{3}}\stackrel{~}{f}_3(t)+6\stackrel{~}{f}_2(t){\displaystyle \frac{22}{15}}\stackrel{~}{f}_1)`$
$`\stackrel{~}{H}_6(t)={\displaystyle _o^t}𝑑t^{}(\stackrel{~}{H}_2)^2E(t^{}),\stackrel{~}{H}_8(t)=\stackrel{~}{\alpha }_G({\displaystyle \frac{8}{3}}\stackrel{~}{f}_3(t)+\stackrel{~}{f}_2(t){\displaystyle \frac{1}{3}}\stackrel{~}{f}_1)`$ (46)
The squark and slepton masses are also affected by nonuniversality. For the up squarks in the first two generations one finds
$`m_{\stackrel{~}{u}_{iL}}^2(t)=m_0^2+m_{ui}^2+\stackrel{~}{\alpha }_G[{\displaystyle \frac{8}{3}}\stackrel{~}{f}_3+{\displaystyle \frac{3}{2}}\stackrel{~}{f}_2+{\displaystyle \frac{1}{30}}\stackrel{~}{f}_1]m_{\frac{1}{2}}^2+({\displaystyle \frac{1}{2}}{\displaystyle \frac{2}{3}}\mathrm{sin}^2\theta _W)M_Z^2\mathrm{cos}2\beta `$
$`m_{\stackrel{~}{u}_{iR}}^2(t)=m_0^2+m_{ui}^2+\stackrel{~}{\alpha }_G[{\displaystyle \frac{8}{3}}\stackrel{~}{f}_3+{\displaystyle \frac{8}{15}}\stackrel{~}{f}_1]m_{\frac{1}{2}}^2+{\displaystyle \frac{2}{3}}\mathrm{sin}^2\theta _WM_Z^2\mathrm{cos}2\beta `$ (47)
The analysis of the first two generations of the down quarks and of the sleptons is similar. Finally we discuss the third generation squarks. Here one finds
$`m_{\stackrel{~}{b}_L}^2(t)={\displaystyle \frac{1}{2}}m_0^2+m_b^2+{\displaystyle \frac{1}{2}}m_U^2++\stackrel{~}{\alpha }_G[{\displaystyle \frac{4}{3}}\stackrel{~}{f}_3+{\displaystyle \frac{1}{15}}\stackrel{~}{f}_1]m_{\frac{1}{2}}^2+({\displaystyle \frac{1}{2}}+{\displaystyle \frac{1}{3}}\mathrm{sin}^2\theta _W)M_Z^2\mathrm{cos}2\beta `$
$`m_{\stackrel{~}{b}_R}^2(t)=m_0^2+m_b^2+\stackrel{~}{\alpha }_G[{\displaystyle \frac{8}{3}}\stackrel{~}{f}_3+{\displaystyle \frac{2}{15}}\stackrel{~}{f}_1]m_{\frac{1}{2}}^2{\displaystyle \frac{1}{3}}\mathrm{sin}^2\theta _WM_Z^2\mathrm{cos}2\beta `$
$`m_{\stackrel{~}{t}_L}^2(t)=m_Q^2+m_t^2+({\displaystyle \frac{1}{2}}{\displaystyle \frac{2}{3}}\mathrm{sin}^2\theta _W)M_Z^2\mathrm{cos}2\beta `$
$`m_{\stackrel{~}{t}_R}^2(t)=m_U^2+m_t^2+{\displaystyle \frac{2}{3}}\mathrm{sin}^2\theta _WM_Z^2\mathrm{cos}2\beta `$ (48)
where $`m_U^2`$ and $`m_Q^2`$ are defined by
$`m_U^2={\displaystyle \frac{1}{3}}m_0^2+{\displaystyle \frac{2}{3}}\stackrel{~}{f}A_0m_{\frac{1}{2}}{\displaystyle \frac{2}{3}}kA_0^2+{\displaystyle \frac{2}{3}}hm_0^2+[{\displaystyle \frac{2}{3}}\stackrel{~}{e}+\stackrel{~}{\alpha }_G({\displaystyle \frac{8}{3}}\stackrel{~}{f}_3\stackrel{~}{f}_2+{\displaystyle \frac{1}{3}}\stackrel{~}{f}_1)]`$
$`m_Q^2={\displaystyle \frac{2}{3}}m_0^2+{\displaystyle \frac{1}{3}}\stackrel{~}{f}A_0m_{\frac{1}{2}}{\displaystyle \frac{1}{3}}kA_0^2+{\displaystyle \frac{1}{3}}hm_0^2+[{\displaystyle \frac{1}{3}}\stackrel{~}{e}+\stackrel{~}{\alpha }_G({\displaystyle \frac{8}{3}}\stackrel{~}{f}_3+\stackrel{~}{f}_2{\displaystyle \frac{1}{15}}\stackrel{~}{f}_1)]`$ (49)
All the sparticle mass relations limit to the universal case when we set
$`\alpha _i(0)\stackrel{~}{m}_i(0)/\alpha _Gm_{\frac{1}{2}}=1`$.
Appendix B : nonuniversality effects on $`\mu `$
Since $`\mu `$ is determined via the constraint of electro-weak symmetry breaking it is sensitive to the gaugino mass nonuniversality. Thus one has
$$\mu ^2=(m_{H_1}^2m_{H_2}^2\mathrm{tan}\beta ^2)(\mathrm{tan}\beta ^21)^1\frac{1}{2}M_Z^2+\mathrm{\Delta }\mu ^2$$
(50)
where $`\mathrm{\Delta }\mu ^2`$ is the loop correction. From the above one finds
$$\frac{\mu ^2}{c_r}=(t^21)^1(m_{\frac{1}{2}}^2g_r^{}t^2(m_{\frac{1}{2}}^2e_r^{}+A_0m_0m_{\frac{1}{2}}f_r^{}))+\frac{\mathrm{\Delta }\mu ^2}{c_r}$$
(51)
Here $`g_r^{}=\frac{\stackrel{~}{g}}{c_r}`$ where $`\stackrel{~}{g}=\stackrel{~}{\alpha }_G(\frac{3}{2}\stackrel{~}{f}_2+\frac{3}{10}\stackrel{~}{f}_1)`$, and $`f_r^{}=\frac{\stackrel{~}{f}}{c_r}`$. For large $`\mathrm{tan}\beta `$ Eq.(51) reduces down to
$$\frac{\mu ^2}{c_r}=(m_{\frac{1}{2}}^2e_r^{}+A_0m_0m_{\frac{1}{2}}f_r^{})+\frac{\mathrm{\Delta }\mu ^2}{c_r}$$
(52)
These results lead to Eq.(34).
Tables:
| Table 1: Nonuniversalities at $`M_X`$. | | | |
| --- | --- | --- | --- |
| SU(5) rep | $`n_1^r`$ | $`n_2^r`$ | $`n_3^r`$ |
| 1 | 1 | 1 | 1 |
| 24 | -1 | $`3`$ | $`2`$ |
| 75 | -5 | 3 | $`1`$ |
| 200 | 10 | 2 | 1 |
| Table 2: Uncertainties in $`\sigma _{\pi N}`$, $`\sigma _0`$ and x | | | | | |
| --- | --- | --- | --- | --- | --- |
| Ref. | $`\sigma _{\pi ,N}`$ (MeV) | Ref. | $`\sigma _0`$ (MeV) | Ref. | x |
| | $`45\pm 8`$ | | $`35\pm 5`$ | $`\overline{\sigma }_0`$/$`\overline{\sigma }_{\pi N}`$ | $`0.74\pm 0.25`$ |
| | $`48\pm 10`$ | | $`36\pm 7`$ | | $`0.64\pm 0.03`$ |
| | $`50\pm 10`$ | $`\overline{\sigma }_0`$ | $`35.5\pm 6`$ | | $`0.57\pm 0.1`$ |
| | $`49\pm 8`$ | | | $`\overline{x}`$ | $`0.67\pm 0.18`$ |
| $`\overline{\sigma }_{\pi N}`$ | $`48\pm 9`$ | | | | |
Figure captions
Fig.1: Exhibition of the dependence of $`\sigma _{\chi p}(scalar)`$ on the uncertainties in $`f_u^p`$, $`f_d^p`$ and $`f_s^p`$ for the universal case. The plots are for three sets of $`f_i^p`$ discussed in Sec.3. The parameter space spanned is as discussed in the text.
Fig.2: Plot of gaugino-Higgsino components $`X_{n0}`$ (n=1,2,3,4) as a function of $`c_{24}`$ for the input data $`m_0=51`$ GeV, $`m_\chi `$ =70 GeV, $`\mathrm{tan}\beta =10`$, $`A_t/m_0=7`$ where $`X_{20}=|X_{20}|`$ and $`X_{40}=|X_{40}|.`$
Fig.3: Plot of the maximum and the minimum curves $`\sigma _{\chi p}(scalar)`$ as a function of $`m_\chi `$ for the case when one considers admixtures of $`\mathrm{𝟏}+\mathrm{𝟐𝟒}`$ representations with $`c_{24}=0.1`$ (dashed), $`c_{24}=0`$ (solid) and $`c_{24}=0.08`$ (dotted) when the other parameters are varied over the assumed naturalness range as discussed in the text.
Fig.4: Plot of the maximum and the minimum curves $`\sigma _{\chi p}(scalar)`$ as a function of $`m_\chi `$ for the case when one considers admixtures of $`\mathrm{𝟏}+\mathrm{𝟕𝟓}`$ representations with $`c_{75}=0.06`$ (dashed), $`c_{75}=0`$ (solid) and $`c_{75}=0.04`$ (dotted) when the other parameters are varied over the assumed naturalness range as discussed in the text.
Fig.5: Plot of the maximum and the minimum curves $`\sigma _{\chi p}(scalar)`$ as a function of $`m_\chi `$ for the case when one considers admixtures of $`\mathrm{𝟏}+\mathrm{𝟐𝟎𝟎}`$ representations with $`c_{200}=0.1`$ (dashed), $`c_{200}=0`$ (solid) and $`c_{200}=0.08`$ (dotted) when the other parameters are varied over the assumed naturalness range as discussed in the text.
Fig.6: Plot of the maximum and the minimum curves for $`\sigma _{\chi p}(scalar)`$ as a function of $`m_\chi `$ for the heterotic string model O-II. The range of parameters consists of $`ϵ,ϵ^{}`$ in the range 0.0025-0.01, $`\theta `$ in the range 0.1-1.6, $`\delta _{GS}`$ in the range -1 to -10, $`m_{3/2}`$ in the range up to 2 TeV, and values of $`\mathrm{tan}\beta `$ range up to 25.
Fig.7: Plot of the maximum and the minimum curves for $`\sigma _{\chi p}(scalar)`$ as a function of $`m_\chi `$ for the $`95_1`$ D brane model. The range of parameters consists of $`\theta `$ in the range $`0.11.6`$, $`m_{3/2}`$ in the range up to 2 TeV, and values of $`\mathrm{tan}\beta `$ range up to 25.
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# 1 Preliminaries and review
## 1 Preliminaries and review
Let $`C`$ be a smooth complex projective curve of genus $`g`$, and let $`pC`$ be a distinguished point. For $`n0`$, denote by $`K(n)`$ the line bundle $`K𝒪(np)`$, where $`K`$ is the canonical bundle. A Higgs bundle with values in $`K(n)`$ is a pair $`(E,\varphi )`$ consisting of a holomorphic vector bundle $`E`$ over $`C`$ and a Higgs field $`\varphi H^0(\text{End}EK(n))`$. It is stable if for all proper subbundles $`FE`$ such that $`\varphi (F)FK(n)`$, $`\mathrm{deg}F/\text{rk}F<\mathrm{deg}E/\text{rk}E`$.
The existing results we shall need to recall are few and, except for the authors’ result stated as (1) below, have been known for some time. They can be summarized as follows.
First, there are some elementary facts about stable Higgs bundles \[14, 4.2,4.3\].
(1.1) If $`X`$ parametrizes two families $`(𝐄,𝚽)`$ and $`(𝐄^{},𝚽^{})`$ of stable Higgs bundles, and $`(𝐄_x^{},𝚽_x^{})(𝐄_x^{},𝚽_x^{})`$ for all $`xX`$, then $`𝐄^{}=𝐄L`$ for some line bundle $`L`$ over $`X`$, and $`𝚽^{}=𝚽`$. Moreover, if an action of $`^\times `$ on $`X`$ lifts to $`𝐄`$ and $`𝐄^{}`$ preserving the Higgs fields, then it lifts to $`L`$ also so that $`𝐄^{}=𝐄L`$ equivariantly.
Next, there is a moduli space, constructed by Simpson and Nitsure .
(1.2) There exists a moduli space $`_n`$ of stable Higgs bundles of rank 2 and degree 1 with values in $`K(n)`$, which is a smooth quasi-projective variety. It admits a universal family $`(𝐄,𝚽)`$, and the $`^\times `$-action on $`_n`$ given by $`\lambda (E\varphi )=(E,\lambda \varphi )`$ lifts to this family.
The following alternative interpretation of $`_0`$, due to Corlette , Donaldson , Hitchin , and Simpson , will only be used in §10.
(1.3) For $`n=0`$, the moduli space $`_0`$ is diffeomorphic to the moduli space $``$ of $`\mathrm{GL}(2,)`$-connections of constant central curvature $`i\omega I`$, where $`\omega `$ is a volume form on $`C`$. That is, $`_0\mu ^1(I)/\mathrm{GL}(2,)`$, where $`\mu :\mathrm{GL}(2,)^{2g}\mathrm{GL}(2,)`$ is given by $`\mu (A_j,B_j)=_{j=1}^gA_j^{}B_j^{}A_j^1B_j^1`$, and $`\mathrm{GL}(2,)`$ acts on $`\mu ^1(I)`$ by simultaneous conjugation. The natural determinant maps and universal families coincide under this diffeomorphism.
Let $`\mathrm{\Xi }`$ be a fixed holomorphic line bundle over $`C`$ of degree 1, and let $`_n_n`$ be the subspace consisting of those $`(E,\varphi )_n`$ such that $`\mathrm{\Lambda }^2E\mathrm{\Xi }`$ and $`\text{tr}\varphi =0`$. In the case $`n=0`$, this is the moduli space studied by Hitchin . The discussion so far, and the previous paper of the authors , refer to $`_n`$, but the remainder of this paper will actually work with $`_n`$. This gives equivalent information for the following reason. The group $`\mathrm{\Sigma }=_2^{2g}\text{Jac}C`$ of line bundles with structure group $`_2`$ acts on $`_n`$ by tensor product, and indeed $`_n=(_n\times T^{}\text{Jac}C)/\mathrm{\Sigma }`$. As seen in §1 of our previous paper , $`H^{}(_n)=H^{}(_n)^\mathrm{\Sigma }H^{}(\text{Jac}C)`$ as rings. To describe $`H^{}(_n)`$, therefore, it suffices to describe $`H^{}(_n)^\mathrm{\Sigma }`$. This will be the purpose of the paper. The part of $`H^{}(_n)`$ not invariant under $`\mathrm{\Sigma }`$ is ignored here, but it is completely described in a forthcoming work .
The main result of our paper on the generators is the following.
(1.4) The rational cohomology ring $`H^{}(_n)^\mathrm{\Sigma }`$ is generated by the universal classes, that is, the Künneth components $`\alpha _2`$, $`\beta _2`$, and $`\psi _{2,j}`$ of $`\overline{c}_2(𝐄)=\frac{1}{4}c_2(\text{End}𝐄)`$.
This result has been extended to higher rank Higgs bundles by Markman .
Following the conventions established by Newstead , we will let $`\alpha =\frac{1}{2}\alpha _2`$, $`\beta =\frac{1}{4}\beta _2`$, and $`\psi _j=\psi _{2,j}`$ for $`j=1,\mathrm{},2g`$, so that
$$c_2(\text{End}𝐄)=2\alpha \sigma \beta +4\underset{j=1}{\overset{g}{}}\psi _je_j,$$
(1.5)
where $`e_1,\mathrm{}e_{2g}`$ is the usual basis for $`H^1(C)`$, and $`\sigma H^2(C)`$ is the positive generator.
## 2 Statement of the main result
Our task, then, is to give a complete set of relations between the generators $`\alpha `$, $`\beta `$, and $`\psi _1,\mathrm{},\psi _{2g}`$. To do so, we must first say a little about the action of the symplectic group on $`H^{}(_n)^\mathrm{\Sigma }`$.
The group of orientation-preserving diffeomorphisms of $`C`$ acts on $`H^{}(C)`$ by automorphisms, so it has the automorphism group of $`H^{}(C)`$, namely $`\mathrm{\Gamma }=\mathrm{Sp}(2g,)`$, as a quotient \[19, p. 178\].
(2.1) There is a natural action of $`\mathrm{\Gamma }`$ on $`H^{}(_n)^\mathrm{\Sigma }`$ which fixes $`\alpha `$ and $`\beta `$ but acts on the $`\psi _j`$ as the standard representation.
Proof. In the case of $`_0`$, this follows immediately from (1), but to extend it to $`_n`$ we describe another argument.
Let $`f:CC`$ be any orientation-preserving diffeomorphism. The complex structure typically is not preserved by $`f`$, so pulling it back induces a new complex structure $`C^{}`$ on the same underlying surface. Since Teichmüller space, or the moduli space of curves, is connected, there is a path connecting $`C`$ to $`C^{}`$. The construction of the moduli space $`_n`$, and of the universal pair $`(𝐄,𝚽)`$, can be carried out simultaneously over all the Riemann surfaces in this path. Hence by homotopy invariance there is a topological isomorphism $`𝐄𝐄^{}`$, where $`(𝐄^{},𝚽^{})`$ is a universal pair on $`C^{}`$. The homotopy class of the isomorphism depends only on the isotopy class of $`f`$.
On the other hand, if $`\widehat{f}:_n^{}_n`$ is the map of moduli spaces induced by $`f:C^{}C`$, then $`(\widehat{f}\times f)^{}(𝐄,𝚽)`$ is a universal pair over $`C^{}`$, and so by the uniqueness in (1),
$$(\widehat{f}\times f)^{}\text{End}𝐄\text{End}𝐄^{}\text{End}𝐄.$$
Hence $`(\widehat{f}\times f)^{}c_2(\text{End}𝐄)=c_2(\text{End}𝐄)`$, so $`\widehat{f}^{}\alpha =\alpha `$, $`\widehat{f}^{}\beta =\beta `$, and $`(\widehat{f}\times f)^{}_j\psi _je_j=_j\psi _je_j`$. The action of the diffeomorphism group on $`H^3(_n)=\psi _j`$ is therefore dual to its action on $`H^1(C)=e_j`$; this factors through the standard representation of $`\mathrm{\Gamma }`$, which is self-dual. Moreover, by (1), the action of the diffeomorphism group on all of $`H^{}(_n)^\mathrm{\Sigma }`$ factors through $`\mathrm{\Gamma }`$. . $`\mathrm{}`$
The exterior square of the standard representation of $`\gamma `$ has an invariant element, the symplectic form. So $`\gamma =2_{j=1}^g\psi _j\psi _{j+g}H^6(_n)`$ is a $`\mathrm{\Gamma }`$-invariant element. Since the powers of the symplectic form are the only invariant elements of exterior powers of the standard representation, we deduce the following from (1).
(2.2) The $`\mathrm{\Gamma }`$-invariant part of $`H^{}(_n)^\mathrm{\Sigma }`$ is generated by $`\alpha H^2`$, $`\beta H^4`$, and $`\gamma H^6`$.
Like the exterior square discussed above, the higher exterior powers of the standard representation of $`\mathrm{\Gamma }`$ are reducible. Indeed, let $`\mathrm{\Lambda }^k(\psi )`$ be the $`k`$th exterior power of the standard representation, with basis $`\psi _1,\mathrm{},\psi _{2g}`$. Define the primitive part $`\mathrm{\Lambda }_0^k(\psi )`$ to be the kernel of the natural map $`\mathrm{\Lambda }^k\mathrm{\Lambda }^{2g+2k}`$ given by the wedge product with $`\gamma ^{g+1k}`$. The primitive part is complementary to $`\gamma \mathrm{\Lambda }^{k2}\mathrm{\Lambda }^k`$, and is an irreducible representation of $`\mathrm{\Gamma }`$: this is well-known for $`\mathrm{Sp}(2g,)`$, and so remains true for the Zariski dense subgroup $`\mathrm{\Gamma }`$. Being irreducible, it is generated by $`\psi _1\mathrm{}\psi _k`$.
For any $`g,n0`$, let $`I_n^g`$ be the ideal within the polynomial ring $`[\alpha ,\beta ,\gamma ]`$ generated by $`\gamma ^{g+1}`$ and the polynomials
$$\rho _{r,s,t}^c=\underset{i=0}{\overset{\mathrm{min}(c,r,s)}{}}(ci)!\frac{\alpha ^{ri}}{(ri)!}\frac{\beta ^{si}}{(si)!}\frac{(2\gamma )^{t+i}}{i!},$$
(2.3)
where $`c=r+3s+2t2g+2n`$, for all $`r,s,t0`$ such that
$$r+3s+3t>3g3+n\text{xxx}\text{ and }\text{xxx}r+2s+2t2g2+n.$$
(2.4)
The following is then the main result of the present paper.
(2.5) As a $`\mathrm{\Gamma }`$-algebra,
$$H^{}(_n)^\mathrm{\Sigma }=\underset{k=0}{\overset{g}{}}\mathrm{\Lambda }_0^k(\psi )[\alpha ,\beta ,\gamma ]/I_{n+k}^{gk}.$$
The theorem enunciated in (1.2) of our previous paper is the above in the case $`n=0`$. In that case the relation of lowest degree is $`\rho _{1,g1,0}^g=g\alpha \beta ^{g1}+(g1)\beta ^{g2}(2\gamma )`$. When $`n=1`$, there are two relations of lowest degree, one of which is $`\rho _{0,g,0}^{g+1}=(g+1)\beta ^g`$. When $`n2`$, the lowest degree in which a relation appears is $`2(2g2+n)`$. At least for $`sr`$, the relations in this degree have the particularly simple form $`\beta ^{sr}(\alpha \beta +2\gamma )^r`$.
## 3 Equivariant cohomology
Our main tool for studying the ring structure of $`H^{}(_n)`$ is equivariant cohomology, which we briefly review. For a more leisurely exposition see Atiyah-Bott .
If a group — say a Lie group — acts on a topological space $`M`$, the homotopy quotient $`M_G`$ is defined as the associated bundle over the classifying space $`BG`$ with fiber $`M`$:
$$M_G=\frac{M\times EG}{G}.$$
The equivariant cohomology of $`M`$ is defined to be simply the ordinary cohomology of $`M_G`$:
$$H_G^{}(X)=H^{}(X_G).$$
This is a module over $`H^{}(BG)`$. Restricting to any fiber gives a natural map $`H_G^{}(M)H^{}(M)`$. Note also that if $`G`$ acts trivially on $`M`$, then $`H_G^{}(M)=H^{}(M)H^{}(BG)`$.
If the action of $`G`$ lifts to a linear action on a vector bundle $`E`$, then a vector bundle $`E_G`$ over $`M_G`$ can be defined in the obvious way. Thus a vector bundle equipped with such a lifting possesses well-defined equivariant characteristic classes lying in $`H_G^{}(M)`$.
In our case the group acting is $`T=^\times `$, so that $`BT=^{\mathrm{}}`$, and $`H^{}(BT)=[u]`$ where $`u`$ is a class of degree 2. Kirwan proved the following fundamental results on $`^\times `$-actions.
(3.1) When $`T=^\times `$ acts algebraically on a smooth quasi-projective $`M`$ so that $`lim_{\lambda 0}\lambda x`$ exists for every $`x`$, then
* there is an additive isomorphism
$$H^i(M)\underset{d}{}H^{i+r_d}(F_d),$$
where $`=_dF_d`$ is the decomposition of the fixed-point set into components and $`r_d`$ is the dimension of the subbundle of $`TM|_{F_d}`$ acted on with negative weight by $`T`$;
* the Leray sequence of $`M_TBT`$ degenerates, so that $`H_T^{}(M)H^{}(M)H^{}(BT)`$ additively, and the ring homomorphism $`H_T^{}(M)H^{}(M)`$ is surjective;
* the restriction to the fixed-point set
$$H_T^{}(M)H_T^{}()=H^{}()[u]$$
is an injective ring homomorphism.
Statement (i) is perhaps most familiar in a symplectic context as stating that the moment map is a perfect Bott-Morse function. But statement (iii) is equally crucial for us since it respects the ring structure. Together with (ii), it will tell us that a polynomial in $`\alpha `$, $`\beta `$, and $`\psi _j`$ is a relation on $`_n`$ if and only if it is the value at $`u=0`$ of a polynomial in $`\alpha `$, $`\beta `$, $`\psi _j`$, and $`u`$ — the equivariant extension of the relation — whose restriction to $`H_G^{}(F_d)=H^{}(F_d)[u]`$ is a relation for each $`d`$.
## 4 Fixed points of the circle action
We will study the action of $`T=^\times `$ on $`_n`$ given simply by $`\lambda (E,\varphi )=(E,\lambda \varphi )`$. By (1) this lifts to the universal bundle, and hence the universal classes extend to equivariant classes, which by abuse of notation, we continue to denote $`\alpha `$, $`\beta `$, and $`\psi _j`$. They are canonical by the uniqueness in (1).
In light of (3), it is vital to determine the fixed-point set for this action. As discussed in (10.5) of our previous paper , this would be somewhat tricky in arbitrary rank. But now that we have restricted attention to rank 2 (and fixed determinant), it is not so hard. The lemma below is proved by Hitchin \[15, 7.1\] for $`_0`$, but his proof generalizes directly to $`_n`$.
(4.1) The components of the fixed-point set $``$ for the $`T`$-action on $`_n`$ are as follows.
* A component $`F_0`$ isomorphic to $`𝒩`$, the moduli space of stable bundles $`E`$ with $`\mathrm{\Lambda }^2E\mathrm{\Xi }`$. It parametrizes Higgs bundles of the form $`(E,0)`$.
* Components $`F_1,\mathrm{},F_{g+\left[\frac{n1}{2}\right]}`$ which are fibered products
$$F_d=C_{2g+n12d}\times _{\text{Jac}^{2d}C}\text{Jac}^dC,$$
where the maps $`C_{2g+n12d}\text{Jac}^{2d}C`$ and $`\text{Jac}^dC\text{Jac}^{2d}C`$ are given by $`DK\mathrm{\Xi }(n)(D)`$ and $`LL^2`$ respectively. These parametrize Higgs bundles $`(E,\varphi )`$ of the form $`E=L\mathrm{\Xi }L^1`$, $`\varphi =\left(\begin{array}{cc}0& 0\\ s& 0\end{array}\right),`$ where $`s`$ is the section of $`KL^2\mathrm{\Xi }(n)`$ vanishing at $`D`$.
Hitchin went on to compute the cohomology of the fixed components of type (ii) as follows. By the Leray sequence
$$H^{}(F_d)=\underset{i\mathrm{\Sigma }}{}H^{}(C_{2g+n12d},_i),$$
(4.2)
where the right-hand side consists of cohomology with local coefficients, and $`_i`$ runs over the flat line bundles with structure group $`_2`$ pulled back from $`\text{Jac}^{2d}C`$. If $`_i`$ is the trivial bundle this is simply the ordinary cohomology $`H^{}(C_{2g+n12d})`$. Otherwise, $`H^k(C_{2g+n12d},_i)=\mathrm{\Lambda }^kH^1(C,L_i)`$ if $`k=2g+n12d`$, and 0 if not. Here $`L_i`$ runs over the flat line bundles with structure group $`_2`$ pulled back from $`\text{Jac}^1C`$ to $`C`$ by the Abel-Jacobi map. Hitchin shows that for $`L_i`$ non-trivial, $`H^0(C,L_i)=H^2(C,L_i)=0`$, and $`H^1(C,L_i)`$ has dimension $`2g2`$.
The action of $`\mathrm{\Sigma }`$ on $`_n`$ commutes with the $`T`$-action, and induces the trivial action on $`H^{}(𝒩)`$ since it is generated by universal classes . It acts on the remaining $`F_d`$ as the Galois group of the unbranched cover $`F_dC_{2g+n12d}`$, and the splitting (4.2) is exactly the decomposition of the cohomology into weight spaces.
Consequently, the $`\mathrm{\Sigma }`$-invariant part of $`H^{}()`$ is
$$H^{}(𝒩)\underset{d=1}{\overset{g+\left[\frac{n1}{2}\right]}{}}H^{}(C_{2g+n12d}).$$
## 5 Symmetric products of a curve
The symmetric products of the curve $`C`$ thus enter into our considerations. So let us review some facts about the cohomology of such a symmetric product. Good references are the paper of Macdonald or the book of Arbarello et al. .
In $`C_m\times C`$, there is a universal divisor $`\mathrm{\Delta }`$ such that $`\mathrm{\Delta }(\{D\}\times C)=D`$. Write its Poincaré dual in terms of Künneth components as
$$m\sigma +\eta +\underset{j=1}{\overset{2g}{}}\xi _je_j,$$
where $`\sigma `$ and $`e_1,\mathrm{},e_{2g}`$ are generators of $`H^2(C)`$ and $`H^1(C)`$ respectively, so that $`\eta H^2(C_m)`$ and $`\xi _1,\mathrm{},\xi _{2g}H^1(C_m)`$. A theorem of Macdonald asserts that the cohomology ring $`H^{}(C_m)`$ is generated by $`\eta `$ and the $`\xi _j`$. It is convenient to introduce $`\theta _j=\xi _j\xi _{j+g}`$ and $`\theta =_{j=1}^g\theta _jH^2(C)`$.
The group of orientation-preserving diffeomorphisms of $`C`$ acts on $`C_m\times C`$. It preserves $`\mathrm{\Delta }`$, and hence the Künneth components of its Poincaré dual. Hence it leaves $`\eta `$ invariant. Moreover, its action on the linear span of the $`\xi _j`$, which is $`H^1(C_m)`$, is dual to its action on $`H^1(C)`$, and hence factors through the quotient $`\mathrm{\Gamma }=\mathrm{Sp}(2g,)`$. There is therefore a surjective homomorphism of $`\mathrm{\Gamma }`$-algebras $`\mathrm{\Lambda }^{}(\xi )[\eta ]H^{}(C_m)`$. Here $`\mathrm{\Lambda }^{}(\xi )`$ denotes the exterior algebra of the standard $`2g`$-dimensional representation of $`\mathrm{\Gamma }`$, with basis vectors $`\xi _1,\mathrm{},\xi _{2g}`$.
The class $`\theta `$ represents the symplectic form. So in terms of the primitive parts $`\mathrm{\Lambda }_0^k`$ introduced in §2, the surjective homomorphism above is better expressed as
$$\underset{k=0}{\overset{g}{}}\mathrm{\Lambda }_0^k(\xi )[\eta ,\theta ]H^{}(C_m).$$
In particular, the $`\mathrm{\Gamma }`$-invariant part of $`H^{}(C_m)`$ is generated by $`\eta `$ and $`\theta `$.
The following result on $`H^{}(C_m)`$ will be of key importance for us. Note that we use the term total degree to mean half the ordinary degree of a cohomology class.
(5.1) Let $`l`$, $`m`$, $`p`$, and $`q`$ be non-negative integers. If $`mg+ql`$ and $`g+pq<l`$, then
$$(\frac{\eta ^p\mathrm{exp}\theta }{(1+\eta )^q})_l=0$$
in $`H^{}(C_m)`$, where the subscript $`l`$ denotes the part in total degree $`l`$.
Proof. Since the cup product is a homomorphism of $`\mathrm{\Gamma }`$-modules, Poincaré duality holds for the $`\mathrm{\Gamma }`$-invariant part. It therefore suffices to check that the product of this expression with any monomial in $`\eta `$ and $`\theta `$ evaluates to 0 on the fundamental class of $`C_m`$.
It follows from Macdonald’s results that any monomial $`\eta ^v_j\xi _j^{w_j}`$ of total degree $`m`$ evaluates on the fundamental class of $`C_m`$ to 1 if $`w_j=w_{j+g}1`$ for each $`jg`$, and 0 otherwise. As pointed out by Zagier , this implies that for any formal power series $`A(x)`$ and $`B(x)`$,
$$A(\eta )\mathrm{exp}(\theta B(\eta ))[C_m]=\underset{\eta =0}{\text{Res}}\frac{A(\eta )(1+\eta B(\eta ))^gd\eta }{\eta ^{m+1}}.$$
We multiply our expression by the generating function $`\mathrm{exp}(s\theta )/(1+t\eta )`$ for the monomials in $`\eta `$ and $`\theta `$, and ask the coefficient of $`s^it^j`$ to vanish whenever $`i+j=ml`$:
$`\underset{s^it^j}{\text{Coeff}}{\displaystyle \frac{\eta ^p\mathrm{exp}((s+1)\theta )}{(1+\eta )^q(1+t\eta )}}[C_m]`$ $`=`$ $`\underset{s^it^j}{\text{Coeff}}\underset{\eta =0}{\text{Res}}{\displaystyle \frac{\eta ^p(1+\eta +s\eta )^gd\eta }{(1+\eta )^q(1+t\eta )\eta ^{m+1}}}`$
$`=`$ $`\text{const.}\underset{\eta =0}{\text{Res}}{\displaystyle \frac{\eta ^p(1+\eta )^{gi}\eta ^i\eta ^jd\eta }{(1+\eta )^q\eta ^{m+1}}}`$
$`=`$ $`\text{const.}\underset{\eta =0}{\text{Res}}\eta ^{i+j+pm1}(1+\eta )^{giq}d\eta .`$
Now since $`giqgq(ml)0`$ by hypothesis, the second factor is a polynomial of degree $`giq`$. All terms therefore have degree at most
$$(i+j+pm1)+(giq)p+(ml)m1+gq=pl1+gq,$$
which is less than $`1`$ by hypothesis. . $`\mathrm{}`$
## 6 Restriction of the universal classes to the fixed-point set
In order to apply (3)(iii), we need to know how the equivariant universal classes restrict to each component of the fixed-point set. The lowest component $`F_0=𝒩`$ is easy. The universal pair over $`\times C`$ restricts to a universal bundle over $`𝒩\times C`$, and the $`T`$-action restricts to a trivial action. So $`\alpha `$, $`\beta `$ and $`\psi _j`$ restrict to classes on $`𝒩`$ defined in a like manner, and bearing the same names. The relations between these classes on $`𝒩`$ have been studied by many authors, notably Zagier ; we will have occasion to use his results later.
The components $`F_d`$ for $`d>0`$ are handled by the following lemma.
(6.1) For $`d>0`$, the restrictions of the universal classes to $`F_d`$ are pulled back from the symmetric product $`C_{2g+n12d}`$; indeed,
* $`\alpha |_{F_d}=(2d1)(\eta u)+\theta `$;
* $`\beta |_{F_d}=(\eta u)^2`$;
* $`\psi _j|_{F_d}=\frac{1}{2}(\eta u)\xi _j`$;
* $`\gamma |_{F_d}=\frac{1}{2}(\eta u)^2\theta `$.
Proof. We first construct an equivariant universal family $`(𝐄,𝚽)`$ of Higgs bundles over $`F_d\times C`$. Since $`\text{End}𝐄`$ is unique as an equivariant bundle by (1), the universal classes must restrict to the Künneth components of its second Chern class.
Take the following three ingredients. First, the line bundle $`K\mathrm{\Xi }(n)`$ over $`C`$. Second, the universal divisor $`\mathrm{\Delta }C_{2g+n12d}\times C`$, or rather its associated line bundle $`𝒪(\mathrm{\Delta })`$. Third, any Poincaré line bundle $``$ over $`\text{Jac}^dC\times C`$. Now pull all three back to $`F_d\times C`$. By the definition of the fibered product, $`^2`$ and $`K\mathrm{\Xi }(n)(\mathrm{\Delta })`$ are isomorphic when restricted to any fiber of the projection $`F_d\times CF_d`$. So by the push-pull formula, $`^2K^1\mathrm{\Xi }^1(n)(\mathrm{\Delta })`$ is the pull-back from $`F_d`$ of a line bundle, say $`M`$. There is then an element $`sH^0(F_d\times C,M^2K\mathrm{\Xi }(n))`$ vanishing precisely on the inverse image of $`\mathrm{\Delta }`$.
Let $`𝐄=M\mathrm{\Xi }^1`$, and let $`𝚽H^0(\text{End}𝐄K(n))`$ be given by $`𝚽=\left(\begin{array}{cc}0& 0\\ s& 0\end{array}\right)`$ with respect to the splitting. Then $`(𝐄,𝚽)`$ parametrizes the pairs of the form described in (4)(ii). It is hence a universal family. Moreover, if $`T`$ acts on the two factors with weights $`1`$ and $`0`$ respectively, then it acts on $`𝚽`$ by scalar multiplication. By (1), then, $`\text{End}𝐄`$ is equivariantly isomorphic to the restriction of its counterpart from $`_n\times C`$.
Since $`𝐄`$ splits as a direct sum, $`c_2(\text{End}𝐄)=(c_1(M\mathrm{\Xi }^1)c_1())^2=c_1(M\mathrm{\Xi }^2)^2=c_1(K^1(n)(\mathrm{\Delta }))^2`$. For this to be correct equivariantly, we must include the weights of the $`T`$-action, so the equivariant $`c_1`$ is the non-equivariant $`c_1`$ minus $`u`$. It is well-known that
$$c_1(𝒪(\mathrm{\Delta }))=(2g+n12d)\sigma +\eta +\underset{j=1}{\overset{2g}{}}\xi _je_j:$$
see for example Arbarello et al. . Hence
$$c_2(\text{End}𝐄)=((22gn)\sigma +(2g+n12d)\sigma +\eta +\underset{j=1}{\overset{2g}{}}\xi _je_ju)^2.$$
Using the identity $`\left(\xi _je_j\right)^2=2\theta \sigma `$ and comparing coefficients with those of (1.5) yields the result. . $`\mathrm{}`$
All our weapons are now prepared, and we are ready to attack the main result. It is not a frontal assault, however. Rather, we begin by computing some relations quite different from the $`\rho `$-classes.
## 7 Some recursively defined polynomials in $`𝜶`$, $`𝜷`$, $`𝜸`$ and $`𝒖`$
The $`\mathrm{\Gamma }`$-invariant subring is at the heart of the larger ring containing it; its structure is the key to that of the whole. Our strategy will therefore be to look first for relations between $`\alpha `$, $`\beta `$, and $`\gamma `$. The method is curiously roundabout. First, certain complicated polynomials, defined recursively here in §7, are shown in §8 to be relations, by writing down their equivariant extensions explicitly. Then they are shown to be expressible in terms of the much simpler polynomials $`\rho _{r,s,t}^c`$ by a purely algebraic argument, given in §9. Not until §12 does a dimension count show that the $`\rho `$-classes must all be relations.
One relation which holds automatically in all $`_n`$ is $`\gamma ^{g+1}=0`$. This is simply because $`\gamma =2_{j=1}^g\psi _j\psi _{j+g}`$ and each $`\psi _j^2=0`$ by skew-commutativity. It is therefore convenient to view our polynomials as belonging to the ring $`R=[\alpha ,\beta ,\gamma ]/(\gamma ^{g+1})`$.
Define polynomials $`\xi _r^k`$ by $`\xi _r^k=0`$ for $`r<0`$, $`\xi _0^k=1`$, and
$$(r+1)\xi _{r+1}^k=\alpha \xi _r^k+(r2k)\beta \xi _{r1}^k+2\gamma \xi _{r2}^k$$
(7.1)
for $`r>0`$. When $`k=0`$, these are the polynomials $`\xi _r`$ defined by Zagier , and his generating function for the $`\xi _r`$ extends readily.
(7.2) Define $`F_0^k(x)=_{r=0}^{\mathrm{}}\xi _r^kx^rR[[x]]`$. Then
$$F_0^k(x)=(1\beta x^2)^{(2k1)/2}e^{2\gamma x/\beta }\left(\frac{1+x\sqrt{\beta }}{1x\sqrt{\beta }}\right)^{(\alpha \beta +2\gamma )/2\beta \sqrt{\beta }}.$$
Proof. From Proposition 4 of Zagier we know that $`F_0^0`$ satisfies the differential equation
$$(1\beta x^2)(F_0^0)^{}(x)=(\alpha +\beta x+2\gamma x^2)F_0^0(x).$$
Now (7.1) is equivalent to
$$(r+1)\xi _{r+1}^k(r1)\beta \xi _{r1}^k=\alpha \xi _r^k+(12k)\beta \xi _{r1}+2\gamma \xi _{r2},$$
which shows that $`F_0^k`$ satisfies the differential equation
$$(1\beta x^2)(F_0^k)^{}(x)=(\alpha +(12k)\beta x+2\gamma x^2)F_0^k(x),$$
with initial condition $`F_0^k(0)=1`$. But $`(1\beta x^2)^kF_0^0`$ satisfies the same differential equation
$`(1\beta x^2)\left((1\beta x^2)^kF_0^0\right)^{}`$ $`=`$ $`(1\beta x^2)^k(1\beta x^2)(F_0^0)^{}k\beta x(1\beta x^2)^kF_0^0`$
$`=`$ $`(\alpha +\beta x+2\gamma x^2)(1\beta x^2)^kF_0^02k\beta x(1\beta x^2)^kF_0^0`$
$`=`$ $`(\alpha +(12k)\beta x+2\gamma x^2)(1\beta x^2)^kF_0^0,`$
and certainly $`(1\beta 0^2)^kF_0^0(0)=1`$, so we conclude that
$$F_0^k(x)=(1\beta x^2)^kF_0^0(x).$$
Now substitute Zagier’s generating function for $`F_0^0`$. . $`\mathrm{}`$
(7.3) The polynomial $`\xi _r^k`$ is a relation on $`𝒩`$ whenever $`rg+2k`$.
Proof. An equivalent form of (7) is $`\xi _r^k=_{i=0}^k\xi _{r2i}^0(\beta )^i`$, and Zagier shows that $`\xi _r^0`$ is a relation on $`𝒩`$ for $`rg`$. . $`\mathrm{}`$
Now define an expression with one more index:
$$\xi _{r,s}^k=\underset{i=0}{\overset{s}{}}\left(\genfrac{}{}{0pt}{}{r2k+si}{r2k}\right)\beta ^{si}\frac{(2\gamma )^i}{i!}\xi _{ri}^k.$$
(7.4)
Note that this is $`0`$ if $`r<2k`$. Moreover, the $`i`$th term in the sum vanishes when $`i>r`$, and also (as an element of $`R`$) when $`i>g`$. Hence in particular $`\xi _{2k,g+l}^k=\beta ^l\xi _{2k,g}^k`$ in $`R`$ for all $`l0`$.
(7.5) Let $`F^k(x,y)=_{r,s=0}^{\mathrm{}}\xi _{r,s}^kx^ry^sR[[x,y]]`$. Then
$$F^k(x,y)=\left((1\beta y)^2\beta x^2\right)^{(2k1)/2}e^{2\gamma x/\beta }\left(\frac{1+x\sqrt{\beta }\beta y}{1x\sqrt{\beta }\beta y}\right)^{(\alpha \beta +2\gamma )/2\beta \sqrt{\beta }}.$$
Proof. For fixed $`r`$ we have
$`{\displaystyle \underset{s0}{}}\xi _{r,s}^ky^s`$ $`=`$ $`{\displaystyle \underset{i=0}{\overset{r}{}}}{\displaystyle \frac{(2\gamma y)^i}{i!}}\xi _{ri}^k{\displaystyle \underset{s=i}{\overset{\mathrm{}}{}}}\left(\genfrac{}{}{0pt}{}{r2k+si}{r2k}\right)(\beta y)^{si}`$
$`=`$ $`{\displaystyle \frac{1}{(1\beta y)^{r2k+1}}}{\displaystyle \underset{i=0}{\overset{r}{}}}{\displaystyle \frac{(2\gamma y)^i}{i!}}\xi _{ri}^k.`$
Multiplying by $`x^r`$ and summing over $`r0`$ we obtain
$$F^k(x,y)=(1\beta y)^{2k1}e^{2\gamma xy/(1\beta y)}F_0^k(\frac{x}{1\beta y}),$$
and the desired result follows by substituting the formula given in (7). . $`\mathrm{}`$
(7.6) For $`r,s0`$ we have
$$\xi _{r,s}^k=\underset{l=0}{\overset{s}{}}(1)^{sl}\left[\left(\genfrac{}{}{0pt}{}{r+l}{r}\right)+\left(\genfrac{}{}{0pt}{}{r+l1}{r}\right)\right]\xi _{sl}^k\xi _{r+s+l}^k,$$
where the binomial coefficient $`\left(\genfrac{}{}{0pt}{}{1}{0}\right)`$ is to be taken as 0.
Proof. Similar to Zagier’s Proof 1 of his Theorem 4. . $`\mathrm{}`$
## 8 Proof that the recursively defined polynomials are relations
After this algebraic preparation, we now find some relations between $`\alpha `$, $`\beta `$, and $`\gamma `$ that can be expressed in terms of the classes $`\xi _{r,s}^k`$ introduced above. We make fundamental use of Kirwan’s theorem (3)(iii), which tells us that any polynomial in the generators that vanishes on the equivariant cohomology of each component of the fixed-point set must be a relation.
(8.1) For $`n0`$, let $`pR[u]`$ be an equivariant relation on $`_{n+2}`$, that is, an element of the kernel of the natural map to $`H_T^{}(_{n+2})`$. Then $`p/u`$ is an equivariant relation on $`_n`$.
Proof. By Kirwan’s theorem (3)(iii), it suffices to show that $`p/u`$ restricts to an equivariant relation on each component of the fixed-point set of $`_n`$.
For $`F_0=𝒩`$, this is obvious, since $`p`$ is also a relation on $`H_T^{}(𝒩)=H^{}(𝒩)[u]`$.
As for $`F_d`$ with $`d>0`$, we may work in $`H^{}(C_m)`$, where $`m=2g+n12d`$. The relation $`p`$ restricts to a relation in $`H^{}(C_{m+2})[u]`$; moreover, since by (6) the restrictions of $`\alpha `$, $`\beta `$, $`\gamma `$ are polynomials in $`\eta u`$ and $`\theta `$, the relation can be expressed as a polynomial $`r(\varphi ,\theta ,u)|_{\varphi =\eta u}`$ such that $`r/u(\varphi ,\theta ,u)|_{\varphi =\eta u}`$ is the restriction of $`p/u`$, which we want to vanish.
The assignment $`u\eta \varphi `$ induces an isomorphism $`H^{}(C_{m+2})[u]H^{}(C_{m+2})[\varphi ]`$, so $`r(\varphi ,\theta ,\eta \varphi )`$ is a relation in $`H^{}(C_{m+2})[\varphi ]`$. Observe now that the derivative with respect to $`\eta `$ of any relation in $`H^{}(C_{m+2})`$ is a relation in $`H^{}(C_m)`$. This follows directly from the list of relations given by Macdonald \[20, 6.21\]. Therefore $`/\eta \text{(}r(\varphi ,\theta ,\eta \varphi )\text{)}`$ is a relation in $`H^{}(C_m)[\varphi ]`$, and hence
$$\frac{}{\eta }\text{(}r(\varphi ,\theta ,\eta \varphi )\text{)}|_{\varphi =\eta u}=\frac{r}{u}(\varphi ,\theta ,u)|_{\varphi =\eta u}$$
is a relation in $`H^{}(C_m)`$, as desired. . $`\mathrm{}`$
(8.2) For $`n0`$, let $`pR[u]`$ be an equivariant relation on $`_n`$. Then $`(u^2\beta )p`$ is an equivariant relation on $`_{n+1}`$.
Proof. Let $`F_d`$ be any component of the fixed-point set of $`_{n+1}`$. We show that restricting $`(u^2\beta )p`$ to $`F_d`$ yields $`0`$. It is clearly true on $`F_0=𝒩`$ since this is contained in $`_n`$. For $`d>0`$, the restriction of $`u^2\beta `$ to $`F_d`$ is $`\eta (2u\eta )`$ by (6). On the other hand $`p`$ restricted to $`F_d_n`$ is supposed to be zero. But the image of the inclusion $`F_d_nF_d`$ is Poincaré dual to $`\eta `$, since it is an étale cover of the inclusion $`C_{2g+n12d}C_{2g2d+n}`$. Hence $`\eta `$ times a relation on $`F_d_n`$ is a relation on $`F_d`$. . $`\mathrm{}`$
These results suggest that, even if we are interested only in the relations on $`_0`$, it is useful to study $`_n`$ for all $`n`$.
(8.3) For $`n0`$ and $`k=0,\mathrm{},[n/2]`$, the equivariant class
$$F^k(u,1)_{2g+2n}=\underset{r=0}{\overset{g+n}{}}\xi _{r,g+nr}^ku^r$$
is an equivariant relation on $`_{n+2}`$.
Proof. By Kirwan’s theorem (3)(iii), it suffices to show that it restricts to a relation on each component of the fixed-point set of $`_{n+2}`$.
For the first component $`F_0`$, namely $`𝒩`$, this follows immediately from (7) and (7).
For the remaining components $`F_d`$ with $`d>0`$, use (6) to restrict (7) to $`F_d`$. This yields
$$F^k(u,1)_{2g+2n}=\left(e^{\theta u}\frac{(1(\eta u)(\eta 2u))^{d+k1}}{(1\eta (\eta u))^{dk}}\right)_{2g+2n},$$
where the subscript, as in the past, denotes the part in total degree $`2g+2n`$, that is, in ordinary degree $`2(2g+2n)`$. To show that this vanishes in $`H^{}(C_{2g2d+1+n})H^{}(F_d)`$, express it as
$`\left(e^{\theta u}{\displaystyle \frac{(1(\eta u)(\eta 2u))^{d+k1}}{(1+\eta u)^{dk}\left(1\frac{\eta ^2}{1+\eta u}\right)^{dk}}}\right)_{2g+2n}`$
$`=`$ $`\left({\displaystyle \underset{i=1}{\overset{\mathrm{}}{}}}\left(\genfrac{}{}{0pt}{}{dk+i}{i}\right){\displaystyle \frac{\eta ^{2i}e^{\theta u}}{(1+\eta u)^{dk+i}}}(1(\eta u)(\eta 2u))^{d+k1}\right)_{2g+2n}.`$
It follows immediately from (5) that
$$\left(\frac{e^{\theta u}(\eta u)^{2i}}{(1+\eta u)^{dk+i}}\right)_{2(gd+n+i+1k)+j}=0$$
for $`j0`$ (the 2 appearing in the subscript since $`\eta u`$ and $`\theta u`$ are substituted for $`\eta `$ and $`\theta `$), and hence that
$$\left(\frac{e^{\theta u}\eta ^{2i}}{(1+\eta u)^{dk+i}}\right)_{2(gd+n+1k)+j}=0$$
for $`n0`$. Consequently each term in the sum above vanishes in total degree $`2g+2n`$. . $`\mathrm{}`$
* (8.4) For even $`n0`$ and $`k=0,\mathrm{},n/2`$, the equivariant class
$$\left((2+u^2\beta )^{n/2k}F^k(u,1)\right)_{2g+n+2k}$$
is an equivariant relation on $`_{n+2}`$ divisible by $`u^{2k}`$.
* For odd $`n0`$ and $`k=0,\mathrm{},(n1)/2`$, the equivariant class
$$\left((1+u^2\beta )(2+u^2\beta )^{(n1)/2k}F^k(u,1)\right)_{2g+n+2k+1}$$
is an equivariant relation on $`_{n+2}`$ divisible by $`u^{2k+1}`$.
Proof. Since $`u`$ is not a zero-divisor in $`H_T^{}(_{n+2})`$, to show the expression in (a) is a relation it suffices to do the same for the part in total degree $`2g+2n`$ of
$`u^{n2k}(2+u^2\beta )^{n/2k}F^k(u,1)`$
$`=`$ $`\left((1+u^2\beta )^2((1\beta )^2\beta u^2)\right)^{n/2k}F^k(u,1)`$
$`=`$ $`{\displaystyle \underset{i}{}}\left(\genfrac{}{}{0pt}{}{\frac{n}{2}k}{i}\right)(1+u^2\beta )^{2i}\left((1\beta )^2\beta u^2\right)^{n/2ki}F^k(u,1)`$
$`=`$ $`{\displaystyle \underset{i}{}}\left(\genfrac{}{}{0pt}{}{\frac{n}{2}k}{i}\right)(1+u^2\beta )^{2i}F^{n/2i}(u,1)`$
$`=`$ $`{\displaystyle \underset{i,j}{}}\left(\genfrac{}{}{0pt}{}{\frac{n}{2}k}{i}\right)\left(\genfrac{}{}{0pt}{}{2i}{j}\right)(u^2\beta )^jF^{n/2i}(u,1).`$
By (8), $`F^{n/2i}(u,1)_{2g+2n2j}`$ is a relation on $`_{nj+2}`$; hence by (8),
$$\left((u^2\beta )^jF^{n/2i}(u,1)\right)_{2g+2n}$$
is a relation on $`_{n+2}`$.
The statement about divisibility is easy, since $`F^k(u,1)=\xi _{r,s}^ku^r`$ and $`\xi _{r,s}^k=0`$ for $`r<2k`$.
The proof of (b) is similar: first multiply by $`u^{n12k}`$, compute as before
$`u^{n12k}(1+u^2\beta )(2+u^2\beta )^{(n1)/2k}F^k(u,1)`$
$`=`$ $`{\displaystyle \underset{i,j}{}}\left(\genfrac{}{}{0pt}{}{\frac{n1}{2}k}{i}\right)\left(\genfrac{}{}{0pt}{}{2i+1}{j}\right)(u^2\beta )^jF^{(n1)/2i}(u,1),`$
then apply (8) and (8).
As in (a), this is clearly divisible by $`u^{2k}`$, but the quotient is further divisible by another factor of $`u`$. This is because the coefficient of $`u^0`$ in the quotient is
$`\left((1\beta )(2\beta )^{(n1)/2k}{\displaystyle \underset{s=0}{\overset{\mathrm{}}{}}}\xi _{2k,s}^k\right)_{2g+n+1}`$
$`=`$ $`\left((1\beta )(2\beta )^{(n1)/2k}{\displaystyle \underset{s=g}{\overset{\mathrm{}}{}}}\xi _{2k,s}^k\right)_{2g+n+1}`$
$`=`$ $`\left((1\beta )(2\beta )^{(n1)/2k}\xi _{2k,g}^k{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}\beta ^l\right)_{2g+n+1}`$
$`=`$ $`\left((2\beta )^{(n1)/2k}\xi _{2k,g}^k\right)_{2g+n+1}`$
$`=`$ $`0`$
using $`\xi _{2k,g+l}^k=\beta ^l\xi _{2k,g}^k`$. . $`\mathrm{}`$
(8.5) For $`n2`$ even (resp. odd), $`\xi _{r,s}^k`$ (resp. $`\xi _{r,s}^k\beta \xi _{r,s1}^k`$) is a relation in the ordinary cohomology of $`_{n+2}`$
* for $`k=[n/2]i`$, $`r=n2i`$, and $`s=g+i`$, where $`i=0,\mathrm{},[n/2]`$;
* for $`k=[n/2]+j`$, $`r=n+3j`$, and $`s=gj`$, where $`j=1,\mathrm{},g`$.
Proof. Suppose first that $`n`$ is even. For (i), just take the formula from (1)(a) with $`k=n/2i`$, plug in $`F^k(u,1)=\xi _{r,s}^ku^r`$, divide by $`u^{2k}`$ and set $`u=0`$. Then compute using the definition of $`\xi _{r,s}^k`$, the binomial theorem, and $`\gamma ^{g+1}=0`$. For (ii), take the same formula on $`_{n+2j}`$ with $`k=n/2+j`$, apply (8) $`j`$ times, and proceed as in (i). Now suppose that $`n`$ is odd. It suffices to prove the same statement for the class $`\xi _{r,s}^k\beta \xi _{r,s1}^k+\xi _{r2,g+1}^k`$, because in case (i) the last term vanishes altogether, and in case (ii) it was shown in the even case to be a relation on $`_{n+3}_{n+2}`$. Then everything is similar to the even case. . $`\mathrm{}`$
## 9 Expressing the $`𝝃`$-classes in terms of the $`𝝆`$-classes
We now have many relations on $`_n`$. We cannot yet show that the simple polynomials $`\rho _{r,s,t}^c`$ of the main theorem are relations, but at least we can show that the relations we do have are linear combinations of them. Hence the goal of this section is to prove the following purely algebraic result.
(9.1) (a) For $`r2k`$, $`\xi _{r,s}^k`$ is a linear combination of those $`\rho _{u,v,w}^{r2k+vw}`$ with $`wr2k`$ and $`u+3wr`$. (b) For $`r2k+1`$, $`\xi _{r,s}^k\beta \xi _{r,s1}^k`$ is a linear combination of those $`\rho _{u,v,w}^{r2k+1+vw}`$ with $`wr2k+1`$ and $`u+3wr`$.
It is an easy matter to check, using high-school algebra and the equality of total degrees $`r+2s+3t=u+2v+3w`$, that when $`n`$, $`k`$, $`r`$ and $`s`$ are as in (8), the conditions (2.4) of membership in $`I_n^g`$ are satisfied by the $`\rho `$-classes named above. Hence the relations of (8) belong to $`I_n^g`$.
The proof of (9) will use a generating function for the $`\xi _r^k`$ which generalizes a formula for the $`\xi _r`$ stated without proof in Zagier’s paper . Zagier kindly communicated a proof to us, and it goes through almost verbatim for the generalization.
(9.2) If
$$\varphi _m^k(r,p)=\underset{x^m}{\text{Coeff}}\frac{1}{\mathrm{cosh}^{2k}\sqrt{3x}}\frac{\sqrt{3x}}{\mathrm{sinh}\sqrt{3x}}(\frac{\sqrt{3x}}{\mathrm{tanh}\sqrt{3x}})^r(\frac{1}{x}\frac{\mathrm{tanh}\sqrt{3x}}{x\sqrt{3x}})^p,$$
then for $`r0`$,
$$\xi _r^k=\underset{m,p}{}\frac{\varphi _m^k(r,p)}{3^{m+p}(r2m3p)!p!}\alpha ^{r2m3p}\beta ^m(2\gamma )^p.$$
Proof. The formula for $`\varphi _m^k(r,p)`$ may be abbreviated as $`\text{Coeff}_{x^m}A(x)B(x)^rC(x)^p`$. This directly gives a generating function for these numbers with $`r`$ and $`p`$ fixed, but to compute
$$F_0^k(t)=\underset{r}{}\xi _r^kt^r$$
we need instead a generating function for $`\varphi _m^k(l+2m+3p,p)`$ with $`l`$ and $`p`$ fixed and $`m`$ variable. The passage from one to the other, as usual, is by residue calculus: write $`\varphi _m^k(r,p)`$ as $`\text{Res}_{x=0}(A(x)B(x)^rC(x)^pdx/x^{m+1})`$ and change variables to $`y=x/B(x)^2=(1/3)\mathrm{tanh}^2(\sqrt{3x})`$ to get
$$\varphi _m^k(l+2m+3p,p)=\underset{y=0}{\text{Res}}(a(y)b(y)^lc(y)^pdy/y^{m+1})$$
with $`a(y)=A(x(y))x^{}(y)/B(y)^2`$, $`b(y)=B(x(y))`$, $`c(y)=C(x(y))B(x(y))^3`$. In other words,
$$\underset{m}{}\varphi _m^k(l+2m+3p,p)y^m=a(y)b(y)^lc(y)^p.$$
Then we need to verify
$`F_0^k(t)`$ $`=`$ $`{\displaystyle \underset{l,m,p0}{}}\varphi _m^k(l+2m+3p,p){\displaystyle \frac{(\alpha t)^l}{l!}}(\beta t^2/3)^m{\displaystyle \frac{(2\gamma t^3/3)^p}{p!}}`$
$`=`$ $`{\displaystyle \underset{l,p0}{}}a(y){\displaystyle \frac{(\alpha tb(y))^l}{l!}}{\displaystyle \frac{(\gamma t^3c(y)/3)^p}{p!}}`$
$`=`$ $`a(y)\mathrm{exp}\text{(}\alpha tb(y)+\gamma t^3c(y)/3\text{)}`$
with $`y=\beta t^2/3`$. Substituting for $`a`$, $`b`$, and $`c`$ the formulas above, we find
$$F_0^k(t)=\mathrm{cosh}^{12k}(\sqrt{3x})\mathrm{exp}\text{(}(\alpha \beta +2\gamma )\sqrt{3x/\beta ^3}2\gamma \mathrm{tanh}(\sqrt{3x})/\beta ^{3/2}\text{)}$$
which, since the new variable $`t`$ is related to the original variable $`x`$ by $`t=\sqrt{3y/\beta }=\beta ^{1/2}\mathrm{tanh}(\sqrt{3x})`$, is equivalent to (7). . $`\mathrm{}`$
Proof of (9). Consider first part (a). Regarded as a polynomial in $`\alpha `$ and $`\gamma `$ only, each $`\rho _{u,v,w}^{r2k+vw}`$ is homogeneous of degree $`u+w`$. So let us decompose $`\xi _{r,s}^k`$ likewise into its homogeneous summands relative to this $`\alpha ,\gamma `$-grading. They are nonzero only in $`\alpha ,\gamma `$-degree $`r2m`$ for $`m0`$. Indeed, using (9) and (7.4), we find that the part of $`\xi _{r,s}^k`$ in $`\alpha ,\gamma `$-degree $`r2m`$ equals
$$\frac{1}{3^m}\underset{i,j}{}\frac{(r2k+si)!}{(r2k)!}\frac{\alpha ^{r2mj}}{(r2mj)!}\frac{\beta ^{s+mj}}{(si)!}\frac{(2\gamma )^j}{i!(ji)!}\varphi _{mj+i}^k(ri,ji),$$
where we adopt the convention of summing over those indices where the factorials all have non-negative arguments.
The $`\rho `$-classes having total degree $`u+2v+3w=r+2s`$ and $`\alpha ,\gamma `$-degree $`u+w=r2m`$ are of the form $`\rho _{r2mw,s^{}w,w}^{r^{}+s^{}2w}`$ for $`w=0,\mathrm{},\mathrm{min}([(r^{}+s^{})/2],s^{},r2m)`$, where we have introduced the abbreviations $`r^{}=r2k`$ and $`s^{}=s+m`$. Using their definition (2.3), we may express any linear combination of these $`\rho `$-classes as
$`{\displaystyle \frac{1}{3^m}}{\displaystyle \underset{i,w}{}}a_w(q2wi)!{\displaystyle \frac{\alpha ^{ri2mw}}{(ri2mw)!}}{\displaystyle \frac{\beta ^{s^{}iw}}{(s^{}iw)!}}{\displaystyle \frac{(2\gamma )^{i+w}}{i!}}`$
$`=`$ $`{\displaystyle \frac{1}{3^m}}{\displaystyle \underset{j,w}{}}a_w(qwj)!{\displaystyle \frac{\alpha ^{r2mj}}{(r2mj)!}}{\displaystyle \frac{\beta ^{s^{}j}}{(s^{}j)!}}{\displaystyle \frac{(2\gamma )^j}{(jw)!}},`$
where $`a_w`$ are arbitrary scalars, $`q=r^{}+s^{}`$, and the factor of $`1/3^m`$ is inserted for convenience.
At least when $`s`$ is large enough that $`s^{}`$ and $`[(r^{}+s^{})/2]r2m`$, these span all the polynomials in $`\alpha ,\beta ,\gamma `$ of the given total degree and $`\alpha ,\gamma `$-degree. It is therefore possible to write the part of $`\xi _{r,s}^k`$ in $`\alpha ,\gamma `$-degree $`r2m`$ as a linear combination of this kind. The goal is to show that $`a_w=0`$ whenever either $`w>r^{}`$ or $`u+3w>r`$, that is, $`w>m`$.
The reader may worry that this will only prove the desired result for $`s`$ large compared to $`r`$. But, according to (2.3) and (7.4), the coefficient, in all of the polynomials we are concerned with, of the monomial $`\alpha ^a\beta ^{sb}\gamma ^c`$ for fixed $`a,b,c`$ is a rational function of $`s`$. So if a linear dependence between these polynomials can be established for sufficiently large $`s`$, then it holds for all $`s`$.
To determine the scalars $`a_w`$, set the coefficients of $`\alpha ^{r2mj}\beta ^{s^{}j}(2\gamma )^j`$ in the last two equations to be equal:
$$\underset{i}{}\frac{(r^{}+si)!(s^{}j)!}{r^{}!(si)!i!(ji)!}\varphi _{mj+i}^k(ri,ji)=\underset{w}{}\frac{a_w(qwj)!}{(jw)!}.$$
Let $`b_j`$ be the left-hand side, and $`L`$ the lower triangular matrix whose $`(j,w)`$ entry is $`(qwj)!/(jw)!`$. Here $`j,w`$ index the rows and columns, and run from $`0`$ to $`r2m`$. In vector notation, the equation above then says $`(b_j)=L(a_w)`$.
The inverse of $`L`$ is the lower triangular matrix whose $`(w,j)`$ entry is
$$(1)^{w+j}\frac{(q+12w)}{(wj)!(q+1wj)!}.$$
Indeed, the sum that needs to be demonstrated is
$$\underset{j=w}{\overset{w^{}}{}}(1)^{w^{}+j}\frac{(q+12w^{})(qwj)!}{(w^{}j)!(q+1w^{}j)!(jw)!}=\delta _{w,w^{}}.$$
This is obvious for $`ww^{}`$. For $`w<w^{}`$, if the summand is denoted $`N_j`$, then as Shalosh B. Ekhad has pointed out ,
$$(q+1ww^{})(w^{}w)N_j=(jw)(q+1wj)N_j(j+1w)(qwj)N_{j+1};$$
since the coefficient on the left is independent of $`j`$, and is nonzero for $`s`$ and hence $`q`$ large, the sum telescopes.
Hence
$$a_w=\underset{i,j}{}\frac{(1)^{wj}(q+12w)(r^{}+si)!(s^{}j)!}{(wj)!(q+1wj)!r^{}!(si)!i!(ji)!}\varphi _{mj+i}^k(ri,ji).$$
Now sum over all variables, and group the factorials together as binomial coefficients, to create the grand generating function
$`{\displaystyle \underset{m,q,s,w}{}}a_w{\displaystyle \frac{(q+1w)!}{(q+12w)(s^{}w)!}}M^mQ^qS^sW^w`$
$`=`$ $`{\displaystyle \underset{i,j,m,q,s,w}{}}(1)^{wj}\left(\genfrac{}{}{0pt}{}{r^{}+si}{r^{}}\right)\left(\genfrac{}{}{0pt}{}{s^{}j}{wj}\right)\left(\genfrac{}{}{0pt}{}{q+1w}{j}\right)\left(\genfrac{}{}{0pt}{}{j}{i}\right)\varphi _{mj+i}^k(ri,ji)M^mQ^qS^sW^w.`$
Using the binomial series, we can successively eliminate the sums over $`q`$, $`w`$, and $`s`$, to obtain
$$\underset{i,p,m}{}\left(\genfrac{}{}{0pt}{}{i+p}{i}\right)\frac{(1WQ)^m(WQ^2)^{i+p}M^{p+m}S^i}{Q(1Q)^{i+p+1}(1S(1WQ))^{r^{}+1}}\varphi _m^k(ri,p).$$
Substituting (9) for the sum over $`m`$ yields
$`{\displaystyle \frac{1}{Q(1S(1WQ))^{r^{}+1}}}{\displaystyle \underset{i,p}{}}\left(\genfrac{}{}{0pt}{}{i+p}{i}\right){\displaystyle \frac{(WQ^2)^{i+p}M^pS^i}{(1Q)^{i+p+1}}}`$
$``$ $`{\displaystyle \frac{1}{\mathrm{cosh}^{2k}\sqrt{X}}}{\displaystyle \frac{\sqrt{X}}{\mathrm{sinh}\sqrt{X}}}({\displaystyle \frac{\sqrt{X}}{\mathrm{tanh}\sqrt{X}}})^{ri}({\displaystyle \frac{3}{X}}{\displaystyle \frac{3\mathrm{tanh}\sqrt{X}}{X\sqrt{X}}})^p,`$
where $`X=3(1WQ)M`$. Applying the binomial theorem again and simplifying transforms this to
$`{\displaystyle \frac{1}{Q(1S(1WQ))^{r^{}+1}}}{\displaystyle \frac{1WQ}{1QWQ+WQ^2(1S(1WQ))\frac{\mathrm{tanh}\sqrt{X}}{\sqrt{X}}}}`$ (9.3)
$`{\displaystyle \frac{1}{\mathrm{cosh}^{2k}\sqrt{X}}}{\displaystyle \frac{\sqrt{X}}{\mathrm{sinh}\sqrt{X}}}({\displaystyle \frac{\sqrt{X}}{\mathrm{tanh}\sqrt{X}}})^r.`$
The goal is to show that the coefficient of $`M^mQ^qS^sW^w`$ vanishes in the above for $`q=r^{}+s^{}`$ and $`w>\mathrm{min}(r^{},m)`$. Since $`X=3(1WQ)M`$, it is equivalent to multiply the generating function (9.3) by $`3^m(1WQ)^m`$ and take the coefficient of $`X^mQ^qS^sW^w`$. But the second line of (9.3) is a power series in $`X`$ only, so this coefficient is a linear combination of the coefficients, for $`nm`$, of $`X^nQ^qS^sW^w`$ in
$$\frac{1}{Q(1S(1WQ))^{r^{}+1}}\frac{(1WQ)^{m+1}}{1QWQ+WQ^2(1S(1WQ))\frac{\mathrm{tanh}\sqrt{X}}{\sqrt{X}}}.$$
We will show that these all vanish.
In fact, we may replace $`\mathrm{tanh}\sqrt{X}/\sqrt{X}`$ by $`1+X`$ in the above. For this can be undone by substituting a power series of the form $`c_1X+c_2X^2+\mathrm{}`$ for $`X`$; hence the coefficients of $`X^nQ^qS^sW^w`$ in the former are linear combinations of $`X^pQ^qS^sW^w`$ in the latter, for $`pm`$.
Taking coefficients of $`X^p`$, $`S^s`$, $`W^w`$, and $`Q^{r^{}+s^{}}`$ in the resulting rational function yields
$$\underset{i=0}{\overset{s}{}}(1)^{w+i}\left(\genfrac{}{}{0pt}{}{r^{}+spi}{r^{}p}\right)\left(\genfrac{}{}{0pt}{}{p+i}{i}\right)\left(\genfrac{}{}{0pt}{}{s^{}pi}{wpi}\right)\left(\genfrac{}{}{0pt}{}{r^{}+s^{}+1w}{n+i}\right).$$
Let $`F(s,i)`$ be the $`i`$th term in the sum. As pointed out by Shalosh B. Ekhad , if we define $`G(s,i)`$ by
$$\frac{i(r^{}+s+1pi)(s^{}+1pi)(r^{}+s^{}+2w)(r^{}+s+s^{}+3pwi)}{(s+1i)(s^{}+1w)(r^{}+s^{}+2pwi)}F(s,i),$$
then by high-school algebra,
$`G(s,i+1)G(s,i)=`$
$`(r^{}+s+1w)(r^{}+s^{}+2w)F(s,i)(s+1)(r^{}+s^{}+2pw)F(s+1,i),`$
and so, summing over $`i`$, we deduce that the sum satisfies a linear recurrence relation in $`s`$:
$$(r^{}+s+1w)(r^{}+s^{}+2w)\underset{i=0}{\overset{s}{}}F(s,i)(s+1)(r^{}+s^{}+2pw)\underset{i=0}{\overset{s+1}{}}F(s+1,i)=0.$$
If $`w>r^{}`$ and $`pm`$, the coefficient of the first sum in the recurrence is $`0`$ for $`s=wr^{}1`$, and the coefficient of the second sum is nonzero for all subsequent $`s`$. Hence the sum vanishes for all $`s`$ sufficiently large, namely $`wr^{}`$.
If $`w>m`$ and $`pm`$, then every term in the sum is easily seen to vanish for $`s=0`$, and for $`s=r^{}m1+p+w`$ if this is positive. The recurrence then implies that the sum is $`0`$ for all positive $`s`$.
This completes the proof of part (a); the proof of (b) is similar. Because
$$\left(\genfrac{}{}{0pt}{}{r^{}+si}{r^{}}\right)\left(\genfrac{}{}{0pt}{}{r^{}+s1i}{r^{}}\right)=\left(\genfrac{}{}{0pt}{}{r^{}1+si}{r^{}1}\right),$$
the grand generating function has $`r^{}1`$ substituted for $`r^{}`$; hence the same is true for all subsequent formulas. . $`\mathrm{}`$
## 10 The relations divisible by $`𝜸`$
Many $`\mathrm{\Gamma }`$-invariant relations on $`_n`$ were computed in §8, but none of these relations were divisible by $`\gamma `$. To find some $`\mathrm{\Gamma }`$-invariant relations that are divisible by $`\gamma `$, we revive the space $``$ of flat connections, which is diffeomorphic to $`_0`$ as described in §2 of our previous paper . We will find a relationship between the cohomology at genus $`g`$ and genus $`g1`$. Accordingly, we will write $`^g`$ to indicate the dependence of $``$ on the genus. Let $`G=\mathrm{SL}(2,)`$, and define $`\mu _g:G^{2g}G`$ by $`\mu _g(A_j,B_j)=A_j^{}B_j^{}A_j^1B_j^1`$. Then $`^g=\mu _g^1(I)/G`$, where $`G`$ acts by simultaneous conjugation. The goal of this section is to prove the following
(10.1) Let $`\rho [\alpha ,\beta ,\psi _j]`$ be a relation on $`^{g1}`$. Then $`\psi _j\psi _{j+g}\rho `$ for each $`jg`$, and hence $`\gamma \rho `$, are relations on $`^g`$.
The proof will involve the following lemma.
(10.2) The only critical value of $`\mu _g`$ is the identity matrix $`I`$.
Proof. The derivative $`d\mu _g:𝔤^{2g}𝔤`$ at $`(A_j,B_j)G^{2g}`$ is easy to compute explicitly: see for example Goldman or Gunning \[10, Lemma 26\]. It is a sum of $`g`$ terms, the $`k`$th being conjugate to $`(a_j,b_j)(I\text{Ad}A_k^1)b_k(I\text{Ad}B_k^1)a_k`$. At a critical point, then, all $`g`$ of these maps must fail to surject.
For $`A\pm IG=\mathrm{SL}(2,)`$, it is easy to check by hand, using Jordan canonical form, that the image of $`(I\text{Ad}A^1):𝔤𝔤`$ is a 2-dimensional subspace from which the eigenspaces of $`A`$ can be recovered, and is a subalgebra if and only if $`A`$ is not diagonalizable. If the maps are not surjective, then either these 2-dimensional subspaces must coincide, or one of $`A_k`$ or $`B_k`$ is $`\pm I`$. In the former case, $`A_k`$ and $`B_k`$ must have a common eigenspace (if they are not diagonalizable) or eigenspaces (if they are). In any case, they must commute, and so $`\mu _g(A_j,B_j)=I`$. . $`\mathrm{}`$
Notice that the arguments of the last paragraph are special to $`\mathrm{SL}(2,)`$; the situation for $`\mathrm{SL}(r,)`$ with $`r>2`$ is more complicated.
Proof of (10). Let $`K=\mathrm{SU}(2)`$, and let $`LG`$ be the locus of matrices of the form $`U^1DU`$, where $`UK`$ and $`D=\text{diag}(\lambda ,1/\lambda )`$ for some positive real $`\lambda `$. Then $`L`$ is a smooth, contractible submanifold of $`G`$ whose tangent space at the identity is $`i`$ times that of $`K`$.
Let $``$ be the intersection $`\mu _g^1(I)(G^{2g2}\times L^2)`$. This is preserved by the $`K`$-action, and there are inclusions
$$\mu _{g1}^1(I)\times \{I\}\times \{I\}\mu _g^1(I).$$
Dividing by the $`K`$-action yields inclusions
$$\stackrel{~}{}^{g1}/K\stackrel{~}{}^g.$$
Here $`\stackrel{~}{}^g=\mu _g^1(I)/K`$, which is a $`G/K`$-bundle over $`^g`$. Since $`G/K`$ is contractible, this is homotopy equivalent to $`^g`$. It is not hard to check, using the definition of the universal classes in §1 of our previous paper , that $`\alpha `$, $`\beta `$, and $`\gamma `$, regarded as classes on $`\stackrel{~}{}^g`$, restrict to their counterparts on $`\stackrel{~}{}^{g1}`$.
It therefore suffices to prove the following two claims: that $`\stackrel{~}{}^{g1}/K`$ induces an isomorphism on cohomology; and that $`/K\stackrel{~}{}^g`$ is Poincaré dual to $`\psi _g\psi _{2g}`$. For then $`\psi _g\psi _{2g}\rho `$ must be a relation on $`\stackrel{~}{}^g`$. The result follows by symmetry, since the action of $`\mathrm{\Gamma }`$ on $`[\alpha ,\beta ,\psi _j]`$ certainly preserves the ideal of relations, and there is an element of $`\mathrm{\Gamma }`$ taking $`\psi _g\psi _{2g}`$ to $`\psi _j\psi _{j+g}`$ for each $`j`$.
To prove the first claim, first note that $``$ can be regarded as the fibered product $`G^{2g2}\times _G(L\times L)`$, where the map $`G^{2g2}G`$ is $`\mu _{g1}`$ and the map $`L\times LG`$ is $`(A,B)ABA^1B^1`$. A direct computation shows that no two $`A,BL`$ have $`ABA^1B^1=I`$; hence by (10) the map $`L\times LG`$ never touches a critical value of $`\mu _{g1}`$. So $``$ is locally trivial over $`L\times L`$ with fiber $`\mu _{g1}^1(I)`$. Since $`L\times L`$ is contractible, this implies that $``$ is homeomorphic to $`L\times L\times \mu _{g1}^1(I)`$.
It follows that $`\stackrel{~}{}^{g1}`$ and $`/K`$ are homotopy equivalent. Indeed, they are homotopy equivalent to the homotopy quotients $`(\mu _{g1}^1(I)\times EK)/K`$ and $`(\times EK)/K`$ respectively, and the latter retracts onto the former, since $`BK`$ is a direct limit of manifolds and $`L\times L`$ is contractible.
To prove the second claim, first note that $`L`$ meets $`KG`$ transversely at the single point $`I`$. It is therefore Poincaré dual to the standard generator of $`H^3(G,)`$. We can now either imitate the argument given by the second author \[31, Prop. 19.3\] for the $`\mathrm{SU}(2)`$ space $`𝒩^g=(\mu _g^1(I)K^{2g})/K`$, or simply use that result. It tells us that the natural maps in the top row of the diagram
$$\begin{array}{ccccc}K^{2g}& & \mu _g^1(I)K^{2g}& & 𝒩^g\\ & & & & \\ G^{2g}& & \mu _g^1(I)& & \stackrel{~}{}^g\end{array}$$
induce isomorphisms on $`H^3`$ under which the generator of the $`j`$th copy of $`H^3(K,)`$ corresponds to $`\psi _j`$. Since the outer columns also induce isomorphisms on $`H^3`$, so does every map in the diagram.
Since the pull-back by inclusion is Poincaré dual to transverse intersection, it now suffices to check that $`G^{2g2}\times L^2`$ is transverse to $`\mu _g^1(I)`$, or equivalently, that at every point of $``$ the derivative $`d\mu _g`$ remains surjective when restricted to the tangent space to $`G^{2g2}\times L^2`$. But this is true even if we restrict further to the tangent space to $`G^{2g2}`$, since as stated before we are at a regular value of $`\mu _{g1}`$. . $`\mathrm{}`$
## 11 The cohomology not fixed by $`𝚪`$
Everything so far has been about the $`\mathrm{\Gamma }`$-invariant part of $`H^{}(_n)^\mathrm{\Sigma }`$, generated by $`\alpha `$, $`\beta `$, and $`\gamma `$. Now it is time to say something about the non-invariant part and the classes $`\psi _j`$.
We begin with a result relating the non-invariant parts of the cohomology of the symmetric product $`C_m`$ at genus $`g`$ to the invariant part at lower genera.
(11.1) As a $`\mathrm{\Gamma }`$-module, the cohomology of the symmetric product $`C_m^g`$ has the form
$$H^{}(C_m^g)=\underset{k=0}{\overset{g}{}}\mathrm{\Lambda }_0^k(\xi )[\eta ,\theta ]/I(C_{mk}^{gk}),$$
where $`I(C_{mk}^{gk})`$ is the ideal of relations between $`\eta `$ and $`\theta `$ in $`C_{mk}^{gk}`$, with the convention that $`I(C_{mk}^{gk})=[\eta ,\theta ]`$ if $`mk<0`$.
Proof. As shown in §5, there is a surjection of $`\mathrm{\Gamma }`$-algebras
$$\underset{k=0}{\overset{g}{}}\mathrm{\Lambda }_0^k(\xi )[\eta ,\theta ]H^{}(C_m^g),$$
where $`\mathrm{\Lambda }_0^k`$, being irreducible, is spanned by the orbit of $`\xi _1\xi _2\mathrm{}\xi _k`$ under $`\mathrm{\Gamma }`$. It therefore suffices to show that a polynomial $`p(\eta ,\theta )`$ is a relation on $`C_{mk}^{gk}`$ if and only if $`p(\eta ,\theta )\xi _1\xi _2\mathrm{}\xi _k`$ is a relation on $`C_m^g`$.
By Poincaré duality the latter is true if and only if for all polynomials $`q(\eta ,\xi _j)`$,
$$p(\eta ,\theta )\xi _1\xi _2\mathrm{}\xi _kq(\eta ,\xi _j)[C_m^g]=0.$$
(11.2)
Now it follows from the description of $`H^{}(C_m^g)`$ in Macdonald that $`_{j=1}^{2g}\xi _j^{p_j}\eta ^q[C_m^g]=1`$ if $`_{j=1}^gp_j+q=m`$ (so that the degrees match) and $`p_j=p_{j+g}1`$ for each $`jg`$ (so that it becomes a monomial in $`\eta `$ and $`\theta _j=\xi _j\xi _{j+g})`$. Otherwise it equals 0. Hence in (11.2) we only need to consider the case
$$q(\eta ,\xi _j)=\xi _{g+1}\mathrm{}\xi _{g+k}r(\eta ,\theta _{k+1},\mathrm{},\theta _g).$$
But $`q`$ can be averaged with its images under all permutations of $`\theta _{k+1},\mathrm{},\theta _g`$ without changing the value of (11.2). Hence we only need to consider the case where $`r`$ is a polynomial in $`\eta `$ and $`\theta _{k+1}+\mathrm{}+\theta _g`$. But then
$`p(\eta ,\theta )\xi _1\xi _2\mathrm{}\xi _kq(\eta ,\xi _j)[C_m^g]`$ $`=`$ $`(1)^kp(\eta ,\theta )\theta _1\theta _2\mathrm{}\theta _kr(\eta ,\theta _{k+1}+\mathrm{}+\theta _g)[C_m^g]`$
$`=`$ $`(1)^kp(\eta ,\theta _{k+1}+\mathrm{}+\theta _g)\theta _1\theta _2\mathrm{}\theta _kr(\eta ,\theta _{k+1}+\mathrm{}+\theta _g)[C_m^g].`$
This always vanishes if $`k>m`$. Otherwise, it equals $`(1)^kp(\eta ,\theta )r(\eta ,\theta )[C_{mk}^{gk}]`$. This vanishes for all $`r`$ if and only if $`p(\eta ,\theta )q(\eta ,\xi _j)[C_{mk}^{gk}]`$ vanishes for all polynomials $`q`$ in $`\eta `$ and $`\xi _j`$, since it does not alter the latter expression to replace $`q`$ with its projection on the $`\mathrm{\Gamma }`$-invariant part, which is a polynomial in $`\eta `$ and $`\theta `$. Again by Poincaré duality, this is equivalent to $`p(\eta ,\theta )=0`$ in $`C_{mk}^{gk}`$. . $`\mathrm{}`$
(11.3) As a $`\mathrm{\Gamma }`$-module, the $`\mathrm{\Sigma }`$-invariant part of the $`T`$-equivariant cohomology of $`_n^g`$ has the form
$$H_T^{}(_n^g)^\mathrm{\Sigma }\underset{k=0}{\overset{g}{}}\mathrm{\Lambda }_0^k(\psi )[\alpha ,\beta ,\gamma ,u]/I_T(_{n+k}^{gk}).$$
Consequently, as a $`\mathrm{\Gamma }`$-module, the $`\mathrm{\Sigma }`$-invariant part of the ordinary cohomology of $`_n^g`$ has the form
$$H^{}(_n^g)^\mathrm{\Sigma }\underset{k=0}{\overset{g}{}}\mathrm{\Lambda }_0^k(\psi )[\alpha ,\beta ,\gamma ]/I(_{n+k}^{gk}).$$
Proof. First we show that if $`\lambda \mathrm{\Lambda }_0^k(\psi )`$ and $`rI_T(_{n+k}^{gk})`$, then $`\lambda rI_T(_n^g)`$. Certainly $`r`$ restricts to relations between $`\alpha `$, $`\beta `$, $`\gamma `$ and $`u`$ on $`𝒩^g=F_0^g`$ and between $`\eta `$, $`\theta `$ and $`u`$ on the remaining fixed components $`F_d^g_n^g`$. On the other hand, by Kirwan’s theorem (3)(iii) it suffices to show that $`\lambda r`$ restricts to similar relations on the fixed components $`F_d^{gk}`$ of $`_{n+k}^{gk}`$.
The case of $`d>0`$ follows immediately from the lemma. As for $`d=0`$, Proposition 2.5 of King-Newstead asserts that if $`\lambda \mathrm{\Lambda }_0^k(\psi )`$ and $`r`$ is a relation between $`\alpha `$, $`\beta `$, $`\gamma `$ on $`𝒩^g`$, then $`\lambda r`$ is a relation on $`𝒩^{gk}`$. This is exactly what is needed.
The left-hand side is therefore a quotient of the right-hand side. To complete the proof, it remains only to check that the $`\mathrm{\Sigma }`$-invariant, $`T`$-equivariant Poincaré polynomials agree. But if $`P=_it^idimH^i`$, then
$`(1t^2)P_T^\mathrm{\Sigma }(_n^g)`$
$`=`$ $`P^\mathrm{\Sigma }(_n^g)`$
$`=`$ $`P^\mathrm{\Sigma }(𝒩^g)+{\displaystyle \underset{d=1}{\overset{g+\left[\frac{n1}{2}\right]}{}}}t^{2g+2d2}P(C_{2g+n12d}^g)`$
$`=`$ $`{\displaystyle \underset{k=0}{\overset{g}{}}}\left(\left(\genfrac{}{}{0pt}{}{2g}{k}\right)\left(\genfrac{}{}{0pt}{}{2g}{k2}\right)\right)(t^{3k}P^\mathrm{\Sigma }(𝒩^{gk})+{\displaystyle \underset{d=1}{\overset{g+\left[\frac{n1}{2}\right]}{}}}t^{2g+2d2}t^kP(C_{2g+n12dk}^{gk}))`$
$`=`$ $`{\displaystyle \underset{k=0}{\overset{g}{}}}\left(\left(\genfrac{}{}{0pt}{}{2g}{k}\right)\left(\genfrac{}{}{0pt}{}{2g}{k2}\right)\right)t^{3k}(P^\mathrm{\Sigma }(𝒩^{gk})+{\displaystyle \underset{d=1}{\overset{g+\left[\frac{n1}{2}\right]}{}}}t^{2(gk)+2d2}P(C_{2(gk)+n12d+k}^{gk}))`$
$`=`$ $`{\displaystyle \underset{k=0}{\overset{g}{}}}\left(\left(\genfrac{}{}{0pt}{}{2g}{k}\right)\left(\genfrac{}{}{0pt}{}{2g}{k2}\right)\right)t^{3k}P^\mathrm{\Sigma }(_{n+k}^{gk})`$
$`=`$ $`(1t^2){\displaystyle \underset{k=0}{\overset{g}{}}}\left(\left(\genfrac{}{}{0pt}{}{2g}{k}\right)\left(\genfrac{}{}{0pt}{}{2g}{k2}\right)\right)t^{3k}P_T^\mathrm{\Sigma }(_{n+k}^{gk}),`$
using Kirwan’s theorem on the Leray sequence (3)(ii) in steps 1 and 6, the perfection of the Morse stratification (3)(i) in steps 2 and 5, the lemma and the result of King-Newstead in step 3, and high-school algebra in step 4. . $`\mathrm{}`$
## 12 Wrap-up
At last we can show that the polynomials $`\rho _{r,s,t}^c`$ of (2.3) are relations, using every tool at our disposal: (9), (10), (11), and a dimension count. Recall that $`I_n^g`$ is the ideal of $`\rho `$-classes introduced in §2.
(12.1) Every element of $`I_n^g`$ is a relation on $`_n^g`$.
Proof. It is actually more convenient to work with $`n+2`$ than $`n`$, so let $`n2`$. Then every element of $`I_{n+2}^g`$ has total degree $`2g+n`$. In degree $`2g+n`$, $`I_{n+2}^g`$ is spanned by the elements $`\rho _{n2i,g+i,0}^{g+i}`$ for $`i=0,\mathrm{},[n/2]`$. The degrees with respect to $`\alpha `$ are all different, so these are linearly independent. The relations $`\xi _{n2i,g+i}^{[n/2]i}`$ of type (i) given in (8) are all of degree $`2g+n`$, also number $`[n/2]+1`$, and are linearly independent for the same reason. By (9) they are in the linear span of the $`\rho _{n2i,g+i,0}^{g+i}`$. Hence this equals the linear span of the relations of type (i), so all the $`\rho _{n2i,g+i,0}^{g+i}`$ are relations.
The diagram below is intended to help the reader visualize the ideal of relations. Each dot represents one of the $`\rho _{r,s,t}^c`$. The total degree is the vertical coordinate, and $`r`$ is the horizontal coordinate. The two edges of the dotted region reflect the two constraints imposed by (2.4). To avoid having to draw a third dimension, only those relations with $`t=0`$ have been shown. The dotted region for any fixed $`t>0`$ would look similar, only translated in a northwesterly direction. The relations established in the previous paragraph are those in the bottom row.
The rest of the proof proceeds by induction on the total degree. We have already seen that the part of $`I_n^g`$ in total degree $`2g+n`$ consists entirely of relations. Now fix $`j>0`$ and consider the part of $`I_{n+2}^g`$ in degree $`2g+n+j`$. Assume by induction that for all $`g`$ and $`n`$, the parts of $`I_{n+2}^g`$ in degree $`<2g+n+j`$ are known to be relations on $`_{n+2}^g`$.
In particular, if $`\rho _{r,s,t}^cI_{n+2}^g`$ has degree $`2g+n+j`$ and $`t>0`$, then $`\rho _{r,s,0}^cI_{n+2}^{gt}`$ has degree $`2(gt)+n+(jt)`$, and hence is a relation on $`_{n+2}^{gt}`$. By (11), $`\psi _1\mathrm{}\psi _{n+2}\rho _{r,s,0}^c`$ is a relation on $`_0^{gt+n+2}`$, so by (10), $`\psi _1\mathrm{}\psi _{n+2}\rho _{r,s,t}^c`$ is a relation on $`_0^{g+n+2}`$, and by (11) again, $`\rho _{r,s,t}^c`$ is a relation on $`_{n+2}^g`$. Because these relations have $`t>0`$, they are not shown in the diagram.
Only the $`\rho _{r,s,t}^c`$ with $`t=0`$ remain. Consider first those relations of the form $`\rho _{0,s,0}^c=c!/s!\beta ^sI_{n+2}^g`$. These are the relations at the left-hand edge of the diagram. If the degree equals $`2g+n+j`$, then since $`j>0`$,
$$2s=2g+n+j2g2+(n+3)$$
and
$$3s=3g+\frac{3}{2}n+\frac{3}{2}j>3g3+(n+3),$$
so in fact $`\rho _{0,s,0}^cI_{n+3}^g`$. Since it has degree $`2g+n+j=2g+(n+1)+(j1)`$, by the induction hypothesis again it is a relation on $`_{n+3}^g`$, and hence on $`_{n+2}^g_{n+3}^g`$.
What if $`t=0`$ but $`r>0`$? It is easily checked that
$$r\rho _{r,s,0}^c=c\alpha \rho _{r1,s,0}^{c1}+(cr)\rho _{r1,s1,1}^{c2}.$$
Now if $`r+2s=2g+n+j`$ and $`1rn+3j2`$ (so that we are not at the right-hand edge of the diagram), we know from the induction hypothesis that $`\rho _{r1,s,0}^{c1}`$ is a relation. And certainly $`\rho _{r,s1,1}^{c2}`$ is a relation, since it belongs to $`I_{n+2}^g`$ and has $`t>0`$. Hence under these circumstances $`\rho _{r,s,0}^c`$ is a relation. This fills in the interior of the diagram; the arrows depict multiplication by $`\alpha `$ (modulo $`\gamma `$), which takes $`\rho _{r1,s,0}^{c1}`$ to $`\rho _{r,s,0}^c`$.
Only one class of total degree $`2g+n+j`$ remains in $`I_{n+2}^g`$. This is $`\rho _{n+3j,gj}^g`$, at the right-hand edge of the diagram. We know from (8)(ii) that the class $`\xi _{n+3j,gj}^{[n/2]+j}`$ is a relation if $`n`$ is odd, and that $`\xi _{n+3j,gj}^{[n/2]+j}\beta \xi _{n+3j,gj1}^{[n/2]+j}`$ is a relation if $`n`$ is even. Since the leading term with respect to $`\alpha `$ of $`\xi _{r,s}^k`$ is $`\left(\genfrac{}{}{0pt}{}{r2k+s}{r2k}\right)\alpha ^k/k!`$, the monomial $`\alpha ^{n+3j}\beta ^{gj}`$ appears in these relations with a nonzero coefficient. On the other hand, by (9) these relations can be expressed as a linear combination of $`\rho _{r,s,t}^cI_{n+2}^g`$. Since $`\rho _{n+3j,gj}^g`$ is the class of maximal degree $`n+3j`$ with respect to $`\alpha `$ among these, its coefficient in this combination must be nonzero. It is therefore indeed a relation. This completes the proof. . $`\mathrm{}`$
(12.2) Every relation between $`\alpha `$, $`\beta `$, $`\gamma `$ on $`_n^g`$ is an element of $`I_n^g`$.
Proof. Since the converse has just been shown, it suffices to show that $`dimH^I(_n^g)=dim[\alpha ,\beta ,\gamma ]/I_n^g`$, where $`H^I`$ denotes the part of $`H^{}`$ invariant under the action of both $`\mathrm{\Sigma }`$ and $`\mathrm{\Gamma }`$.
(In the proofs of the analogous statement for $`𝒩^g`$ , it has been customary to show that all the Betti numbers match up. This can certainly be done for $`_n^g`$, but the cruder statement about overall dimension is clearly sufficient.)
Now since by (3)(i) the $`^\times `$ action is perfect, the dimension of $`H^I(_n^g)`$ may be expressed as a sum over the fixed components enumerated in (4):
$$dimH^I(_n^g)=dimH^I(𝒩^g)+\underset{d=1}{\overset{g+\left[\frac{n1}{2}\right]}{}}dimH^I(C_{2g+n12d}).$$
The Poincaré polynomial of $`H^I(𝒩^g)`$ is
$$\frac{(1t^{2g})(1t^{2g+2})(1t^{2g+4})}{(1t^2)(1t^4)(1t^6)};$$
this is clear, for example, from the presentation with three generators and three relations . To find $`dimH^I(𝒩^g)`$ we want to substitute $`t=1`$. Of course 0 appears in the denominator, but the limit as $`t1`$ can easily be evaluated to $`\left(\genfrac{}{}{0pt}{}{g+2}{3}\right)`$ by substituting $`t^2=1+ϵ`$, then using $`(1+ϵ)^k=1+kϵ+O(ϵ^2)`$.
As for $`dimH^I(C_m)`$, it follows easily from the discussion in Arbarello et al. \[1, VII B\] that this equals $`\left[\frac{m+2}{2}\right]\left[\frac{m+3}{2}\right]`$ if $`m2g1`$ and $`(g+1)(mg+1)`$ if $`m2g1`$. Hence
$`dimH^I(_n^g)`$
$`=`$ $`\left(\genfrac{}{}{0pt}{}{g+2}{3}\right)+{\displaystyle \underset{d=1}{\overset{\left[\frac{n1}{2}\right]}{}}}(g+1)(g2d+n)+{\displaystyle \underset{d=\left[\frac{n1}{2}\right]+1}{\overset{\left[\frac{n1}{2}\right]+g}{}}}\left(gd+1+\left[\frac{n1}{2}\right]\right)\left(gd+1+\left[\frac{n}{2}\right]\right).`$
On the other hand,
$$dim[\alpha ,\beta ,\gamma ]/I_n^g=\underset{r,s,t}{}\mathrm{\hspace{0.17em}1}$$
where the right-hand sum runs over all non-negative $`r,s,t`$ such that $`tg`$ and $`r+3s+3t3g3+n`$ or $`r+2s+2t<2g2+n`$. This can be re-written as
$`{\displaystyle \underset{t=0}{\overset{g}{}}}{\displaystyle \underset{s=0}{\overset{g1+\left[\frac{n}{2}\right]t}{}}}{\displaystyle \underset{r=0}{\overset{\genfrac{}{}{0pt}{}{\mathrm{max}(3g3+n3s3t,}{2g3+n2s2t)}}{}}}1`$
$`=`$ $`{\displaystyle \underset{t=0}{\overset{g}{}}}\left({\displaystyle \underset{s=0}{\overset{g1t}{}}}(3g2+n3s3t)+{\displaystyle \underset{s=0}{\overset{g1+\left[\frac{n}{2}\right]t}{}}}(2g2+n2s2t)\right).`$
It is straightforward, using high-school algebra and the identities $`_{d=1}^kd=k^2/2+k/2`$ and $`_{d=1}^kd^2=k^3/3+k^2/2+k/6`$, to show that this equals the above formula for $`dimH^I(_n^g)`$. . $`\mathrm{}`$
Proof of (2). The two results above show that the ideal $`I(_n^g)`$ of $`\mathrm{\Gamma }`$-invariant relations on $`_n^g`$ is precisely $`I_n^g`$. Now apply (11). . $`\mathrm{}`$
## 13 Relationship with other papers
The present paper is closely related to several other works by the authors; in this final section we indicate briefly a few of the points of contact.
The first author has constructed a compactification $`\overline{}_n`$ of the moduli space of Higgs bundles , by adding a divisor $`Z_n`$ at infinity which is a quotient by $`T=^\times `$ of an open subset of $`_n`$. Indeed, $`\overline{}_n`$ itself is a quotient by $`T`$ of $`_n\times `$. Many of the constructions given herein apply to $`\overline{}_n`$ and $`Z_n`$. In particular, there are direct limits $`\overline{}_{\mathrm{}}`$ and $`Z_{\mathrm{}}`$. Just as $`_{\mathrm{}}`$ is shown in (9.7) of our previous paper to be homotopy equivalent to $`B\overline{𝒢}`$, we expect $`\overline{}_{\mathrm{}}`$, and also $`Z_{\mathrm{}}`$, to be homotopy equivalent to $`B𝒢`$, the classifying space of the full gauge group.
The cohomology rings of $`\overline{}_{\mathrm{}}`$ and $`Z_{\mathrm{}}`$, and hence those of $`\overline{}_n`$ and $`Z_n`$, will have generators like those of $`_n`$, but with one additional generator $`hH^2`$, corresponding to the class discarded in the proof of (10.1) of our previous paper. Indeed, the quotient map $`H_T^{}(_n)H^{}(Z_n)`$ is surjective, and $`h`$ is the image of $`u`$. It can be shown that the kernel is the image in $`H_T^{}(_n)`$ of the compactly supported cohomology. The relations between the generators in $`H^{}(Z_n)`$ are therefore of two types: those coming from relations in $`H_T^{}(_n)`$, and those coming from the compactly supported cohomology. The former are covered by the results of this paper. As for the latter, they ought to be determined by the results of another paper of the first author , in which the intersection pairings between the cohomology and the compactly supported cohomology of $`_0`$ were computed (and shown to vanish).
By studying the stratification of $`_{\mathrm{}}`$ by the HN type of the underlying bundle $`E`$ — not $`(E,\varphi )`$ — in the rank 2 case, the first author has been able to extract the relations in the cohomology of the lowest stratum, which retracts onto $`𝒩`$. This recovers the description of the ring $`H^{}(𝒩)`$, given by several authors , in an especially simple and geometrical fashion. It will be described in a forthcoming paper , where it will also be shown how consideration of the Mumford conjecture leads to a natural geometrical proof of the generation theorem for the moduli space of Higgs bundles in ranks 2 and 3.
The second author has studied moduli problems providing smooth resolutions of the upward and downward flows from the components of the fixed-point set of $`_n`$ in the rank 2 case. The downward flow is intriguing because it can be interpreted as a master space of Bradlow pairs, but the upward flow is particularly related to the present paper. It is relatively simple to describe, but it contains the parts of $`H^{}(_n)`$ not invariant under $`\mathrm{\Sigma }=_2^{2g}`$, in the sense that none of them are killed by the restriction map. These upward flows can therefore be used to complete the description of the cohomology rings $`H^{}(_n)`$, by characterizing the part not invariant under $`\mathrm{\Sigma }`$. In fact this is not so difficult, since for dimension reasons these classes have square 0, and are killed by the $`\psi _j`$. So the products with $`\alpha `$ and $`\beta `$ are all that must be computed. Details will appear in a forthcoming paper .
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# Symmetric Classical Mechanics
## 1 Introduction
The relation between classical and quantum mechanics is beset by a number of quite subtle problems. On one side the quantum formulation of a given dynamical system is typically obtained by applying a quantization prescription to the classical formulation of the system. There are several different quantization prescriptions that one might use. However, all of them display order problems, which means that there is an ambiguity in the choice of the quantum system that corresponds to the original classical one. On the other hand one expects that for an adequate choice of quantum initial data, satisfying some general conditions, the quantum formulation of most dynamical systems should be able to reproduce the predictions of the original classical formulation . Nevertheless, it is still not completely clear what these general conditions should be . This has been an active field of research. Firstly because it constitutes a fundamental problem. Quantum mechanics is believed to provide the most fundamental description of all physical systems. Many of these systems display a classical behaviour. Therefore, quantum mechanics should be able to explain the emergence of a classical domain and moreover to reproduce the predictions of classical mechanics. Secondly, because grasping the conditions that determine the emergence of a classical domain is expected to play an important role in several different fields of research like for instance, quantum cosmology, quantization of closed dynamical systems and semiclassical gravity , just to name a few.
This paper concerns the problem of the semiclassical limit of quantum mechanics. More precisely, we want to investigate whether the semiclassical limit of quantum mechanics might be correctly described by a classical dynamical structure different from ordinary classical mechanics. Our motivation comes from a set of results presented elsewhere . There we were able to prove that when the quantum initial data of an arbitrary dynamical system satisfy a set of classicality criteria the quantum predictions will be consistent (in some precise sense) with the predictions of the classical formulation of the system. The derivation of the criteria points out an interesting fact: for a general set of quantum initial data satisfying the criteria the classical predictions that display the highest degree of consistency with the quantum predictions are not the ones obtained by using the standard formulation of classical mechanics. In this paper the aim is to derive the dynamical framework that provides these new classical predictions. As a result a new theory of classical mechanics will be presented. The theory will be named symmetric classical mechanics and its properties will be studied thoroughly. In particular we will see that: i) The new theory displays a fully consistent canonical structure, ii) the quantization prescription for symmetric classical mechanics is an isomorphism between the classical and quantum algebras of observables and is then not riddled with ordering ambiguities, iii) continuous canonical transformations and, in particular, the time evolution are generated by unitary operators, iv) the time-evolution unitary operator is the solution of a classical version of the Schrödinger equation, v) the limit of symmetric classical mechanics as $`\mathrm{}0`$ is standard classical mechanics and finally, vi) we venture the possibility that symmetric classical mechanics could be trivially coupled to quantum mechanics to obtain a consistent theory of hybrid classical-quantum dynamics.
Most of the previous properties are a direct consequence of the fact that the dynamical structure of symmetric classical mechanics is given by the Moyal bracket. This is a Lie bracket that can be obtained by a deformation of the Poisson bracket . It was first derived as the dynamical structure for the Wigner distribution function formulation of quantum mechanics and has been used as the starting point for a number of semiclassical approximation procedures namely in the context of the quantum dynamics of classically chaotic systems . Extensive reviews of the Moyal-Weyl-Wigner formulation of quantum mechanics are given in . Here, instead, we derive the Moyal bracket as the dynamical structure of the semiclassical limit of quantum mechanics proving, as a by-product, that classical and quantum dynamics can be formulated in terms of the same bracket structure.
## 2 Classicality Criteria
In a previous paper we developed two different classicality criteria providing a measure of the degree of classicality of an arbitrary quantum system. The formalism presented can be summarized in three main steps:
1) Let us consider an arbitrary dynamical system with $`N`$ degrees of freedom. Let $`(q_i,p_i)`$, $`1=1..N`$ or simply $`O_i`$, $`i=\mathrm{1..2}N`$ be a set of canonical variables spanning the phase space of the system. Classical and quantum mechanics provide two alternative descriptions of the configuration of the system at an arbitrary time $`t_0`$. The classical description is given by a set of values $`O_i^0=O_i(t_0)`$ for the canonical variables plus the associated error margins $`\delta _i`$. The quantum description is given by the physical wave function $`|\varphi >`$ belonging to the physical Hilbert space $``$. The first step is to provide a measure of the consistency of these two descriptions. In we proposed two such consistency criteria. Let us review one of them:
Let $`0p<1`$ be an arbitrary probability. Let $`M`$ be some positive integer and let us consider the set of intervals of the type:
$$I_i(p,M)=[O_i^0\frac{\delta _i}{(1p)^{1/(2M)}},O_i^0+\frac{\delta _i}{(1p)^{1/(2M)}}],$$
(1)
associated to each classical observable $`O_i`$. In each of the previous intervals we can calculate the probability $`p_i`$ generated by the wave function $`|\varphi >`$, in the representation of the corresponding quantum observable $`\widehat{O}_i`$:
$$p_i(p,M)=\underset{a_iI_i(p,M),k}{}|<a_i,k|\varphi >|^2,$$
(2)
where $`|a_i,k>`$ is the general eigenvector of $`\widehat{O}_i`$ with eigenvalue $`a_i`$ and degeneracy index $`k`$. For given values of the classical and quantum data $`O_i^0`$, $`\delta _i`$ and $`|\varphi >`$, this probability is an exclusive function of $`p`$ and $`M`$. We can now state the consistency criterion:
Definition 1 \- Consistency Criterion
The classical and quantum data, describing a given configuration of the dynamical system will be $`M`$-order consistent if and only if for all $`0p<1`$ and all $`i=1,\mathrm{},2N`$ the condition $`p_i(p,M)p`$ is satisfied, i.e.;
$$\underset{a_iI_i(p,M),k}{}|<a_i,k|\varphi >|^2p,p[0,1[,i=1,\mathrm{},2N,$$
(3)
where $`I_i(p,M)`$ is given by eq.(1).
This criterion provides a measure of how peaked the wave function is - in the representation of each of the quantum observables - around the classical error margin of the corresponding classical observable. Notice that, in particular, the higher the degree of consistency the bigger the probability that a quantum measurement provides a value inside the corresponding classical error interval.
2) The second and main step in developing the classicality criterion was the derivation of the following expansion:
$`\widehat{A}A`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{2N}{}}}{\displaystyle \frac{A}{O_i}}(\widehat{O}_iO_i)+{\displaystyle \frac{1}{2}}{\displaystyle \underset{i,j=1}{\overset{2N}{}}}{\displaystyle \frac{^2A}{O_iO_j}}(\widehat{O}_iO_i)(\widehat{O}_jO_j)+\mathrm{}.`$ (4)
where $`\widehat{A}`$ is a general operator, $`\widehat{A}=F(\widehat{O}_i)`$ and $`A`$ is some classical version (which version is yet to be discussed) of $`\widehat{A}`$, $`A=G(O_i)`$. If the expansion (4) is valid then we can easily obtain the following relation:
$`<E^m(\widehat{A},\varphi ,A^0)|E^m(\widehat{A},\varphi ,A^0)>`$ (5)
$`{\displaystyle \underset{i_1=1}{\overset{2N}{}}}\mathrm{}{\displaystyle \underset{i_m=1}{\overset{2N}{}}}{\displaystyle \underset{j_1=1}{\overset{2N}{}}}\mathrm{}{\displaystyle \underset{j_m=1}{\overset{2N}{}}}{\displaystyle \underset{k=1}{\overset{m}{}}}{\displaystyle \underset{s=1}{\overset{m}{}}}{\displaystyle \frac{A}{O_{i_k}}}|_{O_{ik}=O_{ik}^0}\left({\displaystyle \frac{A}{O_{j_s}}}\right)^{}|_{O_{js}=O_{js}^0}<E_{O_{i1},\mathrm{},O_{im}}|E_{O_{j1},\mathrm{},O_{jm}}>+\mathrm{}`$
where $`A^0=G(O_i^0)`$ and $`|E_X^m>=|E^m(\widehat{X},\varphi ,X^0)>`$ is named the mth-order error ket of the operator $`\widehat{X}`$ around the classical value $`X^0`$ and is given by $`|E_X^m>=(\widehat{X}X^0)^m|\varphi >`$ for all $`m𝒩`$. Moreover, $`|E_{O_{j1},\mathrm{},O_{jm}}>=(\widehat{O}_{j_1}O_{j_1}^0)\mathrm{}.(\widehat{O}_{j_m}O_{j_m}^0)|\varphi >`$ is named the $`m`$-order mixed error ket. The most important property of the error ket framework is the following: if $`<E_X^m|E_X^m>\mathrm{\Delta }^{2m}`$ then, in the representation of $`\widehat{X}`$ the wave function $`|\varphi >`$ has at least a probability $`p`$ confined to the interval $`I_m=[X^0\frac{\mathrm{\Delta }}{(1p)^{1/2m}},X^0+\frac{\mathrm{\Delta }}{(1p)^{1/2m}}]`$. That is, $`|\varphi >`$ is m-order consistent with the classical interval $`[X^0\mathrm{\Delta },X^0+\mathrm{\Delta }]`$ (check for the consistency criterion in step 1). The reader should refer to for a detailed discussion of the error ket formalism.
This property points out a quite obvious way of developing a classicality criterion. Let us assume for the moment that (4) is valid for $`\widehat{A}=\widehat{O}_i(t)`$ and $`A=O_i(t)`$, $`i=\mathrm{1..2}N`$ (where $`O_i(t)`$ is some classical version of $`\widehat{O}_i(t)`$ \- see point 3)) and let us consider all the sequences of observables $`S_{i_k}=O_{i_1},\mathrm{}.,O_{i_n}`$ associated to the sequences of values $`i_k=i_1,..,i_n\{\mathrm{1..2}N\}`$ such that:
$$\frac{A}{S_{ik}}=\frac{^nA}{O_{i1}\mathrm{}O_{in}}0,$$
(6)
for some classical observable $`A=O_i(t),i=\mathrm{1..2}N`$. Let also $`S_{ik}^m`$ be an array of $`m`$ arbitrary sequences $`S_{ik}`$. Moreover, we define $`\delta _{S_{ik}^m}`$ to be the product of all error margins associated with the observables included in $`S_{ik}^m`$. We proposed the set of relations,
$$<E_{S_{ik}^m}|E_{S_{ik}^m}>\delta _{S_{ik}^m}^2,mM,$$
(7)
as a classicality criterion. Notice that given the classical initial data $`(O_i^0,\delta _i)`$, the set of inequalities (7) constitute a set of conditions (classicality conditions) on the functional form of the initial data wave function $`|\varphi >`$. If a certain dynamical system with some given initial data satisfies these relations up to order $`m=M`$, then we say that the system is $`M`$-order classical. In this case it is straightforward to obtain from (5) that:
$$<E_A^M|E_A^M>\left(\underset{i=1}{\overset{2N}{}}\left|\frac{A}{O_i}\right|_{O_i=O_i^0}\delta _i+\mathrm{}\right)^{2M}=\delta _A^{2M}$$
(8)
and thus we can state that if $`|\varphi >`$ satisfies (7) up to order $`M`$ then not only $`|\varphi >`$ is $`M`$-order consistent with the initial time classical description, but also all future classical and quantum descriptions of the configuration of the system are $`M`$-order consistent.
3) The key step to obtain the classicality criterion was the derivation of the relation (4). Given $`\widehat{A}`$ it is clear that this relation is not valid for an arbitrary $`A`$. In fact $`A`$ should be obtained from $`\widehat{A}`$ by following a well defined procedure that was named dequantization. In , we proved that if this dequantization is the inverse of the Dirac quantization map , then the general relation (4) is not exactly valid, being however a very good approximation. More precisely, we proved that for $`A`$ obtained from the operator $`\widehat{A}`$ by applying the Dirac dequantization map, the difference between the right- and left-hand sides of eq.(4) is, at most proportional to a factor of $`\mathrm{}^2`$.
On the other hand, if $`\widehat{A}`$ in eq.(4) is the time-evolution of some initial observable $`\widehat{O}_i(0)`$ (i.e. $`\widehat{A}=\widehat{O}_i(t)`$ ), then $`A`$, the observable obtained by applying the Dirac dequantization map to $`\widehat{A}`$, is just the time evolution of the classical observable $`O_i(0)`$, that is $`A=O_i(t)`$. To proceed, let us define the semiclassical limit of the quantum operator $`\widehat{A}`$ to be the classical observable $`A_S`$ that fully validates the expansion (4). We can now restate our statement made in the previous paragraph in the following terms: The observable $`O_i(t)`$ is not the semiclassical limit of the quantum observable $`\widehat{O}_i(t)`$ (i.e. $`O_i(t)A_S`$) and hence does not allow for the exact statement that if the classical and quantum initial data satisfy the set of relations (7), then the classical and quantum predictions will be $`M`$-order consistent at all times (8). Still, $`O_i(t)`$ provides a prediction that diverges from the correct semiclassical limit one by a factor proportional to $`\mathrm{}^2`$ at most (i.e. $`O_i(t)=A_S+𝒪(\mathrm{}^2)`$). This means that the standard classical prediction is well inside the error interval associated with the semiclassical limit prediction.
Consequently, the results of allow us to conclude that: i) Classical mechanics does not provide the exact semiclassical limit for quantum systems with a set of initial data satisfying (7). ii) Still, classical mechanics provides a very similar prediction to the exact semiclassical limit one.
In this paper we want to obtain the exact form of the classical observables that fully validates the expansion (4). More important, the aim is to derive a consistent dynamical framework able to provide those semiclassical limit observables directly from the classical initial data.
## 3 Symmetric Dequantization
The complete set of fundamental hermitian operators $`S=\{\widehat{O}_i,i=1,\mathrm{},2N\}`$ spans the algebra $`\widehat{𝒜}`$ of operators acting on the Hilbert space $``$,
$$\widehat{𝒜}\{\widehat{A}:\widehat{A}=\underset{i=1}{\overset{n}{}}c_i\underset{j=1}{\overset{m}{}}\widehat{O}_{ij};n,m𝒩,c_i𝒞\},$$
(9)
where $`\widehat{O}_{ij}S`$. Consequently, the set
$$\widehat{}\left\{\widehat{O}_{i_1i_2\mathrm{}i_k}\widehat{O}_{i_1}\widehat{O}_{i_2}\mathrm{}\widehat{O}_{i_k};1i_1,i_2,\mathrm{},i_k2N,k𝒩\right\}$$
(10)
generates all the elements of $`\widehat{𝒜}`$. We stress the fact that the order of the operators in the product $`\widehat{O}_{i_1i_2\mathrm{}i_k}`$ is meaningful.
$`\widehat{𝒜}`$ is thus an infinite dimensional complex vector space together with the bracket rule,
$$[\widehat{A},\widehat{B}]\widehat{A}\widehat{B}\widehat{B}\widehat{A},\widehat{A},\widehat{B}\widehat{𝒜}.$$
(11)
The set of fundamental operators $`S`$ is decomposed into $`N`$ pairs of variables, $`(\widehat{q}_i,\widehat{p}_i)`$, with $`i=1,\mathrm{},N`$ and:
$$\begin{array}{cc}\widehat{O}_i=\widehat{q}_i,\hfill & i=1,\mathrm{},N,\hfill \\ \widehat{O}_i=\widehat{p}_{iN},\hfill & i=N+1,\mathrm{},2N.\hfill \end{array}$$
(12)
These pairs are canonically conjugate in the sense that:
$$\begin{array}{ccc}[\widehat{q}_i,\widehat{q}_j]=\hfill & [\widehat{p}_i,\widehat{p}_j]=0;\hfill & \hfill [\widehat{q}_i,\widehat{p}_j]=i\mathrm{}\delta _{ij}.\end{array}$$
(13)
Let $`T^{}M`$ be the phase space of the corresponding classical system and $`𝒜`$ the algebra of classical observables, $`𝒜\{fC^{\mathrm{}}:T^{}M𝒞\}`$. Let us define the dequantization map, $`V`$: $`\widehat{𝒜}𝒜`$, such that $`V\mathrm{\Lambda }=1`$, where $`\mathrm{\Lambda }`$ is the Dirac quantization map . This map attributes a classical variable to every quantum operator.
Let then $`A=V(\widehat{A})`$. We saw in that if $`A`$ is obtained from $`\widehat{A}`$ using the map $`V`$ then, in general, the expansion (4) is not exactly valid, the right and left hand sides displaying a difference of the order of $`\mathrm{}^2`$ at most. Consequently, all subsequent results, including eq.(8) display an imprecision of the order of $`\mathrm{}^2`$. Therefore the conclusion that if a given set of classical and quantum initial data are $`M`$-order classical, then their time evolution will always be $`M`$-order consistent, is only valid up to a correction - of the classical prediction - by a term proportional to $`\mathrm{}^2`$. Moreover, the map $`V`$ is beset by order problems: it is neither univocous nor injective, as illustrated by the two following examples: let $`\widehat{q}`$, $`\widehat{p}`$ be the two fundamental operators of a one dimensional quantum system and consider the two following hermitian operators: $`\widehat{A}=\widehat{q}\widehat{p}^2\widehat{q}`$ and $`\widehat{B}=1/2(\widehat{q}^2\widehat{p}^2+\widehat{p}^2\widehat{q}^2)`$. Clearly, $`\widehat{A}\widehat{B}`$ and yet one possible Dirac dequantization yields:
$$V(\widehat{A})=q^2p^2\text{and}V(\widehat{B})=q^2p^2$$
and thus the map $`V`$ is not injective. Using the same example we also have: $`\widehat{A}=\widehat{B}+\mathrm{}^2`$ and yet:
$$V(\widehat{A})=q^2p^2V(\widehat{B}+\mathrm{}^2)=q^2p^2+\mathrm{}^2$$
and so the map $`V`$ is not univocous either.
We see that there are many different ways of dequantizing a system. This should come as no surprise since the Dirac quantization map is itself not one to one. As a consequence the classical algebra $`𝒜`$ which is spanned by:
$$\left\{O_{i_1i_2\mathrm{}i_k}O_{i_1}O_{i_2}\mathrm{}O_{i_k};1i_1,i_2,\mathrm{},i_k2N,k𝒩\right\},$$
(14)
where this time the order of the observables is immaterial, shows no straightforward relation with $`\widehat{}`$.
The whole framework described in the last section (and in particular the expansion (4)) is exactly valid if the classical observables $`A`$ are obtained (through a trivial substitution of the fundamental observables $`\widehat{O}_i`$ by the classical ones $`O_i`$) from a fully symmetric form of $`\widehat{A}`$ . Consequently, the aim now is to present a dequantization map that, given an arbitrary quantum observable $`\widehat{A}`$, yields the classical observable $`A_S`$.
Let us thus introduce the ”symmetric dequantization” prescription. Any operator in the algebra $`\widehat{𝒜}`$ can be cast in the form of a linear combination of fully symmetrized polynomia of the fundamental operators. Consider e.g. $`\widehat{p}\widehat{q}`$:
$$\widehat{p}\widehat{q}=\frac{1}{2}(\widehat{p}\widehat{q}+\widehat{q}\widehat{p})+\frac{1}{2}[\widehat{p},\widehat{q}](\widehat{p}\widehat{q})_++\frac{1}{2}[\widehat{p},\widehat{q}],$$
(15)
where $`(\widehat{p}\widehat{q})_+1/2(\widehat{p}\widehat{q}+\widehat{q}\widehat{p})`$ is the completely symmetrized product. Since $`[\widehat{p},\widehat{q}]=i\mathrm{}`$ is a c-number, we have achieved the expansion of $`\widehat{p}\widehat{q}`$ in the basis of the fully symmetrized operators.
H.Weyl was the first to suggest that any operator could be expanded as a sum of completely symmetric terms, . He proposed the following rule. A given operator $`\widehat{b}(\widehat{\stackrel{}{q}},\widehat{\stackrel{}{p}})`$ is represented in the form:
$$\widehat{b}(\widehat{\stackrel{}{q}},\widehat{\stackrel{}{p}})=𝑑\stackrel{}{x}𝑑\stackrel{}{y}\beta (\stackrel{}{x},\stackrel{}{y})\mathrm{exp}\left(i\stackrel{}{x}\widehat{\stackrel{}{q}}+i\stackrel{}{y}\widehat{\stackrel{}{p}}\right),$$
where $`\stackrel{}{x}`$ and $`\stackrel{}{y}`$ are two $`N`$-dimensional vectors whose components are c-numbers, $`\widehat{\stackrel{}{q}}=(\widehat{q}_1,\mathrm{},\widehat{q}_N)`$, $`\widehat{\stackrel{}{p}}=(\widehat{p}_1,\mathrm{},\widehat{p}_N)`$ and $`\beta (\stackrel{}{x},\stackrel{}{y})=\beta ^{}(\stackrel{}{x},\stackrel{}{y})`$ is some numerical function (possibly singular). The previous condition ensures the hermiticity of the operator. If the position $`\widehat{q}_j`$ appears $`n`$ times and the momentum $`\widehat{p}_k`$ appears $`m`$ times in a given operator, then we include terms of the form<sup>3</sup><sup>3</sup>3$`\delta ^{(n)}(x)`$ is the $`n`$-th derivative of the delta-function with respect to its argument. $`i^n\delta ^{(n)}(x_j)`$ and $`i^m\delta ^{(m)}(y_k)`$, respectively, in the function $`\beta (\stackrel{}{x},\stackrel{}{y})`$. We equally include factors $`\delta (x_l)`$ and $`\delta (y_s)`$ for all $`\widehat{q}_l`$ and $`\widehat{p}_s`$ that are absent in the operator.
As an example, consider a system with one degree of freedom ($`N=1`$) and some operator where $`\widehat{q}`$ appears twice and $`\widehat{p}`$ once. From the previous rule we get the completely symmetric operator:
$$\widehat{b}(\widehat{q},\widehat{p})=𝑑x𝑑y\left[i^2\delta ^{^{\prime \prime }}(x)\right]\left[i\delta ^{^{}}(y)\right]\mathrm{exp}\left(ix\widehat{q}+iy\widehat{p}\right)=\frac{1}{3}\left(\widehat{q}^2\widehat{p}+\widehat{q}\widehat{p}\widehat{q}+\widehat{p}\widehat{q}^2\right).$$
There is a clear advantage in using the symmetric version of quantum mechanics. The set $`\widehat{}_+`$ of symmetrized products,
$$\widehat{}_+\left\{\left(\widehat{O}_{i_1i_2\mathrm{}i_k}\right)_+\left(\widehat{O}_{i_1}\widehat{O}_{i_2}\mathrm{}\widehat{O}_{i_k}\right)_+;1i_1,i_2,\mathrm{},i_k2N,k𝒩\right\},$$
(16)
constitutes a basis for $`\widehat{𝒜}`$. In fact it is easy to check that its elements are linearly independent and moreover that all elements of $`\widehat{𝒜}`$ can be expanded in terms of the elements of $`\widehat{}_+`$. On the other hand, in the completely symmetric product $`\left(\widehat{O}_{i_1}\widehat{O}_{i_2}\mathrm{}\widehat{O}_{i_k}\right)_+`$ the order of the operators is immaterial and therefore $`\widehat{}_+`$ and $``$ have the same number of elements. As a side remark we conclude that $`\widehat{}`$ is overcomplete. We can thus define the dequantization map $`V_S`$ in the following way:
Definition 2 \- Symmetric Dequantization
The dequantization map $`V_S:`$ $`\widehat{𝒜}𝒜`$ is defined by the following rules:
1) $`V_S`$ is a linear map;
2) $`V_S`$ maps the identity to the identity;
3) $`V_S\left(\left(\widehat{O}_{i_1}\widehat{O}_{i_2}\mathrm{}\widehat{O}_{i_k}\right)_+\right)=O_{i_1}O_{i_2}\mathrm{}O_{i_k},\text{for all }\left(\widehat{O}_{i_1}\widehat{O}_{i_2}\mathrm{}\widehat{O}_{i_k}\right)_+\widehat{}_+`$.
Let us study some of the properties of $`V_S`$:
a) For a generic operator $`\widehat{A}`$ we have:
$$V_S(\widehat{A})A=V_S\left(\underset{i=1}{\overset{n}{}}c_i\left(\underset{j=1}{\overset{m}{}}\widehat{O}_{ij}\right)_+\right)=\underset{i=1}{\overset{n}{}}c_i\underset{j=1}{\overset{m}{}}O_{ij},$$
(17)
where $`_{i=1}^nc_i\left(_{j=1}^m\widehat{O}_{ij}\right)_+`$ is the expansion of the operator $`\widehat{A}`$ in the basis $`\widehat{}_+`$.
b) $`V_S\left([\widehat{O}_i,\widehat{O}_j]\right)=i\mathrm{}\{O_i,O_j\}`$, where $`\{,\}`$ is the Poisson bracket defined by:
$$\{A,B\}\underset{i=1}{\overset{N}{}}\left[\frac{A}{q_i}\frac{B}{p_i}\frac{A}{p_i}\frac{B}{q_i}\right],A,B𝒜$$
(18)
c) From the expansion (17) and the definition of $`V_S`$, we conclude that:
$$V_S(\widehat{A}^{})=\left[V_S(\widehat{A})\right]^{}A^{}.$$
(19)
d) Since the two basis $`\widehat{}_+`$ and $``$ have the same number of elements, the map $`V_S`$ is bijective and univocous.
e) From property a), we see that $`A=V_S(\widehat{A})`$ is the appropriate classical observable to be used in the expansion (4) of section 2. In that case the expansion (4) is exactly valid. Consequently, all the results of section 2 concerning the consistency between the classical and the quantum description of the time evolution of a general observable $`\widehat{A}`$ become exactly valid.
f) In general $`V_S(\widehat{A}\widehat{B})V_S(\widehat{A})V_S(\widehat{B})`$ and $`V_S\left([\widehat{A},\widehat{B}]\right)i\mathrm{}\{V_S(\widehat{A}),V_S(\widehat{B})\}`$. Let us then define a new product and a new bracket in $`𝒜`$:
Definition 3 \- The product $``$ and the bracket $`[,]_M`$
The new product is the map:
$$:𝒜\times 𝒜𝒜;V_S(\widehat{A})V_S(\widehat{B})=V_S(\widehat{A}\widehat{B}),\widehat{A},\widehat{B}\widehat{𝒜}.$$
(20)
Using this product we can define the new classical bracket:
$$[,]_M:𝒜\times 𝒜𝒜;[A,B]_M=ABBA,$$
(21)
and it is straightforward to check that the new bracket satisfies the identity: $`[V_S(A),V_S(B)]_M=V_S\left([\widehat{A},\widehat{B}]\right)`$.
The aim of the next section is to derive the explicit formula of the product $`V_S(\widehat{A}\widehat{B})`$ and of the commutator $`V_S([\widehat{A},\widehat{B}])`$ and, in the sequel, study their properties. We will see that the product $``$ is well defined: $`(𝒜,+,)`$ is a ring and the bracket $`[,]_M`$ is a true Lie bracket. We can thus anticipate the last property of the symmetric dequantization map:
g) The dequantization map $`V_S`$ is an isomorphism between the Lie algebras $`(\widehat{𝒜},,[,])`$ and $`(𝒜,,[,]_M)`$.
## 4 Dynamical structure of Symmetric Classical Mechanics
The purpose of this section is thus to derive an explicit formula for the dequantization of the product $`V_S(\widehat{A}\widehat{B})`$ of two operators $`\widehat{A},\widehat{B}\widehat{𝒜}`$ and, as a by-product, of the commutator $`V_S([\widehat{A},\widehat{B}])`$. This will allow us to establish the canonical structure of symmetric classical mechanics and therefore to make predictions about the evolution of an arbitrary classical system without having to refer to its quantum formulation.
### 4.1 The product $``$
First of all, notice that performing a symmetrization in the sense of section 3, coincides with the procedure of normal ordering bosonic fields by applying Wick’s theorem, providing we define the ”propagator”:
$$<\widehat{p}_i\widehat{q}_j>=<\widehat{q}_j\widehat{p}_i>\frac{1}{2}[\widehat{p}_i,\widehat{q}_j]=i\frac{\mathrm{}}{2}\delta _{ij}.$$
(22)
The following example illustrates this method:
$$\begin{array}{c}\widehat{O}_i\left(\widehat{O}_j\widehat{O}_k\right)_+=\left(\widehat{O}_i\widehat{O}_j\widehat{O}_k\right)_++<\widehat{O}_i\widehat{O}_j>\widehat{O}_k+<\widehat{O}_i\widehat{O}_k>\widehat{O}_j=\\ \\ =\left(\widehat{O}_i\widehat{O}_j\widehat{O}_k\right)_++\frac{1}{2}[\widehat{O}_i,\widehat{O}_j]\widehat{O}_k+\frac{1}{2}[\widehat{O}_i,\widehat{O}_k]\widehat{O}_j,\end{array}$$
where $`\left(\widehat{O}_i\widehat{O}_j\widehat{O}_k\right)_+\frac{1}{3!}\left(\widehat{O}_i\widehat{O}_j\widehat{O}_k+\text{permutations}\right)`$. We stress that this identity only holds, because $`[\widehat{O}_i,\widehat{O}_j]`$ is a c-number.
Let us now try to obtain $`V_S(\widehat{A}\widehat{B})`$ and $`V_S([\widehat{A},\widehat{B}])`$ for generic operators $`\widehat{A},\widehat{B}\widehat{𝒜}`$, given $`A=V_S(\widehat{A})`$, $`B=V_S(\widehat{B})`$. If we assume that $`\widehat{A}`$ and $`\widehat{B}`$ are already completely symmetrized, then we get (cf.(4)):
$$\{\begin{array}{c}\widehat{A}=_{n=0}^{\mathrm{}}\frac{1}{n!}_{1i_1,\mathrm{},i_n2N}\frac{^nA}{O_{i_1}\mathrm{}O_{i_n}}\left(\widehat{M}_{i_1}\mathrm{}\widehat{M}_{i_n}\right)_+\hfill \\ \\ \widehat{B}=_{m=0}^{\mathrm{}}\frac{1}{m!}_{1j_1,\mathrm{},j_m2N}\frac{^mB}{O_{j_1}\mathrm{}O_{j_m}}\left(\widehat{M}_{j_1}\mathrm{}\widehat{M}_{j_m}\right)_+,\hfill \end{array}$$
(23)
where the monomials $`\widehat{M}_i`$ are defined by: $`\widehat{M}_i\widehat{O}_iO_i`$.
The subtlety resides in noticing that in general $`(\widehat{A}\widehat{B})_+1/2(\widehat{A}\widehat{B}+\widehat{B}\widehat{A})`$. So the whole problem reduces to symmetrizing $`\widehat{A}\widehat{B}`$ properly. If we carry out the product of the two expansions (23), we shall have to dequantize terms of the form:
$$V_S\left((\widehat{M}_{i_1}\mathrm{}\widehat{M}_{i_n})_+(\widehat{M}_{j_1}\mathrm{}\widehat{M}_{j_m})_+\right).$$
(24)
Using Wick’s theorem we have:
$$\begin{array}{c}(\widehat{M}_{i_1}\mathrm{}\widehat{M}_{i_n})_+(\widehat{M}_{j_1}\mathrm{}\widehat{M}_{j_m})_+=\left(\widehat{M}_{i_1}\mathrm{}\widehat{M}_{i_n}\widehat{M}_{j_1}\mathrm{}\widehat{M}_{j_m}\right)_++\text{terms with one contraction}+\\ \\ +\text{terms with two contractions}+\mathrm{}+\text{terms with }min\{n,m\}\text{ contractions}.\end{array}$$
(25)
On the other hand,
$$\begin{array}{c}V_S\left((\widehat{M}_{i_1}\mathrm{}\widehat{M}_{i_k})_+\right)=V_S(\widehat{M}_{i_1})\mathrm{}V_S(\widehat{M}_{i_k})=V_S(\widehat{O}_{i_1}O_{i_1})\mathrm{}V_S(\widehat{O}_{i_k}O_{i_k})=\\ \\ =(O_{i_1}O_{i_1})\mathrm{}(O_{i_k}O_{i_k})=0\end{array}$$
Consequently, if $`nm`$ then the term (24), (25) will yield a vanishing contribution to $`V_S(\widehat{A}\widehat{B})`$. We are left with:
$$\begin{array}{c}V_S(\widehat{A}\widehat{B})=_{n=0}^{\mathrm{}}\frac{1}{(n!)^2}_{1i_1,\mathrm{},i_n2N}_{1j_1,\mathrm{},j_n2N}\frac{^nA}{O_{i_1}\mathrm{}O_{i_n}}\frac{^nB}{O_{j_1}\mathrm{}O_{j_n}}\times \\ \\ \times V_S[<\widehat{M}_{i_1}\widehat{M}_{j_1}><\widehat{M}_{i_2}\widehat{M}_{j_2}>\mathrm{}<\widehat{M}_{i_n}\widehat{M}_{j_n}>+\\ \\ +<\widehat{M}_{i_1}\widehat{M}_{j_2}><\widehat{M}_{i_2}\widehat{M}_{j_1}>\mathrm{}<\widehat{M}_{i_n}\widehat{M}_{j_n}>+\text{permutations}],\end{array}$$
where the expression inside the bracket includes all the terms with $`n`$ contractions. Since the derivative $`(^nB)/(O_{j_1}\mathrm{}O_{j_n})`$ is symmetric with respect to swapping any two indices, we conclude that all the terms yield the same contribution. There are $`n!`$ permutations of the $`j`$ indices. We hence get:
$$\begin{array}{c}V_S(\widehat{A}\widehat{B})=_{n=0}^{\mathrm{}}\frac{1}{n!}_{1i_1,\mathrm{},i_n2N}_{1j_1,\mathrm{},j_n2N}\frac{^nA}{O_{i_1}\mathrm{}O_{i_n}}\frac{^nB}{O_{j_1}\mathrm{}O_{j_n}}\times \\ \\ \times V_S\left(\frac{1}{2^n}[\widehat{M}_{i_1},\widehat{M}_{j_1}][\widehat{M}_{i_2},\widehat{M}_{j_2}]\mathrm{}[\widehat{M}_{i_n},\widehat{M}_{j_n}]\right).\end{array}$$
(26)
Notice that $`V_S\left([\widehat{M}_i,\widehat{M}_j]\right)=V_S\left([\widehat{O}_iO_i,\widehat{O}_jO_j]\right)=i\mathrm{}\{O_i,O_j\}`$. Let us now define the derivative $`\stackrel{}{}`$ which obeys the antisymmetric Leibnitz rule:
$$\begin{array}{cc}\stackrel{}{}(AB)=\hfill & (A)BAB,\hfill \\ \stackrel{}{}\{A,B\}=\hfill & \{A,B\}\{A,B\}.\hfill \end{array}$$
(27)
We equally define the following ”Liouvillian” operator;
$$\widehat{}\frac{1}{2}\underset{i=1}{\overset{N}{}}\left(\frac{}{p_i}\frac{\stackrel{}{}}{q_i}\frac{}{q_i}\frac{\stackrel{}{}}{p_i}\right).$$
(28)
We represent eq.(26) in the form,
$$V_S(\widehat{A}\widehat{B})=\underset{n=0}{\overset{\mathrm{}}{}}_n.$$
(29)
Let us now prove by induction that:
Lemma
$$_n\frac{1}{n!}\left(\frac{i\mathrm{}}{2}\right)^n\widehat{}^nAB.$$
(30)
Proof:
From (26) we get $`_0=AB`$ and
$$_1=\frac{i\mathrm{}}{2}\underset{i,j=1}{\overset{2N}{}}\frac{A}{O_i}\frac{B}{O_j}\{O_i,O_j\}=\frac{i\mathrm{}}{2}\underset{i=1}{\overset{N}{}}\left(\frac{A}{q_i}\frac{B}{p_i}\frac{A}{p_i}\frac{B}{q_i}\right)=\frac{i\mathrm{}}{2}\{A,B\}.$$
It is easy to check that:
$$\{A,B\}=\widehat{}AB.$$
(31)
And so $`_1=i\mathrm{}/2\widehat{}AB`$, in agreement with (29), (30).
Let us now assume that (30) holds for some $`n`$. We then have from (26):
$$\begin{array}{c}_{n+1}=\frac{1}{(n+1)!}_{1i_1,\mathrm{},i_{n+1}2N}_{1j_1,\mathrm{},j_{n+1}2N}\left(\frac{i\mathrm{}}{2}\right)^{n+1}\frac{^{n+1}A}{O_{i_1}\mathrm{}O_{i_{n+1}}}\times \\ \\ \times \frac{^{n+1}B}{O_{j_1}\mathrm{}O_{j_{n+1}}}\times \{O_{i_1},O_{j_1}\}\mathrm{}\{O_{i_{n+1}},O_{j_{n+1}}\}=\\ \\ =\frac{1}{(n+1)!}\left(\frac{i\mathrm{}}{2}\right)^{n+1}_{1i_1,\mathrm{},i_n2N}_{1j_1,\mathrm{},j_n2N}\{O_{i_1},O_{j_1}\}\mathrm{}\{O_{i_n},O_{j_n}\}\times \\ \\ \times \{\frac{^nA}{O_{i_1}\mathrm{}O_{i_n}},\frac{^nB}{O_{j_1}\mathrm{}O_{j_n}}\}=\\ \\ =\frac{1}{(n+1)!}\left(\frac{i\mathrm{}}{2}\right)^{n+1}_{1i_1,\mathrm{},i_n2N}_{1j_1,\mathrm{},j_n2N}\{O_{i_1},O_{j_1}\}\mathrm{}\{O_{i_n},O_{j_n}\}\times \\ \\ \times \left[\widehat{}\left(\frac{^nA}{O_{i_1}\mathrm{}O_{i_n}}\frac{^nB}{O_{j_1}\mathrm{}O_{j_n}}\right)\right].\end{array}$$
In the last step we used (31). Since $`\{O_i,O_j\}`$ are c-numbers, they commute with the operator $`\widehat{}`$, and we get:
$$\begin{array}{c}_{n+1}=\frac{1}{n+1}\frac{i\mathrm{}}{2}\widehat{}[\frac{1}{n!}\left(\frac{i\mathrm{}}{2}\right)^n_{1i_1,\mathrm{},i_n2N}_{1j_1,\mathrm{},j_n2N}\\ \\ \{O_{i_1},O_{j_1}\}\mathrm{}\{O_{i_n},O_{j_n}\}\times \frac{^nA}{O_{i_1}\mathrm{}O_{i_n}}\frac{^nB}{O_{j_1}\mathrm{}O_{j_n}}]=\\ \\ =\frac{1}{n+1}\frac{i\mathrm{}}{2}\widehat{}\left[_n\right]=\frac{1}{(n+1)!}\left(\frac{i\mathrm{}}{2}\right)^{n+1}\widehat{}^{n+1}AB,\end{array}$$
in agreement with (30). In summary, we proved the following theorem:
Theorem 1: Let $`\widehat{A},\widehat{B}\widehat{𝒜}`$ and $`A=V_S(\widehat{A}),B=V_S(\widehat{B})`$. The dequantization of the product $`\widehat{A}\widehat{B}`$ is given by:
$$V_S(\widehat{A}\widehat{B})AB=\mathrm{exp}\left(\frac{i\mathrm{}}{2}\widehat{}\right)AB.$$
(32)
This product is more commonly found in the form:
$$AB=A\mathrm{exp}\left(\frac{i\mathrm{}}{2}\widehat{𝒥}\right)B,$$
(33)
where $`\widehat{𝒥}`$ is the Janus operator:
$$\widehat{𝒥}\underset{i=1}{\overset{N}{}}\left(\frac{\stackrel{}{}}{q_i}\frac{\stackrel{}{}}{p_i}\frac{\stackrel{}{}}{p_i}\frac{\stackrel{}{}}{q_i}\right).$$
(34)
This product has the following properties for ($`A,B,C𝒜;`$ $`a,b𝒞)`$:
$$\begin{array}{cc}1)\hfill & \text{Linearity: }(aA+bB)C=a(AC)+b(BC),\hfill \\ 2)\hfill & A_{\mathrm{}}B=B_{\mathrm{}}A,\hfill \\ 3)\hfill & \text{Associativity: }(AB)C=A(BC),\hfill \\ 4)\hfill & \text{Identity: }A1=1A=A,\hfill \\ 5)\hfill & p_ip_j=p_ip_j;q_iq_j=q_iq_j;q_ip_j=q_ip_j+\frac{i\mathrm{}}{2}\delta _{ij}\hfill \\ 6)\hfill & (AB)^{}=B^{}A^{}\hfill \end{array}$$
(35)
Property 2) means that changing the order of the variables $`A`$, $`B`$ is tantamount to performing the substitution $`\mathrm{}\mathrm{}`$ in formula (33). This can be proved immediately by substituting the identity:
$$B\widehat{𝒥}^nA=(1)^nA\widehat{𝒥}^nB,$$
(36)
in equation (30). Property 3) is a trivial consequence of dequantizing the product of three operators. The remaining properties are straightforward to prove using the formula (33).
### 4.2 The bracket $`[,]_M`$
Theorem 2: The classical bracket $`[A,B]_MV_S([\widehat{A},\widehat{B}])`$ is given by:
$$[A,B]_M=2iA\mathrm{sin}\left(\frac{\mathrm{}}{2}\widehat{𝒥}\right)B,$$
(37)
for any $`\widehat{A},\widehat{B}\widehat{𝒜}`$ and $`A=V_S(\widehat{A})`$, $`B=V_S(\widehat{B})`$.
Proof:
$$\begin{array}{c}[A,B]_MV_S([\widehat{A},\widehat{B}])=V_S(\widehat{A}\widehat{B})V_S(\widehat{B}\widehat{A})=ABBA=A\left(e^{\frac{i\mathrm{}}{2}\widehat{𝒥}}e^{\frac{i\mathrm{}}{2}\widehat{𝒥}}\right)B=\\ \\ =_{n=0}^{\mathrm{}}\frac{1}{n!}\left(\frac{i\mathrm{}}{2}\right)^n\left[1(1)^n\right]A\widehat{𝒥}^nB=2iA\mathrm{sin}\left(\frac{\mathrm{}}{2}\widehat{𝒥}\right)B,\end{array}$$
where we used eq.(21) and property 2) in eq.(33). This expression is the celebrated Moyal bracket, . This formula might appear awkward at first sight. One would expect $`V_S([\widehat{A},\widehat{B}])=i\mathrm{}\{A,B\}`$. To order $`\mathrm{}`$, we have:
$$\begin{array}{cc}V_S(\widehat{A}\widehat{B})=\hfill & A\left(1+\frac{i\mathrm{}}{2}\widehat{𝒥}+𝒪(\mathrm{}^2)\right)B=AB+\frac{i\mathrm{}}{2}\{A,B\}+𝒪(\mathrm{}^2),\hfill \\ & \\ V_S([\widehat{A},\widehat{B}])=\hfill & 2iA\left(\frac{\mathrm{}}{2}\widehat{𝒥}+𝒪(\mathrm{}^3)\right)B=i\mathrm{}\{A,B\}+𝒪(\mathrm{}^3).\hfill \end{array}$$
(38)
To this order we do indeed recover the Poisson bracket. The Moyal bracket has the following properties for $`A,B,C𝒜`$; $`a,b𝒞`$:
$$\begin{array}{cc}1)\hfill & \text{Linearity: }[aA+bB,C]_M=a[A,C]_M+b[B,C]_M,\hfill \\ 2)\hfill & \text{Antisymmetry: }[B,A]_M=[A,B]_M,\hfill \\ 3)\hfill & \text{Jacobi identity: }[[A,B]_M,C]_M+[[B,C]_M,A]_M+[[C,A]_M,B]_M=0,\hfill \\ 4)\hfill & \text{Product (Leibnitz) rule: }[AB,C]_M=A[B,C]_M+[A,C]_MB,\hfill \\ 5)\hfill & \text{Structure constants :}[q_i,q_j]_M=[p_i,p_j]_M=0;[q_i,p_j]_M=i\mathrm{}\delta _{ij}.\hfill \end{array}$$
(39)
All these results follow immediately from the properties (35) of the product $``$ and from (21).
### 4.3 Dynamical Evolution
The time evolution of a quantum operator $`\widehat{A}(t)`$ is given by:
$$\widehat{A}(t)=\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{n!}\left(\frac{it}{\mathrm{}}\right)^n[\widehat{H},[\widehat{H},[\mathrm{},[\widehat{H},\widehat{A}]\mathrm{}]]],$$
(40)
where $`\widehat{H}`$ is the Hamiltonian. If $`V_S(\widehat{H})=H`$, then the previous equation yields upon dequantization:
$$A(t)=\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{n!}\left(\frac{it}{\mathrm{}}\right)^n[H,[H,[\mathrm{},[H,A]_M\mathrm{}]_M]_M]_M,$$
(41)
As a consequence of (41), the observable $`A`$ obeys the differential equation:
$$\dot{A}(t)=\frac{i}{\mathrm{}}[H,A(t)]_M.$$
(42)
Notice that this equation could also be obtained by dequantizing the original quantum dynamical equation for the observable $`\widehat{A}(t)`$. In particular we have:
$$\{\begin{array}{cc}\dot{q}_i(t)=\frac{H}{p_i}\hfill & \\ & \hfill i=1,\mathrm{},N.\\ \dot{p}_i(t)=\frac{H}{q_i}\hfill & \end{array}$$
(43)
Notice that these equations look exactly like the ones in traditional classical mechanics (i.e. with Poisson brackets). However, this similarity is misleading. Indeed, $`H`$ does not look exactly like the traditional classical Hamiltonian. Rather, all products of variables are replaced by the product $``$.
Let us now study some of the properties of the new theory:
1) Time evolution is generated by an unitary transformation. From the quantum theory, we have:
$$\{\begin{array}{c}\widehat{A}(t)=\widehat{U}(t)^1\widehat{A}(0)\widehat{U}(t),\hfill \\ \\ i\mathrm{}\frac{\widehat{U}}{t}=\widehat{H}\widehat{U},\widehat{U}(0)=1.\hfill \end{array}$$
(44)
Applying the map $`V_S`$:
$$\{\begin{array}{c}A(t)=U(t)^1A(0)U(t),\hfill \\ \\ i\mathrm{}\frac{U}{t}=HU,U(0)=1,\hfill \end{array}$$
(45)
where we used the fact that $`V_S(\widehat{U}^1\widehat{U})=V_S(U^1)V_S(\widehat{U})=1`$, and therefore $`V_S(\widehat{U}^1)=V_S(\widehat{U})^1`$.
Moreover, the classical quantity $`U(t)`$ is unitary in the sense that $`U^1=U^{}`$ ($`\widehat{U}^1=\widehat{U}^{}V_S(\widehat{U}^1)=V_S(\widehat{U}^{})V_S(\widehat{U})^1=V_S(\widehat{U})^{}`$) and thus $`U^{}U=1`$.
Substituting (45) into (42), we can verify explicitly that it provides a solution to the equations of motion.
2) All unitary transformations generate canonical transformations. In fact:
$$\begin{array}{c}U^1[A,B]_MU=U^1ABUU^1BAU=\\ \\ =U^1AUU^1BUU^1BUU^1AU=\\ \\ =[U^1AU,U^1BU]_M,\end{array}$$
(46)
and so the bracket structure is preserved under the action of $`U`$. In particular, time evolution is a canonical transformation.
3) The limit $`\mathrm{}0`$ of symmetric classical mechanics is standard classical mechanics. Indeed, the identities,
$$\{\begin{array}{cc}lim_\mathrm{}0AB=AB,\hfill & \\ & \hfill A,B𝒜\\ lim_\mathrm{}0\frac{1}{i\mathrm{}}[A,B]_M=\{A,B\},\hfill & \end{array}$$
(47)
can be checked immediately from the expansions (38). Using these limits in eqs.(41) and (42), we recover the standard version of classical mechanics.
4) Symmetric classical mechanics fully validates the use of expansion (4). In fact, if $`A(0)=V_S\left(\widehat{A}(0)\right))`$, then $`A(t)`$ (obtained by solving (42)) is given by $`V_S\left(\widehat{A}(t)\right)`$. Therefore expansion (4) is exactly valid for $`\widehat{A}(t)A(t)`$, where $`A(t)`$ is the prediction of symmetric classical mechanics for the time evolution of the observable $`A(0)`$, and thus all the results concerning the consistency between the classical and quantum predictions are exactly valid if the classical predictions are those of symmetric classical mechanics.
### 4.4 Symmetric Quantization
Finally, we shall define a quantization prescription for symmetric classical systems:
Definition 4: Symmetric Quantization
The symmetric quantization map $`\mathrm{\Lambda }_S:𝒜\widehat{𝒜}`$ is the Lie algebra isomorphism defined by the following rules:
1) $`\mathrm{\Lambda }_S`$ is linear,
2) $`\mathrm{\Lambda }_S(AB)=\mathrm{\Lambda }_S(A)\mathrm{\Lambda }_S(B).`$
It satisfies the following properties:
1) $`\mathrm{\Lambda }_S`$ maps the identity to the identity:
$$\mathrm{\Lambda }_S(A)=\mathrm{\Lambda }_S(A1)=\mathrm{\Lambda }_S(A)\mathrm{\Lambda }_S(1)\mathrm{\Lambda }_S(1)=1.$$
2) It is the inverse map of $`V_S`$:
$$\mathrm{\Lambda }_S\left(O_{i_1}O_{i_2}\mathrm{}O_{i_k}\right)=\mathrm{\Lambda }_S\left((O_{i_1}O_{i_2}\mathrm{}O_{i_k})_+\right)=\left(\mathrm{\Lambda }_S(O_{i_1})\mathrm{\Lambda }_S(O_{i_2})\mathrm{}\mathrm{\Lambda }_S(O_{i_k})\right)_+$$
where in the last step we used rule 2) from the definition. The previous identity together with property 1) proves that $`\mathrm{\Lambda }_SV_S=1`$.
3) A trivial consequence of rule 2) is the following:
$$\mathrm{\Lambda }_S\left([A,B]_M\right)=[\mathrm{\Lambda }_S(A),\mathrm{\Lambda }_S(B)].$$
## 5 Conclusions
In this paper we presented an alternative formulation of classical physics. The new theory was named symmetric classical mechanics. Its properties were studied thoroughly and, most important, it was shown that symmetric classical mechanics is the exact semiclassical limit of quantum mechanics for an arbitrary quantum system with a set of initial data satisfying the classicality criterion presented in section 2. In other words, the time evolution of a general quantum system is $`M`$-order consistent with the predictions of symmetric classical mechanics, provided the initial data for the two formulations are $`M`$-order classical. Notice that this property is not completely satisfied by standard classical mechanics.
Clearly, symmetric classical mechanics is not the only possible alternative framework for classical mechanics. The entire set of properties of the new theory and in particular its consistent canonical structure are a direct consequence of the definition of the dequantization map or, to go even further, of the choice of a basis for the algebra of quantum observables. Therefore, all the results presented in this paper can be reformulated for other dequantization maps, providing in this fashion other, possibly more interesting, descriptions of classical physics, .
The motivation to develop the theory of symmetric classical mechanics was threefold: 1) The first and most important motivation is theoretical. Symmetric classical mechanics provides a new perspective over the problem of the semiclassical limit of quantum mechanics. Firstly, because it proves that the semiclassical limit might be correctly described (and even more accurately) by another fully consistent dynamical structure and not just by classical mechanics. Secondly, because it clarifies the role of the limit $`\mathrm{}0`$ in deriving the semiclassical limit of quantum mechanics. Symmetric classical mechanics provides a description of the semiclassical limit in which no assumption is made about the magnitude of the Planck constant. Its validity rests exclusively upon a number of conditions that should be satisfied by the initial data wave function. In other words, symmetric classical mechanics would still provide a valid description of dynamics in a world with a huge Planck constant. As a side result we see that standard classical mechanics can be seen as a second limit of quantum mechanics when the set of initial data of the dynamical system satisfies some classicality conditions and the Planck constant can be regarded as being of neglectable magnitude. This result corroborates the argument and the results of . 2) On the other hand, symmetric classical mechanics is formulated in terms of the Moyal bracket and this bracket also provides the dynamics of the Moyal-Weyl-Wigner formulation of quantum mechanics. Therefore, it comes as no surprise that when compared to standard classical mechanics, symmetric classical mechanics displays a clearer relation with quantum mechanics. The quantization map from symmetric classical mechanics to quantum mechanics is one-to-one and there are thus no order problems in the quantization of a symmetric classical system. This property provides a new approach for the analysis of the order ambiguities in quantum mechanics. This analysis can now be enforced at the level of the original classical theory. 3) The last motivation concerns the problem of developing a consistent theory of coupled classical-quantum dynamics . We expect that the symmetric dequantization map might be consistently extended to the case where the purpose is to dequantize only one sector of the original quantum theory. If this is the case and if the extended dequantization map preserves the original set of properties, then it is trivial to obtain a consistent formulation (i.e. a true Lie bracket structure) of hybrid classical-quantum dynamics.
### Acknowledgments
We would like to thank João Marto for several suggestions made through the present work. This work was partially supported by the grants ESO/PRO/1258/98 and CERN/P/Fis/15190/1999.
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# Nonadiabatic noncyclic geometric phase and ensemble average spectrum of conductance in disordered mesoscopic rings with spin-orbit coupling
As is well known, the geometric phase has manifested itself extensively in physics, particularly in mesoscopic systems where quantum interference is extremely important . Recently, Morpurgo et al reported a novel splitting of the main peak (corresponding to the $`hc/e`$ Aharonov-Bohm (AB) oscillations) in the ensemble average Fourier spectrum of the conductance in open mesoscopic rings . The authors conjectured that the observed splitting is due to the spin-orbit(SO)-induced Berry’s phase . It is probably strong experimental evidence showing an important effect of the SO geometric phase on quantum transport. Although it was concluded that, in a mesoscopic ring possessing the time-reversal symmetry in the absence of AB-flux, the SO-dependent transport can be treated formally in the absence of SO-coupling but with an effective magnetic flux , it is unclear yet how to calculate the value of this flux as well as the ensemble average spectrum of conductance; besides, it is not clear either whether the above conclusion is still valid in the presence of an arbitrary local magnetic field on the ring(i.e., the aforementioned time-reversal symmetry is broken), which appears to be the experimental case in . Mal’shukov et al attempted to account for the observed splitting of the main peak, but were not quite successful . Three aspects of the experiment require theoretical explanation. First, the magnitude of the observed splitting is surprisingly large when compared with an estimation based on the adiabatic approximation in a clean mesoscopic ring. For a clean ring with radius $`r`$ subject to a crown-shaped effective magnetic field $`𝐁_{eff}=(B_0cos\phi _r,B_0sin\phi _r,B_z)`$ in the cyclindrical coordinates, Stern predicted that the inverse $`B_z`$ period becomes $`(\mathrm{\Delta }B_z)^1=\pi r^2/\varphi _0\pm 1/(2B_0)`$ at $`B_0>>B_z`$, and the splitting would be $`1.2\times 10^3mT^1`$ for an estimated experimental parameter $`B_00.8T`$, which is at least one order of magnitude less than the observed value. Secondly, the origin of the side structure on the main peak needs to be clarified. Lastly, most existing theoretical estimations are crucially based on the adiabatic or cyclic condition; however, neither adiabatic nor cyclic evolution is well satisfied in the experiment . In view of these facts, we believe that the nonadiabatic noncyclic geometric phase, essentially similar to the effective flux addressed formally in Ref. , could play a crucial role in the system, this is the key point of the present work, being essentially different from some existing theoretical analyses . However, it is still highly nontrival to evaluate the SO-induced geometric phase in the nonadiabatic noncyclic transport and its effect on the ensemble average spectrum.
It is worth pointing out that the relevant geometric phase detected in the experiment is likely induced by a $`U(1)\times SU(2)_{spin}`$ field. In exploring the global geometrical connotations of gauge fields of either the Abelian $`U(1)`$ type or the non-Abelian monopole type, Yang et al showed that the nonintegrable gauge phase factor in the wavefunction gives an intrinsic and complete description of the relevant field, hereafter refered to as Yang’s theory. In this Letter, we first generalize Yang’s theory to the non-Abelian $`U(1)\times SU(2)_{spin}`$ electromagnetic field. Using a simple one-dimensional(1D) continuum model for a quasi-1D mesoscopic ring, we then analyze carefully the nonintegrable phase induced by this field and evaluate its effect on the splitting of the main peak in the ensemble average spectrum of the conductance. Remarkably, we find that the splitting as well as the side structure of the main peak observed by Morpurgo et al stem from the nonadiabatic noncyclic geometric phase.
We consider an electron subject to an electromagnetic field. The corresponding Hamiltonian with $`U(1)_{em}\times SU(2)_{spin}`$ gauge symmetry is given by
$$\widehat{H}=\frac{1}{2m}(𝐩+\frac{e}{c}𝐀\frac{\mu }{c}𝐚)^2eA^0+\mu a^0+V(𝐫),$$
(1)
where $`\mu =g\mu _B/2`$ with $`g`$ the gyromagnetic ratio and $`\mu _B=e\mathrm{}/(2mc)`$ the Bohr magneton. Here $`A^\nu =(A^0,𝐀)`$ represents a $`U(1)_{em}`$ electromagnetic potential, and $`a^\nu =(a^0,𝐚)=(\stackrel{}{\sigma }𝐁,\stackrel{}{\sigma }\times 𝐄/2)`$ is an $`SU(2)_{spin}`$ potential with $`\stackrel{}{\sigma }`$ denoting the Pauli matrix. $`V(𝐫)`$ is an arbitrary spin-independent local potential at the point $`𝐫`$. The Schrödinger equation for the normalized two-component wave function $`\mathrm{\Psi }(x^\nu )`$ reads
$$i\mathrm{}\frac{}{t}\mathrm{\Psi }(x^\nu )=\widehat{H}\mathrm{\Psi }(x^\nu ).$$
(2)
By introducing a new wave function , $`\mathrm{\Psi }_0(x^\nu )=\widehat{U}\mathrm{\Psi }(x^\nu )`$, where
$$\widehat{U}=exp(i\frac{e}{\mathrm{}c}_\mathrm{\Gamma }A_\nu 𝑑x^\nu )\widehat{P}exp(i\frac{\mu }{\mathrm{}c}_\mathrm{\Gamma }a_\nu 𝑑x^\nu )$$
(3)
with $`\widehat{P}`$ the path ordering operator and $`\mathrm{\Gamma }`$ an integration curve from a fixed $`x_0^\nu `$ to $`x^\nu `$, we find that Eq.(2) reduces exactly to
$$i\mathrm{}\frac{}{t}\mathrm{\Psi }_0(x^\nu )=\widehat{H}_0\mathrm{\Psi }_0(x^\nu )$$
(4)
with
$$\widehat{H}_0=\widehat{U}(\widehat{H}i\mathrm{}\frac{}{t})\widehat{U}^1=\frac{(i\mathrm{})^2}{2m}+V(𝐫).$$
(5)
Clearly, $`\widehat{U}`$ is a continuous local gauge transformation. Under this gauge transformation, the Hamiltonian (1) is transformed to a Hamiltonian devoid of electromagnetic fields, but with a phase shift in the wave function as seen in Eq.(3). In this sense, the gauge factor in Eq. (3) is just the nonintegrable phase in Yang’s theory, which can describe completely the $`U(1)\times SU(2)_{spin}`$ electromagnetic field. For a mesoscopic ring where the phase memory is retained by electrons, we may conclude that physical properties of the system in the presence of an electromagnetic field can be expressed in terms of the same quantity in the absence of the electromagnetic field, but with a nonintegrable phase being taken into account. An important application is related to the SO coupling: any spin-independent transport quantity can be expressed in terms of the same quantity in the absence of SO scattering but with an effective magnetic flux, a fact which was shown directly by using the transfer matrix method in a tight-binding form for a mesoscopic ring possessing the time-reversal symmetry in the absence of AB-flux . In fact, with the help of this generalized theory, we are able to study a disordered mesoscopic system subject to an electromagnetic field in a simpler way.
We now focus on the phase factor first. To capture essential physics of geometric phase in the present quasi-1D system, we employ a simple 1D model. For a closed path parameterized by arc length $`s`$, the total phase factor in Eq.(3) is $`\gamma _t=\gamma _{AB}+\stackrel{~}{\gamma }`$, where $`\gamma _{AB}=2\pi \varphi /\varphi _0`$ is the usual AB phase with $`\varphi `$ the magnetic flux and $`\varphi _0=hc/e`$, and $`\stackrel{~}{\gamma }`$ is the second phase factor in Eq.(3), which is determined by a Schrödinger-type equation
$$i\mathrm{}\frac{}{s}|\xi (s)=\mu \stackrel{}{\sigma }(\frac{1}{v}𝐁\frac{1}{2c}\widehat{𝐯}\times 𝐄)|\xi (s).$$
(6)
Here $`\widehat{𝐯}`$ is a unit vector along the direction of the velocity $`𝐯=v\widehat{𝐯}`$ and $`ds=vdt`$. Equation (6) describes the evolution of the spin state $`|\xi `$ governed by the operator $`\widehat{U}`$. The phase associated with Eq.(6) can be further written as $`\stackrel{~}{\gamma }=\gamma _d+\gamma _{AB}^{eff}+\gamma _p`$ with the dynamical phase
$$\gamma _d=\frac{\mu }{\mathrm{}}\xi (s)|\stackrel{}{\sigma }\frac{1}{v}𝐁|\xi (s)ds,$$
(7)
the effective AB phase
$$\gamma _{AB}^{eff}=\frac{\mu }{\mathrm{}}\xi (s)|\frac{1}{2c}\stackrel{}{\sigma }(\widehat{𝐯}\times 𝐄)|\xi (s)ds,$$
(8)
and $`\gamma _p`$ is the Pancharatnam phase, to be addressed in detail later. Here we emphasize that $`\gamma _{AB}^{eff}`$ is a kind of geometric phase, though it seems from Eq.(8) as if it were a dynamical phase related to an ‘effective magnetic field’ $`𝐯\times 𝐄/2c`$. The reason lies in the fact that the two waves propagating in opposite directions in the ring acquire phases with the opposite sign for $`\gamma _{AB}^{eff}`$ (simply because it depends on the velocity direction $`\widehat{𝐯}`$), but the same sign for $`\gamma _d`$ . The geometrical feature of $`\gamma _{AB}^{eff}`$ seems to be ignored in some earlier analyses , which appears to be a minor reason for the existing discrepancy between theory and experiment. In fact, $`\gamma _{AB}^{eff}`$ is just induced by an $`SU(2)_{spin}`$ vector potential $`𝐚`$, and it is clear from Eq.(1) that $`𝐚`$ plays a role similar to that of the $`U(1)_{em}`$ vector potential $`𝐀`$ in the AB effect. As a result, it is expected that an effective AB effect can be induced by this $`SU(2)_{spin}`$ vector potential , as was also shown by Choi et al .
For a unit vector $`𝐧=(n_1,n_2,n_3)=(sin\theta cos\phi ,sin\theta sin\phi ,cos\theta )`$ with $`𝐧`$ a unit sphere $`S^2`$, each $`𝐧`$ corresponds to the spin state $`|\xi =(e^{i\phi /2}cos(\theta /2),e^{i\phi /2}sin(\theta /2))^T`$ via the relation $`𝐧=\xi |\stackrel{}{\sigma }|\xi `$, where $`T`$ represents matrix transposition. The noncyclic Pancharatnam phase accumulated in an evolution of $`𝐧`$ is found to be
$$\gamma _p=\frac{1}{2}_{\mathrm{\Sigma }=C}𝐧𝑑𝚺,$$
(9)
where $`d𝚺`$ is an area element on $`S^2`$, $`C`$ is a specific closed curve on $`S^2`$, which is along the actual path of $`𝐧(s)`$ plus the shorter geodesic curve from the final point $`𝐧(s_f)=(sin\theta _fcos\phi _f,sin\theta _fsin\phi _f,cos\theta _f)`$ to the initial point $`𝐧(0)=(sin\theta _icos\phi _i,sin\theta _isin\phi _i,cos\theta _i)`$. This Pancharatnam phase can be derived as
$$\gamma _p=\frac{1}{2}_0^{t_f}\frac{n_1\dot{n}_2n_2\dot{n}_1}{1+n_3}𝑑t+arctg\frac{sin(\phi _f\phi _i)}{ctg\frac{\theta _f}{2}ctg\frac{\theta _i}{2}+cos(\phi _f\phi _i)},$$
(10)
where $`t_f`$ is the final time, and $`𝐧`$ is determined by the equation
$$\frac{d𝐧}{dt}=\frac{2\mu }{\mathrm{}}(𝐁\frac{1}{2c}𝐯\times 𝐄)\times 𝐧,$$
(11)
which represents a spin-$`\frac{1}{2}`$ particle moving in an effective magnetic field $`(𝐁𝐯\times 𝐄/2c)`$. This phase is not equal to the cyclic Aharonov-Anandan (AA) phase in general , but recovers the AA phase $`\gamma _{AA}=\frac{1}{2}_0^\tau 𝑑t(n_1\dot{n}_2n_2\dot{n}_1)/(1+n_3)`$ for any cyclic evolution with the period $`\tau `$ . It is remarkable that the nonintegrable phase in Eq.(3) can be evaluated by simply computing Eqs.(7), (8), and (10); while it is hard to calculate the value of the effective flux addressed formally in Ref. , particularly in the presence of an arbitrary local magnetic field.
At this stage, as in the experiment , we study a disordered ring with the Rashba SO-interaction(equivalent to an internal electric field $`𝐄=E𝐞_z`$), subject to a local magnetic field $`𝐁=B_z𝐞_z`$ and a magnetic flux $`\varphi =\pi r^2B_z`$. The Hamiltonian, which is in the form of Eq.(1), becomes
$$\widehat{H}=\mathrm{}\omega _r[i\frac{}{\phi _r}+\frac{\varphi }{\varphi _0}\frac{\eta }{2}(\sigma _xcos\phi _r+\sigma _ysin\phi _r)]^2\mu B_z\sigma _z+V(\phi _r),$$
(12)
where $`\omega _r=\mathrm{}/(2mr^2)`$, $`\phi _r`$ is the polar angle, and the normalized electric field strength $`\eta =\mu _BEr/c\mathrm{}=2m\kappa r`$ with the SO coefficient $`\mathrm{}^2\kappa `$.
To account for the experimental results naturally, we investigate the electronic transmission across a disordered ring connected to external current leads, schematically illustrated in Fig.3 in Ref.. In such a system, the electronic transmission is significantly affected by the nonintegrable phase. Using the method originally proposed by Büttiker et al and our generalization of Yang’s theory, the transmission coefficient across the ring is found to be
$$T_g=\frac{ϵ^2}{b^4}\left|(ba,1)\stackrel{~}{T}_+[\frac{e^{i\mathrm{\Delta }\gamma }}{b^2}\left(\begin{array}{cc}(b^2a^2)& a\\ a& 1\end{array}\right)\stackrel{~}{T}_{}\left(\begin{array}{cc}(b^2a^2)& a\\ a& 1\end{array}\right)\stackrel{~}{T}_+\stackrel{~}{1}]^1\left(\begin{array}{c}ba\\ 1\end{array}\right)\right|^2,$$
(13)
where $`\stackrel{~}{T}_+`$ and $`\stackrel{~}{T}_{}`$ are the transfer matrices of the upper and lower branches of the ring, $`\stackrel{~}{1}`$ is the unit matrix, $`a=\pm (\sqrt{12ϵ}1)/2`$, and $`b=\pm (\sqrt{12ϵ}+1)/2`$ with $`0ϵ1/2`$. $`\mathrm{\Delta }\gamma =\gamma _{AB}+\gamma _{AB}^{eff}(𝐧(0))+\gamma _p(𝐧(0))`$ represents the nonadiabatic noncyclic geometrical phase accumulated in the evolution when electron(with the initial spin-state $`𝐧(0)`$) moves one cycle in the clockwise sense . For a beam of electron waves with Fermi wave vector $`k_f`$, the rate for electrons to traverse one round in the ring is $`\omega _f=\mathrm{}k_f/(mr)`$ for ballistic motion, but is estimated approximately to be $`\omega _d=l\omega _f/(2\pi r)`$ for weak diffusive motion , where $`l`$ is the electron mean free path. This rate can be regarded as the angular frequency of the otherwise rotating magnetic field felt by the electron spin , which is given by $`𝐁_{eff}(t)=(B_0^{f,d}cos\omega _{f,d}t,B_0^{f,d}sin\omega _{f,d}t,B_z)`$ with $`B_0^{f,d}=\eta \mathrm{}\omega _{f,d}/2\mu `$. Then from the equation $`d𝐧(t)/dt=(2\mu /\mathrm{})𝐁_{eff}(t)\times 𝐧(t)`$, $`𝐧(t)`$ is derived exactly as
$`𝐧^T(t)`$ $`=`$ $`\left(\begin{array}{ccc}cos\omega _{f,d}t\hfill & sin\omega _{f,d}t& \hfill 0\\ sin\omega _{f,d}t\hfill & cos\omega _{f,d}t& \hfill 0\\ 0\hfill & 0& \hfill 1\end{array}\right)`$ (17)
$`\times `$ $`\left(\begin{array}{ccc}sin^2\chi +cos^2\chi cos\omega _st\hfill & cos\chi sin\omega _st& \hfill \frac{1}{2}sin2\chi (1cos\omega _st)\\ cos\chi sin\omega _st\hfill & cos\omega _st& \hfill sin\chi sin\omega _st\\ \frac{1}{2}sin2\chi (1cos\omega _st)\hfill & sin\chi sin\omega _st& \hfill cos^2\chi +sin^2\chi cos\omega _st\end{array}\right)𝐧^T(0),`$ (21)
where $`\omega _s=\sqrt{\omega _0^2+(\omega _{f,d}+\omega _1)^2}`$ and $`\chi =arctg[\omega _0/(\omega _{f,d}+\omega _1)]`$ with $`\omega _0=2\mu B_0^{f,d}/\mathrm{}`$ and $`\omega _1=2\mu B_z/\mathrm{}`$. On the other hand, we can rewrite Eq.(8) clearly as
$$\gamma _{AB}^{eff}(𝐧(0))=\frac{\eta \omega _{f,d}}{2}_0^{2\pi /\omega _{f,d}}sin\theta cos(\omega _{f,d}t\phi )𝑑t$$
(22)
with $`\theta =arctg(\sqrt{n_1^2+n_2^2}/n_3)`$ and $`\phi =arctg(n_2/n_1)`$. Substituting Eq.(21) into Eqs.(10) and (22), the nonadiabatic noncyclic phases $`\gamma _{AB}^{eff}`$ and $`\gamma _p`$ can be computed, at least numerically.
For simplicity, but without loss of generality, we compute $`\stackrel{~}{T}_+`$ and $`\stackrel{~}{T}_{}`$ in a generalized Kronig-Penny ring consisting of $`N=N_++N_{}`$ uniformly spaced $`\delta `$function barriers with random strengths, where $`N_+`$ $`(N_{})`$ is the number of barriers on the upper (lower) branch. The Hamiltonian for the system in the absence of electromagnetic fields reads $`\widehat{H}_0=(\mathrm{}^2/2m)d^2/dx_\varsigma ^2+_{n_\varsigma =1}^{N_\varsigma }\lambda _{n_\varsigma }\delta (x_\varsigma n_\varsigma a_0)`$, where $`\lambda _{n_\varsigma }`$ is a potential strength parameter, $`a_0`$ is the lattice spacing, and $`\varsigma =+()`$ represents the upper (lower) branch. The spinless electron wave function in the regions where no potentials are present may be written as $`\psi _\varsigma (x_\varsigma )=A_{n_\varsigma }e^{ik_fx_\varsigma }+B_{n_\varsigma }e^{ik_fx_\varsigma }`$. The coefficients A$`_{n_\varsigma }`$ and B$`_{n_\varsigma }`$ across site $`x_\varsigma =n_\varsigma a_0`$ are related through the matrix $`\stackrel{~}{M}_{n_\varsigma }`$,
$$\left(\begin{array}{c}A_{\text{(n+1)}_\varsigma }\\ B_{\text{(n+}1\text{)}_\varsigma }\end{array}\right)=\stackrel{~}{M}_{n_\varsigma }\left(\begin{array}{c}A_{n_\varsigma }\\ B_{n_\varsigma }\end{array}\right)$$
with
$$\stackrel{~}{M}_{n_\varsigma }=\left(\begin{array}{cc}1\frac{iV_{n_\varsigma }}{2k_f}& \frac{iV_{n_\varsigma }}{2k_f}e^{2ik_fn_\varsigma a_0}\\ \frac{iV_{n_\varsigma }}{2k_f}e^{2ik_fn_\varsigma a_0}& 1+\frac{iV_{n_\varsigma }}{2k_f}\end{array}\right),$$
where $`V_{n_\varsigma }=2m\lambda _{n_\varsigma }/\mathrm{}^2`$ is assumed to be distributed uniformly in an interval $`[w/2,w/2]`$. Then, one can find
$$\stackrel{~}{T_\varsigma }=\left(\begin{array}{cc}e^{ik_fN_\varsigma a_0}& 0\\ 0& e^{ik_fN_\varsigma a_0}\end{array}\right)\underset{n_\varsigma =1}{\overset{N_\varsigma }{}}\stackrel{~}{M}_{n_\varsigma }.$$
(23)
Note that $`\stackrel{~}{T_\varsigma }`$ can be further simplified to a $`2\times 2`$ matrix . Substituting the simplified Eq.(23) and $`\mathrm{\Delta }\gamma `$ into Eq.(13), we are able to calculate the transmission coefficient $`T_g`$.
For comparison with the experimental observation, we plot in Fig.1 the calculated ensemble average Fourier spectrum of the conductance for unpolarized electrons , which is defined as
$$|G(\nu )|=|_{B_m}^{B_m}e^{i\nu B_z}G(B_z)𝑑B_z|,$$
(24)
where $`G(B_z)=(e^2/h)_{\pm 𝐧(0)}\overline{T}_g(B_z)`$ with $`\overline{T}_g`$ the average on the initial spin-orientation, $``$ represents the ensemble average. Reasonable parameters in the calculation are determined as follows. As in the experiment, $`B_m=0.35T`$, $`v_f3.0\times 10^5m/s`$, $`g14`$, $`r1.05\mu m`$ (with $`N_\varsigma =4200`$ and $`a_0\frac{\pi }{4}\times 10^9m`$), which leads to the period in a magnetic field $`1.2mT`$; $`w0.267k_f`$, which corresponds to the mean free path $`l=96k_f^2a_0/w^21.0\mu m`$ . The dimensionless coefficient $`\eta 3.5`$, which corresponds to the experimentally reported SO coefficient $`\mathrm{}^2\kappa 5.5\times 10^{10}eVcm`$. Finally, it is typical to consider the case $`ϵ=0.25`$. It is worth emphasizing that the essential feature of Fig.1 is sensitive mainly to the SO coupling parameter $`\eta `$: no clear splitting is present in the main peak if $`\eta `$ is smaller than about $`1.5`$. This implies that the SO-interaction plays a crucial role in the splitting. To clarify the origin of the structure of the main peak, we plot it under both adiabatic and nonadiabatic conditions. The former case is shown in the inset of Fig.1, where $`\gamma _p=\pi cos\theta _i(1cos\chi _a)`$ and $`\gamma _{AB}^{eff}=\eta \pi cos\theta _isin\chi _a`$ with $`\chi _a=arctg(B_0^d/B_z)`$. As $`(B_z/B_0^d)`$ is no longer small, it is very difficult to have a clear analytical understanding of the influences of $`\gamma _p`$ and $`\gamma _{AB}^{eff}`$ on the splitting of the peak. From the inset of Fig.1, we can see that under the adiabatic approximation a somewhat splitting of the main Fourier peak is present only if we include the effective AB phase; however, this splitting feature is obviously not in good agreement with the experimental result. After careful analysis, we understand that the effect of the adiabatic $`\gamma _p`$-phase is too weak to play an important role in causing clearly observable splitting. More remarkably, if ever we take into account both the Pancharatnam phase and the effective AB phase in the nonadiabatic noncyclic case, as shown by the solid line in the main panel of Fig.1, our theoretical result is in excellent agreement with the experimental observation, especially for the splitting and the side structure (two small peaks) of the main peak (see Fig.5 in Ref.). From Fig.1, we can also see that the Pancharatnam phase plays a key role in the main splitting, while the two small side-peaks are closely related to the effective AB phase. It is therefore clarified for the first time that the experimentally observed splitting of the main peak in the ensemble average Fourier spectrum stems from the nonadiabatic noncyclic Pancharatnam phase and the effective AB phase, both being dependent on the SO coupling.
Finally, we remark that the radius of the ring $``$ the mean free path ($`1\mu m`$) and the transport is in the weakly diffusive regime in the experiment. In the present quasi-1D ring, a multi-channel effect, albert weak and secondary, may exhibit in the ensemble average spectrum of the conductance (e.g., broadening and smearing of the peak splitting), which has been ignored in the present work and may deserve for further detailed study in future. Nevertheless, the effect would not affect the present conclusion regarding the splitting qualitatively because the non-adiabaticity of geometric phase is unlikely changed significantly in the present weak diffusive ring.
* To whom correspondence should be addressed. E-mail: zwang@hkucc.hku.hk
Figure Caption
Fig.1 The peak of the ensemble average Fourier spectrum of the conductance in nonadiabatic noncyclic cases. The inset shows the corresponding curves under the adiabatic condition.
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# Color Superconductivity in High Density Quark Matter
## 1 Introduction
In this article I review recent progress on color superconductivity in high density QCD. It is remarkable how much more we know about QCD at high density compared to just two years ago, in 1998. The ultimate fate of baryonic matter at high density is a fundamental property of QCD, with implications for the astrophysics of neutron stars, as well as heavy ion collisions. More generally, we would like someday to understand the QCD phase diagram (figure (1)) in all its complexity as a function of temperature, chemical potential and quark masses.
While lattice studies have been fundamental in determining the behavior of QCD at high temperature, technical difficulties arise in the application of Monte Carlo techniques once a chemical potential is introduced (essentially, the measure of integration is no longer positive definite, due to complex eigenvalues of the Dirac operator). The progress I describe does not rely on brute force techniques, but rather on physical insights associated with the presence of a Fermi surface. The special properties of physics near a Fermi surface are discussed in section 2. Remarkably, rigorous statements can be made in the limit of infinite density.
There are two qualitatively different regimes to be addressed. The regime most likely to be realized in nature (i.e. in the cores of neutron stars or in heavy ion collisions) is that of intermediate density, where the average distance between quarks is still of order a Fermi and the gluons exchanged are rather soft. This regime is strongly coupled, and while we can make qualitative statements about the nature of the ground state using renormalization group techniques, we cannot perform quantitative calculations. The intermediate density case is discussed in section 2.
At extremely high densities the effective coupling $`\alpha _s`$ is small, leading to a weakly coupled liquid of quarks. The only low energy excitations in this liquid are quasi-particles and quasi-holes representing fluctuations near the Fermi surface. There are rather general arguments (see below) which suggest that any attractive interaction, no matter how small, can lead to Cooper pairing in the presence of a Fermi surface. In the case of QCD, this attractive interaction is provided by gluon exchange in the $`\overline{3}`$ channel. Thus, we expect to find condensation of Cooper pairs of quarks in the high density limit.
Actually, even in the high density limit, there are some subtleties which make the problem less than straightforward. Most importantly, although the effective coupling is weak, the exchange of magnetic gluons leads to a long range interaction due to the absence of a magnetic screening mass. Quarks carry electric, rather than magnetic, color charge, and hence are better at screening the timelike component ($`A_0`$) of the gluon field than the spacelike ($`A_i`$). It is only non-perturbative effects which can screen color-magnetic fluctuations, and these are nearly absent at weak coupling. An understanding of dynamic screening due to Landau damping is necessary to control the high density calculations. These issues are discussed in section 4.
The order parameter for color superconductivity is
$$\psi ^TC\mathrm{\Gamma }\psi .$$
(1)
where C is the charge conjugation operator and $`\mathrm{\Gamma }`$ is a matrix in color, flavor and Dirac space. As we will see below, determining the precise form of $`\mathrm{\Gamma }`$ requires some work; it is particularly complicated in the case of three flavors. However, as we will discuss in section 5, in the weak coupling limit the true groundstate can be determined in a controlled approximation. Here we will simply note that single gluon exchange between two quark quasiparticles can be decomposed into an attractive $`\overline{3}`$ channel and a repulsive $`6`$ channel. Thus, at the most naive level we expect an anti-triplet condensate, which breaks $`\mathrm{SU}(3)\mathrm{SU}(2)`$. The main effect of this condensate is that it leads to the Higgs phenomena for at least some subset of the gluons, or equivalently to the Meissner effect and screening of some subset of the color magnetic fields. The phenomenological implications of color superconductivity are still not well understood and will be the subject of investigation for some time to come.
Let me close this introduction with a historical note. The idea that quark matter might be a color superconductor is quite old . The original insight was based on the existence of the attractive $`\overline{3}`$ channel and an analogy with ordinary superconductors. Recent interest in the problem was rekindled by the work of two groups that considered diquark condensation due to instanton-mediated interactions , predicting gaps as large as $`100`$ MeV. These calculations, while uncontrolled, are quite suggestive, and led to the recent progress on the subject. It is often claimed that early investigations predicted tiny gaps, at most of order a few MeV. While this may have been the consensus among the few theorists who had actually worked on the problem, it is actually an unfair characterization of Bailin and Love’s results . A value of the strong coupling large enough to justify the instanton liquid picture of also yields a large gap when substituted in Bailin and Love’s results. After all, instantons are suppressed by an exponential factor $`\mathrm{exp}(2\pi /\alpha _s)`$. Bailin and Love merely suffered from the good taste not to extrapolate their results to large values of $`\alpha _s`$!
## 2 Physics Near a Fermi Surface and the Renormalization Group
Important simplifications arise in the study of cold, dense matter due to the existence of a Fermi surface, which in relativistic systems is likely to be rotationally invariant. The energy of a low-energy excitation (a quasi-particle or -hole) is then independent of the orientation of its momentum $`\stackrel{}{p}`$, and only depends on $`p\mu `$, where $`p=|\stackrel{}{p}|`$ and $`\mu `$ is the chemical potential or Fermi energy. (Here, for simplicity, we will always work with massless quarks.) This leads to a kind of dimensional reduction, so that physics near a Fermi surface is effectively 1+1 dimensional. In particular, arbitrarily weak interactions can lead to non-perturbative phenomena like pair formation.
The renormalization group approach is particularly useful here – we integrate out the modes far from the Fermi surface, leaving only the low-energy quasi-particle and -hole states that are involved in the interesting physics. These excitations might in principle be related to the original quarks in a complicated way, but on quite general grounds must be described by an effective action of the form
$$S_{eff}=𝑑td^3p\psi ^{}\left(i_t(ϵ(p)ϵ_F)\right)\psi +S_{int},$$
(2)
where $`S_{int}`$ contains higher dimensional, local quasi-particle operators. Strictly speaking, this form of the effective action is only valid for models in which the original interactions were local (short ranged). While appropriate for QCD at intermediate densities , where non-perturbative effects are expected to generate screening of magnetic gluons, it must be modified at weak coupling where magnetic fluctuations are long ranged . However, it is the only technique I know of from which we can obtain robust information about the strongly coupled region of the phase diagram. Below I review the results of this analysis, and defer a discussion of the weak coupling phase until the following section.
It can be shown using simple classical scaling arguments that all interactions are irrelevant except for the Cooper pairing interaction (scattering of quasi-particles at opposide sides of the Fermi surface: $`\stackrel{}{p}_1\stackrel{}{p}_2`$) and stricly colinear scattering: $`\stackrel{}{p}_1\stackrel{}{p}_2`$, which can lead to the Overhauser effect (chiral waves) at large-$`N_c`$ . Both of these interactions are classically marginal, so quantum corrections determine their evolution. Here we restrict ourselves to local Cooper pairing operators which are invariant under the full $`SU(3)_L\times SU(3)_R\times U(1)_A`$ chiral symmetry:
$`O_{LL}^0`$ $`=`$ $`(\overline{\psi }_L\gamma _0\psi _L)^2,O_{LR}^0=(\overline{\psi }_L\gamma _0\psi _L)(\overline{\psi }_R\gamma _0\psi _R)`$ (3)
$`O_{LL}^i`$ $`=`$ $`(\overline{\psi }_L\gamma _i\psi _L)^2,O_{LR}^i=(\overline{\psi }_L\stackrel{}{\gamma }\psi _L)(\overline{\psi }_R\stackrel{}{\gamma }\psi _R).`$
These come in both color symmetric ($`\overline{3}`$) and antisymmetric (sextet) combinations. More general operators with different flavor or Dirac structures can be reduced to linear combinations of the basic ones (3), using parity and Fierz rearrangements. This analysis can be extended to operators (such as those induced by instantons) that break the anomalous $`U(1)_A`$ symmetry , yielding a very robust characterization of QCD even at intermediate densities and strong coupling. We will not discuss the details of this more general analysis here, but the results (given reasonable assumptions about the signs and magnitudes of the interactions) are qualitatively similar
The RG evolution of the operators in (3) is determined by quark-quark scattering near the Fermi surface. A bubble graph with four-quark vertices $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_2`$ and external quark lines satisfying the Cooper pairing kinematics yields
$$G_1G_2I(\mathrm{\Gamma }_1)_{i^{}i}(\mathrm{\Gamma }_1)_{k^{}k}\left[(\gamma _0)_{ij}(\gamma _0)_{kl}+\frac{1}{3}(\stackrel{}{\gamma })_{ij}(\stackrel{}{\gamma })_{kl}\right](\mathrm{\Gamma }_2)_{jj^{}}(\mathrm{\Gamma }_2)_{ll^{}}$$
(4)
Here $`I=\frac{i}{8\pi ^2}\mu ^2\mathrm{log}(\mathrm{\Lambda }_{IR}/\mathrm{\Lambda }_{UV})`$, where $`[\mathrm{\Lambda }_{IR},\mathrm{\Lambda }_{UV}]`$ are the upper and lower cutoffs of the momentum shell integrated out. We define the density of states on the Fermi surface to be $`N=\mu ^2/(2\pi ^2)`$ (in weak coupling) and $`t\mathrm{log}(\mathrm{\Lambda }_{IR}/\mathrm{\Lambda }_{UV})`$. The RG flow does not mix $`LL`$ and $`LR`$ operators, nor different color channels. We obtain the following RG equations
$`{\displaystyle \frac{d(G_0^{LL}+G_i^{LL})}{dt}}`$ $`=`$ $`{\displaystyle \frac{N}{3}}(G_0^{LL}+G_i^{LL})^2,`$ (5)
$`{\displaystyle \frac{d(G_0^{LL}3G_i^{LL})}{dt}}`$ $`=`$ $`N(G_0^{LL}3G_i^{LL})^2,`$ (6)
$`{\displaystyle \frac{d(G_0^{LR}+3G_i^{LR})}{dt}}`$ $`=`$ $`0,`$ (7)
$`{\displaystyle \frac{d(G_0^{LR}G_i^{LR})}{dt}}`$ $`=`$ $`{\displaystyle \frac{2N}{3}}(G_0^{LR}G_i^{LR})^2.`$ (8)
The linear combination $`G_{}=G_0^{LL}+G_i^{LL}`$ reaches its Landau pole first, governed by the equation
$$G_{}(t)=\frac{1}{1+(N/3)G_{}(0)t}.$$
(9)
In general, interactions which are attractive at the matching scale, $`G_{}(0)>0`$, will grow during the evolution, and reach a Landau pole at the scale $`t_{}=3/(NG_{}(0))`$. The corresponding energy scale is
$$\mathrm{\Lambda }_{IR}=\mathrm{\Lambda }_{UV}\mathrm{exp}\left(\frac{3}{NG_{}(0)}\right),$$
(10)
which agrees with the usual BCS result. Repulsive interactions, corresponding to negative initial values of the coupling, become weaker near the Fermi surface.
## 3 High Density Limit
In the high density limit the typical momentum transfer between quarks is large, and therefore the effective coupling is small. The properties of this phase can be deduced in a systematic, weak coupling expansion. However, there is one technical problem that must be solved having to do with soft magnetic gluons. In the renormalization of the Cooper pairing interaction there is a region of phase space where the incoming quarks are only slightly deflected by the gluon exchange. This leads to an IR divergence unless arbitrarily soft gluons are screened in some way. While gluons acquire a perturbative electric mass at high density, it can be shown that to all orders in perturbation theory no magnetic mass is generated . The only hope is that Landau damping – a form of dynamic screening affecting the spacelike gluons – is enough to control this IR problem. That this is so was first pointed out by Son , who went on to deduce the following behavior for the diquark gap at weak coupling
$$\mathrm{\Delta }\mu g^5\mathrm{exp}\left(\frac{3\pi ^2}{\sqrt{2}g}\right).$$
(11)
This result has since been confirmed by RG methods as well as Schwinger-Dyson techniques . In this section I discuss some of the details of these calculations, concentrating on the the RG approach.
The magnetic gluon propagator, including vacuum polarization effects from virtual quarks, has the form ($`q_0<<q`$)
$$D_{\mu \nu }^T(q_0,q)=\frac{P_{\mu \nu }^T}{q^2+i\frac{\pi }{2}m_D^2\frac{|q_0|}{q}}.$$
(12)
Strictly speaking $`D_{\mu \nu }^T`$ is gauge dependent, but in our leading order calculations the propagator always appears contracted with gamma matrices next to nearly on-shell external quark lines. Thus the gauge dependent parts are higher order in g due to the equations of motion. The effect of Landau damping is to cut off the small-q divergence in (12) at $`qq_0^{1/3}m_D^{2/3}`$, where $`m_D^2=N_f\frac{g^2\mu ^2}{2\pi ^2}`$ is the Debye screening mass. A common feature of both the Schwinger-Dyson and RG calculations in the weak coupling region is loop integrals dominated by energy transfers of order $`q_0\mathrm{\Delta }`$, and hence momentum transfers of order $`q_{}\mathrm{\Delta }^{1/3}m_D^{2/3}`$. $`q_{}`$ can be made as large as desired by going to high density.
The main technical problem in the RG approach is long range magnetic interactions, or equivalently the presence of soft gluons. At no point can the theory be completely described by quarks with purely local interactions as in (2). The effective Lagrangian contains both quark and gluon excitations (with energies below the cutoff $`\mathrm{\Lambda }`$) and local interactions resulting from integration of higher energy shells. This modifies the form of the RG equations obtained , so that equations (5) and (6) become
$`{\displaystyle \frac{d(G_0^{LL}+G_i^{LL})}{dt}}`$ $`=`$ $`{\displaystyle \frac{N}{3}}(G_0^{LL}+G_i^{LL})^2{\displaystyle \frac{g^2}{9\mu ^2}},`$ (13)
$`{\displaystyle \frac{d(G_0^{LL}3G_i^{LL})}{dt}}`$ $`=`$ $`N(G_0^{LL}3G_i^{LL})^2+{\displaystyle \frac{g^2}{27\mu ^2}}.`$ (14)
The solution of these RG equations leads to a Landau pole in the dominant $`\overline{3}`$, LL and RR channels given by (11). It is worth commenting on the angular momentum of the condensate. The RG equations can be derived for general values of angular momentum l. Naively interpreted, the results suggest that condensates might occur in higher l channels, leading to the breaking of O(3) rotational invariance. A more detailed gap equation analysis shows that this is not the case: a large s-wave gap suppresses the formation of p-wave and higher l gaps. The literature is somewhat confused on this important issue. The papers in address the issue of rotational invariance and do not agree on the size of higher l gaps. However, neither paper addresses the interplay of s-wave and higher l gaps, which is studied in .
## 4 The QCD Groundstate at High Density
In this section I describe the vacuum energy analysis necessary to determine the groundstate of QCD at high density . Neither the RG nor Schwinger-Dyson analyses are sufficient to specify the actual groundstate. Strictly speaking, the former only reveals the energy scale and quantum numbers of the pairing instability, while the latter only identifies extrema of the vacuum energy. As we shall see, there are additional subtleties which can only be resolved by consideration of energetics.
First, let us consider the case of 2 massless flavors. Because the condensate occurs between pairs of either left (LL) or right (RR) handed quarks in the J=L=S=0 channel , and the $`\overline{3}`$ color channel is antisymmetric, the quarks must pair in the isospin singlet (ud - du) flavor channel. However, even in this case there is a subtlety, as the relative color orientations of the LL and RR condensates are not determined by the usual leading order analysis. A misalignment of these condensates violates parity, and further breaks the gauge group beyond $`\mathrm{SU}(3)_\mathrm{c}\mathrm{SU}(2)_\mathrm{c}`$. An analysis of the Meissner effect is necessary to determine the relative orientation , and the effect is higher order in g. There are thus a number of unstable configurations of only slightly higher energy with different color-flavor orientations (and hence different symmetry breaking patterns), leading to the possibility of disorienting the diquark condensate (see figure (3)).
The generalization to three flavors is far from straightforward. Again, one can show that the condensate must occur in the J=L=S=0 and color $`\overline{3}`$ channel. (The sextet condensate is suppressed in the weak coupling limit and I do not discuss it here.) The Pauli principle then requires that the flavor structure again be antisymmetric $`(q_iq_jq_jq_i)`$, for quarks of flavor $`i,j`$. Thus, one can have combinations of condensates which are in the $`\overline{3}`$ of both color and flavor $`\mathrm{SU}(3)_\mathrm{L}`$ or $`\mathrm{SU}(3)_\mathrm{R}`$. Due to the chirality preserving nature of perturbative gluon exchange, there is no mixing of LL and RR condensates, which form independently. One can immediately see that there are a number of possibilities. For example, the condensates for the three flavors and both chiralities might all align in color space, leading to an $`\mathrm{SU}(3)_\mathrm{c}\mathrm{SU}(2)_\mathrm{c}`$ breaking pattern. A more complicated condensate has been proposed called Color Flavor Locking (CFL), in which the $`\overline{3}`$ color orientations are “locked” to the $`\overline{3}`$ flavor orientation. In figure (4) we give a simple picture of CFL condensation.
To determine the nature of the energy surface governing the various color-flavor orientations of the condensate, we can begin by characterizing the color-flavor configuration space of condensates. We consider the ansätz
$$\mathrm{\Delta }_{ij}^{ab}{}_{}{}^{L,R}=A_k^c{}_{}{}^{L,R}ϵ_{}^{abc}ϵ_{ijk},$$
(15)
where a,b are color and i,j flavor indices. L and R denote pairing between pairs of left and right handed quarks, respectively. Under color and flavor A transforms as
$$A^LU_cA^LV^L,$$
(16)
where $`U_c`$ is an element of $`\mathrm{SU}(3)_\mathrm{c}`$ and $`V^L`$ of $`\mathrm{SU}(3)_\mathrm{L}`$. A similar equation holds for $`A^R`$. It is always possible to diagonalize $`A^L`$ by appropriate choice of $`U_c`$ and $`V^L`$:
$$A^L=\left(\begin{array}{ccc}a& 0& 0\\ 0& b& 0\\ 0& 0& c\end{array}\right).$$
(17)
Generically, there does not exist a $`V^R`$ which diagonalizes $`A^R`$ in this basis. In the CFL case, where the diagonalized $`A^L`$ is proportional to the identity, $`a=b=c`$, it is easy to show that one can choose $`V^R`$ such that $`A^R=\pm A^L`$. These two configurations are related by a $`U(1)_A`$ rotation (see section 3). Hence, they are degenerate in the high density limit where gluon exchange dominates. Instanton effects, important at intermediate density, favor $`A^R=A^L`$. Note that parity, if unbroken, requires $`A^L=A^R`$, and hence implies simultaneous diagonalizability.
In we considered the potential vacua parametrized by a,b,c. First, we use the Dyson-Schwinger (gap) equation to determine which of these configurations are energy extrema. Next, we computed the energies of the extrema to determine the true groundstate. A similar analysis has been carried out by Schäfer and Wilczek in the approximation where gluon interactions are replaced by local four fermion interactions. They concluded that the CFL vacuum had the lowest energy. In our analysis, which I summarize below, we included the gluons in the analysis, introducing long range color-magnetic fluctuations (controlled by Landau damping) and Meissner screening into the gap equation and vacuum energy calculations.
At asymptoticaly high densities (weak coupling) the diagrams (a)-(c) in figure (5) give the leading approximation to the effective action. Note that in these diagrams the quark propagators include the diquark condensate (see (20) below), and the gluon propagators include Landau damping, but not the Meissner effect. The latter arises from the condensate-dependence of quark loops in diagrams (c) and (d). The resulting gap equation (figure (6)), with condensate shown explicitly at lowest order in $`\mathrm{\Delta }`$) is given by
$$S^1(q)S_0^1(q)=ig^2\frac{d^4k}{(2\pi )^4}\mathrm{\Gamma }_\mu ^AS(k)\mathrm{\Gamma }_\nu ^BD_{AB}^{\mu \nu }(kq),$$
(18)
where
$$\mathrm{\Gamma }_\mu ^A=\left(\begin{array}{cc}\gamma _\mu T^A& 0\\ 0& C(\gamma _\mu T^A)^TC^1\end{array}\right).$$
(19)
$`D_{AB}^{\mu \nu }`$ is the gluon propagator, including the effects of Landau damping and Debye screening (we assume Feynman gauge throughout).
We will restrict the color group structure in the gap equation to the attractive anti-symmetric $`\overline{3}`$ channel, which projects out the anti-symmetric part of $`S(k)`$ in color space in the gap equation. Here $`S`$ is the fermion propagator for the spinor $`(\psi _a^i,\psi _a^{iC})`$ with $`i`$ a flavor index and $`a`$ a color index.
For the three flavor case $`S`$ can be written explicitly as an $`18\times 18`$ matrix in color flavor space. The inverse propagator may be written
$$S^1(q)=\left(\begin{array}{cc}q/+\mu /& \gamma _0\mathrm{\Delta }^{}\gamma _0\\ \mathrm{\Delta }& q/\mu /\end{array}\right)$$
(20)
where $`\mu /=\mu \gamma _0`$. $`\mathrm{\Delta }`$ is a $`9\times 9`$ matrix which for the ansätz (17) takes the form
$$\mathrm{\Delta }=\left(\begin{array}{ccccccccc}0& 0& 0& 0& c& 0& 0& 0& b\\ 0& 0& 0& c& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& b& 0& 0\\ 0& c& 0& 0& 0& 0& 0& 0& 0\\ c& 0& 0& 0& 0& 0& 0& 0& a\\ 0& 0& 0& 0& 0& 0& 0& a& 0\\ 0& 0& b& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& a& 0& 0& 0\\ b& 0& 0& 0& a& 0& 0& 0& 0\end{array}\right)$$
(21)
Because we are dealing with a diquark condensate the non-trivial part of the gap equation involves the lower left $`9\times 9`$ block. We will refer to this sub-block of the propagator $`S`$ as $`S_{21}`$.
For a particular ansätz $`\mathrm{\Delta }`$ to be a solution to the gap equation we require that the color antisymmetric part of $`T^AS_{21}(k)T^A`$ (corresponding to the $`\overline{3}`$ channel) be proportional in color-flavor space to the off-diagonal 21 submatrix of $`S^1(q)S_0^1(q)`$, or $`\mathrm{\Delta }(q)`$, which appears on the LHS of the gap equation (see for more discussion of this point).
The propagator may be found by inverting the sparse matrix in (20) using Mathematica. Only three ansätze satisfy our condition: $`a=b=c`$; $`a=b,c=0`$; $`b=c=0`$. We refer to these solutions as (111) (color-flavor locking), (110) ($`30`$ breaking) and (100) ($`32`$ breaking) respectively. For these ansätze the color antisymmetric part of $`T^AS_{21}(k)T^A`$ has the form of a constant multiplying the matrix form (21) with $`a,b,c`$ set to 0 or 1 as is appropriate for the ansätz. After contour integration over l, we find the following gap kernels
$`(111)`$ $`:`$ $`{\displaystyle \frac{2}{3}}{\displaystyle \frac{\mathrm{\Delta }}{\sqrt{k_0^2+\mathrm{\Delta }^2}}}+{\displaystyle \frac{1}{3}}{\displaystyle \frac{\mathrm{\Delta }}{\sqrt{k_0^2+4\mathrm{\Delta }^2}}}`$
$`(110)`$ $`:`$ $`{\displaystyle \frac{\mathrm{\Delta }}{2\sqrt{k_0^2+\mathrm{\Delta }^2}}}+{\displaystyle \frac{\mathrm{\Delta }}{2\sqrt{k_0^2+2\mathrm{\Delta }^2}}}`$
$`(100)`$ $`:`$ $`{\displaystyle \frac{\mathrm{\Delta }}{\sqrt{k_0^2+\mathrm{\Delta }^2}}}`$ (22)
These kernels are to be substituted in the following gap equation, which we obtain under the approximation $`q_0<<|\stackrel{}{q}|`$. (We also neglected the anti-particle contributions, which are suppressed by powers of $`1/\mu `$.)
$`\mathrm{\Delta }(p_0)`$ $`=`$ $`{\displaystyle \frac{g^2}{12\pi ^2}}{\displaystyle }dq_0{\displaystyle }d\mathrm{cos}\theta ({\displaystyle \frac{\frac{3}{2}\frac{1}{2}\mathrm{cos}\theta }{1\mathrm{cos}\theta +(G+(p_0q_0)^2)/(2\mu ^2)}}`$
$`+{\displaystyle \frac{\frac{1}{2}+\frac{1}{2}\mathrm{cos}\theta }{1\mathrm{cos}\theta +(F+(p_0q_0)^2)/(2\mu ^2)}})K(q_0),`$
where F and G represent the medium effects on the electric and magnetic gluons, and $`K(q_0)`$ is one of the gap kernels from (22). The Meissner effect makes an additional contribution to G beyond that of Landau damping. In we evaluated the gluon vacuum polarization $`P_{\mu \nu }(q_0,q)`$ in the presence of a diquark condensate. (A more detailed computation of the Meissner effect is given by Rischke , with similar results.) The additional Meissner screening is given by $`\delta G\frac{1}{2}𝒫_{ij}^TP_{ij},`$ where $`𝒫_{ij}^T=\left(\delta _{ij}\widehat{q}_i\widehat{q}_j\right)`$ is the transverse projection operator. At low momenta, $`q_0,q\mathrm{\Delta }`$, $`\delta G(q_0,q)`$ is of order the Debye mass $`m_Dg\mu `$, while at larger energy or momenta the effect is suppressed by a power of $`\frac{\mathrm{\Delta }}{q_0}`$ or $`\frac{\mathrm{\Delta }}{q}`$. We limited ourselves to an estimate of the size of the Meissner effect on the gap solutions. To this end, we used the following approximation for $`\delta G`$:
$$\delta G(q_0,q)m_D^2\frac{\mathrm{\Delta }_0}{\sqrt{q^2+q_0^2+\mathrm{\Delta }_0^2}},$$
(24)
where $`\mathrm{\Delta }_0`$ is the maximum value of the function $`\mathrm{\Delta }(k_0,k)`$. Note we did not introduce any color structure in $`\delta G`$; all gluons experience the same magnetic screening. While this is a crude approximation, it gives the rough size of the Meissner effect on $`\mathrm{\Delta }`$.
We solved the gap equations for all three gap kernels using this form of the Meissner effect, and the results are shown in figure (7) for the case of $`\mu =400`$ MeV. The effect is to decrease the size of the condensate but it is a small perturbation on the solutions obtained without the Meissner effect.
To determine which of the above gaps is the true minimum energy state we must calculate the vacuum energy, which receives contributions from vacuum to vacuum loops of both quarks and gluons (figure 1). We start with the CJT effective potential , which upon extremization wrt appropriate propagators and vertices leads to the Schwinger-Dyson equations. The fermion equation is the gap equation given above, while the gluon equation reproduces Landau damping. We wish to compare energies corresponding to our three solutions to determine which one is the true vacuum (the difference in energies $`V`$ will be gauge invariant, whereas actual values are not). It is easy to show that the value of the effective potential evaluated on the gap solution is given by:
$$V=i\frac{d^4p}{(2\pi )^4}\mathrm{trln}S(p)/S_0(p).$$
(25)
Diagramatically, this is equivalent to the graph of figure 1(a) when evaluated on the gap solution.
The fermion loops are most easily calculated by going to a basis where $`S_0S^1`$ is diagonal in color-flavor space. Note that the gap matrix $`\mathrm{\Delta }`$ has non-trivial Dirac structure that must be accounted for : $`\mathrm{\Delta }=\mathrm{\Delta }_1\gamma _5P_++\mathrm{\Delta }_2\gamma _5P_{}`$, where $`P_\pm `$ are particle and anti-particle projectors. Our analysis has been restricted to the particle gap function $`\mathrm{\Delta }_1`$. The anti-particle gap function $`\mathrm{\Delta }_2`$ has its support near $`k_02\mu `$, and its contribution to the vacuum energy is suppressed. There are 18 eigenvalues, which occur in 9 pairs. The product of each pair is of the form
$$\left(1+a\frac{\mathrm{\Delta }^2(k_o,k)}{k_0^2+(|\stackrel{}{k}|\mu )^2}\right),$$
(26)
where a is an integer. For our three cases we obtain the following sets of eigenvalues:
$`(111)`$ $``$ $`8\times \{a=1\},1\times \{a=4\}`$
$`(110)`$ $``$ $`4\times \{a=1\},2\times \{a=2\}`$
$`(100)`$ $``$ $`4\times \{a=1\}`$ (27)
The binding energy is of order
$`E_q`$ $``$ $`{\displaystyle d^3k𝑑k_0\mathrm{ln}\left[1+a\frac{\mathrm{\Delta }^2(k_0,k)}{k_0^2+(k\mu )^2}\right]}`$ (28)
$``$ $`a\mu ^2\mathrm{\Delta }_0^2,`$ (29)
where $`\mathrm{\Delta }_0`$ is the maximum value of the gap function $`\mathrm{\Delta }(k_0,k)`$, which has rather broad support in both energy and momentum space away from the Fermi surface. A more precise answer than (28) requires numerical evaluation, but it is clear that the result scales with a and has only a weak (logarithmic) dependence on the variations in the shape of $`\mathrm{\Delta }(k_0,k)`$. Substituting our numerical results for the gaps in the three cases, it is easy to establish that
$$E(111)<E(110)<E(100).$$
(30)
We find that the CFL vacua remains the lowest energy state, at least at asymptotically high densities where the calculation is reliable. The Meissner effect is a small correction to the vacuum energy at asymptotic densities. Configurations which satisfy the gap equations but are not the global minimum of energy are presumably saddlepoints, since they are continuously connected to the CFL vacuum via color and flavor rotations.
## 5 Conclusions
In this contribution I have tried to summarize some important progress of the last two years on the theory of cold, dense quark matter. Due to space limitations, I was not able to discuss a number of important issues, such as the low energy effective Lagrangian, continuity of hadronic and quark phases and more phenomenological studies. I list some of the important papers in .
I thank my collaborators N. Evans, J. Hormuzdiar and M. Schwetz and my colleagues D. Hong, R. Pisarski, K. Rajagopal, M. Rho, D. Rischke and T. Schäfer for contributing to my understanding of this subject. My research was supported under DOE contract DE-FG06-85ER40224 and by a JSPS visiting fellowship.
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# The massive stellar content in NGC 604 and its evolutionary state.
## 1 Introduction
Giant extragalactic H ii regions (RH iis) are amongst the brightest objects in galaxies. RH iis have been studied extensively because they are the best indicators of the conditions that lead to massive star formation, and they show the cloud properties immediately after the moment when stars form. Equally important, they are favourable places where to measure the chemical abundances and star formation rates in galaxies, and hence can be used to trace the chemical and star formation history of the universe, when observed in galaxies at different redshift (Madau et al. 1996; Kobulnicky & Zaritsky 1999).
RH iis are characterized to have a size larger than 100 pc and H$`\alpha `$ luminosity brighter than 10<sup>39</sup> erg s<sup>-1</sup> (Kennicutt 1984). Therefore, the nebula requires an ionizing photon luminosity larger than 10<sup>51</sup> s<sup>-1</sup>; this is provided by a stellar cluster that contains more than 100 young massive stars. These characteristics are very similar to those of starburst galaxies. However, they are less luminous than prototypical starbursts, and thus they are referred to as mini-starbursts (Walborn 1991). Like starbursts, RH iis show a nebular emission-line spectrum at optical wavelengths, and an absorption-line spectrum at wavelengths shorter than the Balmer jump (Leitherer 1997; Rosa, Joubert & Benvenuti 1984). This spectral morphology reflects the fact that RH iis are powered by massive stars. These stars emit photons with energies of tens of eV which are absorbed and re-emitted in their stellar winds, producing ultraviolet resonance transitions. However, the stellar wind is optically thin to most of the ultraviolet photons, that can travel tens of parsec from the star before they are absorbed and photoionize the surrounding interstellar medium. Subsequently, this ionized gas cools down via an emission spectrum. This spectral dichotomy picture allows to derive the stellar content and the evolutionary state of the cluster, through the analysis of the ultraviolet (e.g. Vacca et al. 1995; Leitherer et al. 1996), or the optical light (e.g. Cid Fernandes et al. 1992; García-Vargas & Díaz 1994; Stasińska & Leitherer 1996) using evolutionary synthesis and photoionization models. The two techniques have been applied to the prototypical starburst nucleus NGC 7714, giving similar results (González Delgado et al. 1999a).
However, around the Balmer jump the spectra of starbursts (González Delgado et al. 1998) and some H ii regions (Terlevich et al. 1996) show the higher order terms of the Balmer series and He i lines in absorption, formed in the photospheres of massive stars. These stellar lines can be detected in absorption because the strength of the gaseous Balmer lines in emission decreases rapidly with decreasing wavelength, whereas the equivalent width of the stellar absorption lines is almost constant with wavelength (González Delgado, Leitherer & Heckman 1999b). Evolutionary synthesis models that predict the profiles of the higher order terms of the Balmer series and He i lines in absorption can also be used to estimate the evolutionary state of the stellar cluster in H ii regions.
The second more luminous and brightest H ii region in the Local Group of galaxies, after 30 Dor in the LMC, is NGC 604 in M 33. Its distance, 840 kpc (Freedman et al. 1991), allows detailed studies of the individual stars, and also to obtain integrated properties. Thus, NGC 604 is an excellent laboratory in which to explore questions about the effect of star formation, the slope and upper mass limit of the IMF, and to test the consistency between the three techniques described above, that allow to determine the stellar content and the evolutionary state of the ionizing stellar cluster from the spatially integrated ultraviolet and optical spectra of the H ii region.
NGC 604 has been intensively studied in the past (e.g. Peimbert 1970; Israel & van der Kruit 1974; Smith 1975; Hawley & Grandi 1977; Kwitter & Aller 1981; Conti & Massey 1981; Rosa & D’Odorico 1982; Viallefond & Goss 1986). A precise determination of the chemical composition of the gas was made by Díaz et al. (1987); they give an oxygen abundance 12+log(O/H)=8.5. VLA observations indicate that the gas is very tenuous (average electron density, $`rms`$ $`N_\mathrm{e}`$, between a few and $`100`$ cm<sup>-3</sup>), and not very dusty, with a mean visual extinction $`A_V0.5`$ mag (Churchwell & Goss 1999). The morphology of the ionized gas is very complex, showing many filaments and shell structures that are expanding (e.g. Hippelein & Fried 1984; Rosa & Solf 1984; Clayton 1988; Sabalisck et al. 1995; Muñoz-Tuñón et al. 1995; Yang et al. 1995; Medina Tanco et al. 1997). This morphology is a consequence of the violent star formation activity in NGC 604. HST imaging photometry has revealed that the stellar cluster is resolved into $``$200 massive stars in an area of $``$10000 pc<sup>2</sup> (Drissen, Moffat & Shara 1993; Hunter et al. 1996). Evidences that the stellar cluster is evolved come from the detection of Wolf-Rayet stars (Conti & Massey 1981; D’Odorico & Rosa 1981; Drissen, Moffat & Shara 1990, 1993), one supernova remnant (D’Odorico et al. 1980), and one candidate to red supergiant (Terlevich et al. 1996). The existence of these stars as members of the stellar cluster of NGC 604 suggests that the age of the region is in the range 3-5 Myr.
This paper presents spatially integrated ultraviolet and optical spectra of NGC 604. The goal is to derive the massive stellar content and the evolutionary state of the ionizing cluster, by means of evolutionary synthesis and photoionization models applied to the ultraviolet resonance wind stellar lines, to the nebular emission lines and to the higher-order terms of the Balmer series and He i lines in absorption. The consistency between the results obtained with the three techniques, and the similarity with the results obtained from studies based on the detection of individual stars, strengthens the reliability and power of these techniques when they are applied to determine the stellar content of more distant star-forming regions. Section 2 presents the observations. Section 3 describes the photoionization models that fit the emission-line spectrum of the nebula. In section 4, the higher-order terms of the Balmer and He i lines in absorption are analyzed. Evolutionary synthesis models of the ultraviolet resonance wind stellar lines are in section 5. In section 6, the massive stellar content is derived. The summary and conclusions are in section 7.
## 2 Observations
### 2.1 HST images
NGC 604 has been intensively observed with the WFPC2 camera onboard the HST. For the purpose of this work, we have retrieved HST archive images at the ultraviolet (F170W filter) and optical H$`\alpha `$ (filter F656N) wavelengths. Two F170W exposures of 350 s each and two F656N exposures of 1000 s each were combined to produce one final ultraviolet (Fig. 1) and one H$`\alpha `$ image (Fig. 2). The nebula was centred in the PC camera, therefore the ionizing cluster is sampled at 0.046 arcsec pixel<sup>-1</sup>, corresponding to 0.19 pc pixel<sup>-1</sup> (at a distance of 840 kpc, 1 arcsec corresponds to 4.1 pc). The ultraviolet emission is spread in the inner 20$`\times `$20 arcsec. The stellar density looks much lower than that observed in super star clusters. Surface brightness photometry indicates that 75 percent of the ultraviolet flux is produced by the core of the stellar cluster, that is contained within a central nebular hole of diameter $`15`$ arcsec; we shall refer to this as the centre of the nebula. The H$`\alpha `$ image shows many filaments and shell structures that extend out to $`200`$ pc from the centre of the nebula. However, the brightest parts of the nebula trace almost a ring structure at a distance of 40 pc from the centre.
The total ultraviolet flux per unit wavelength in the F170W image is $`1.9\times 10^{13}`$ erg s<sup>-1</sup> cm<sup>-2</sup> Å<sup>-1</sup>. The H$`\alpha `$ flux is $`4.0\times 10^{11}`$ erg s<sup>-1</sup> cm<sup>-2</sup>, in good agreement with the value measured by Churchwell & Goss (1999).
We perform surface brightness photometry in the H$`\alpha `$ image using a circular aperture around the central cluster nebular hole. The cumulative flux is plotted in Fig. 3. This plot shows that 50, 75 and 90 percent of the total flux is contained in the inner 60, 100 and 140 pc radii<sup>1</sup><sup>1</sup>1Unless otherwise specified, distances refer to radial distances from the adopted centre of symmetry, as shown in Figs. 1 and 2., respectively. The H$`\alpha `$ flux distribution of NGC 604 has the characteristic features of a core-halo distribution, where the core can be quantitative and qualitatively explained by a thick shell structure. A detailed modelling of the geometry of NGC 604 will be presented elsewhere. The surface brightness weighted core radius of the region (following the method described in García-Vargas et al. 1997) is 80 pc.
The $`rms`$ electron density is also estimated from the H$`\alpha `$ image as $`(Q/\alpha _BV)^{1/2}`$, where $`Q`$ is the ionizing photon luminosity, $`\alpha _B`$ the total recombination coefficient of H, and $`V`$ the volume of the ionized region. Fig. 3 also shows the $`rmsN_\mathrm{e}`$ as a function of distance, that ranges from values between less than 100 cm<sup>-3</sup> to a few cm<sup>-3</sup>, with $`rmsN_\mathrm{e}10`$ cm<sup>-3</sup> at 100 pc from the center of the nebula.
### 2.2 Ultraviolet spectroscopy
NGC 604 has been observed with the IUE through the large aperture (9.5$`\times `$22 arcsec)<sup>2</sup><sup>2</sup>2The precise form and size of the IUE large aperture is not known; we take a size of $`22\times 9.5`$ arcsec. The actual form is somewhat intermediate between the rectangle shown here and an enclosed ellipse. in low-dispersion mode, with either the SWP (1100–1900 Å) or the LWR (1900–3200 Å) cameras. The spectra were taken with the large aperture located at several different positions across the region. We have retrieved from the IUE Newly Extracted Spectra (INES, Rodríguez-Pascual et al. 1999; Cassatella et al. 2000) archive the spectra taken at the positions shown in Fig. 1 (see also Table 1). Several of the spectra were also retrieved in February 1997 from the IUE ULDA archive. When comparing the two different extractions of the same aperture spectrum, we noticed that the slope of the spectra were slightly steeper in the ULDA archive; both extractions have the same flux at 1500 Å but the ULDA extraction is 10 percent brighter at 1280 Å than the INES extraction. ULDA used a different extraction algorithm (IUESIPS) than INES, and their discrepancies are well documented (González-Riestra et al. 2000). INES represents the state of the art in the extraction of IUE spectra; so this is the data set that we use.
The flux at 1750 Å measured in the spectra SWP5688 (PA=149$`\mathrm{°}`$) and SWP6638 (PA=110$`\mathrm{°}`$) is $`1.6\times 10^{13}`$ erg s<sup>-1</sup> cm<sup>-2</sup> Å<sup>-1</sup> and $`1.3\times 10^{13}`$ erg s<sup>-1</sup> cm<sup>-2</sup> Å<sup>-1</sup>, respectively. The sum total flux in these two apertures represents approximately the total ultraviolet flux emitted by the stellar cluster in NGC 604, and it amounts to $`2.9\times 10^{13}`$ erg s<sup>-1</sup> cm<sup>-2</sup> Å<sup>-1</sup>. This value is 45 percent larger than the total monochromatic flux measured in the F170W HST image ($`\mathrm{F}_{170}=1.9\times 10^{13}\mathrm{erg}\mathrm{s}^1\mathrm{cm}^2\mathrm{\AA }^1`$). Fig. 4 shows the IUE spectra at PA=110$`\mathrm{°}`$ and at PA=32$`\mathrm{°}`$, together with the average spectrum of those at PA=149$`\mathrm{°}`$, 136$`\mathrm{°}`$ and 159$`\mathrm{°}`$.
### 2.3 Optical spectroscopy
Optical spectra were obtained with the 4.2m William Herschel Telescope at the Observatorio del Roque de los Muchachos (La Palma), as part of the GEFE<sup>3</sup><sup>3</sup>3Grupo de Estudios de Formación Estelar, is an international collaboration formed to take advantage of the international time granted by the Comité Científico Internacional at the Observatories in the Canary Islands. collaboration. The details of the observations and data reduction are in Terlevich et al. (1996). For the purpose of this work, we use the scanned spectra at PA=60$`\mathrm{°}`$ and at PA=120$`\mathrm{°}`$, and the single long-slit spectrum at PA=131$`\mathrm{°}`$. The scanned spectra were centred at R.A.=1<sup>h</sup> 31<sup>m</sup> 43<sup>s</sup> and Dec=30$`\mathrm{°}`$ 31′ 52″, and they cover the core of the region by displacing a 1 arcsec wide longslit in steps of 1 arcsec and taking at each position a 1 minute exposure. The process is repeated until an area of 1.75′$`\times `$0.18′ (for the scan at PA=120$`\mathrm{°}`$) and 1.75′$`\times `$0.23′ (for the scan at PA=60$`\mathrm{°}`$) are covered (see Fig. 2). A one dimensional spectrum representative of the inner region of the nebula was obtained by merging the extractions that cover 11″$`\times `$60″ (at PA=120$`\mathrm{°}`$) and 14″$`\times `$60″ (at PA=60$`\mathrm{°}`$). At PA=131$`\mathrm{°}`$, a one-dimensional spectrum was extracted of length 4.6 arcsec, and centered at the position where the continuum is maximum (see fig. 3b in Terlevich et al. 1996). At this position the fraction of the stellar continuum with respect to the nebular emission lines is maximum, because it corresponds to the zone of the nebular hole. As previously noted by Terlevich et al. (1996), the spectrum shows the higher-order terms of the Balmer series and He i lines in absorption, and the continuum represents well the stellar light of the ionizing stellar cluster in NGC 604 (Fig. 5).
## 3 Modelling the nebular emission-lines
The emission-line spectrum of an RH ii depends on the radiation field from the ionizing stellar cluster, and on the density distribution and chemical composition of the gas. A photoionization code takes as input the spectral energy distribution of the cluster, and it solves the ionization-recombination and heating-cooling balances, to predict the ionization structure of the nebula, the electron density and the intensity of the emission lines. By comparing these output with observations, it is possible to obtain information about the ionizing stars and their evolutionary state.
The star formation law, age and massive stellar content of the stellar clusters can be constrained by comparing the observed emission line strengths with the predictions from the photoionization models, when the code uses as input the spectral energy distribution (SED) generated by a stellar evolutionary synthesis code. This technique has been used successfully to study the stellar content in starbursts and RH iis (e.g. García-Vargas, Bressan & Díaz 1995a,b; García-Vargas et al. 1997; González Delgado et al. 1999a; Luridiana, Peimbert & Leitherer 1999; Stasińska & Schaerer 1999). Here, we fit the emission-line spectrum of NGC 604 using the photoionization code CLOUDY (version 90.04, Ferland 1997).
### 3.1 Input parameters
In order to predict the intensities of the emission lines, we fix the geometry of the nebula, the electron density and the chemical composition of the gas. We assume that the constant density gas is ionization bounded and spherically distributed around the ionizing cluster. The inner radius is 20 pc (as seen in the H$`\alpha `$ image), and the outer radius is determined by the ionization front. The results do not depend strongly on the inner radius if this changes by a factor 2 (e.g. Luridiana, Peimbert & Leitherer 1999). We assume that the gas occupies only a fraction of the sphere. Thus the filling factor, $`\varphi `$, is a free parameter with values of 10<sup>-1</sup>, $`5\times 10^2`$, 10<sup>-2</sup>, 10<sup>-3</sup> and 10<sup>-4</sup>. The change of filling factor is equivalent to changing the ionization parameter $`U`$, defined as $`Q/(4\pi RN_\mathrm{e}c`$); where $`Q`$ is the ionizing photon luminosity, $`N_\mathrm{e}`$ the electron density, $`c`$ the speed of light and $`R`$ the distance of the gas to the ionizing source. For a spherical geometry, the average $`U`$ is proportional to $`(\varphi ^2N_\mathrm{e}Q)^{1/3}`$.
The chemical abundances are fixed to those derived by Díaz et al. (1987). We scale the chemical composition of the gas to the value of the oxygen abundance, 12+log(O/H)=8.5, except the abundance of N, Ne, S and He, for which we take the values derived by Díaz et al. (1987; see Table 2). The electron density derived by Díaz et al. (1987) is of the order of 100 cm<sup>-3</sup>. Our measurements of the \[S ii\]6717/6731 ratio indicate the low-density limit, with values below 100 cm<sup>-3</sup> <sup>4</sup><sup>4</sup>4Notice, however, that the atomic parameters for the sulfur have changed in the intervening time. We have computed constant density models at two values of the electron density, 100 cm<sup>-3</sup> and 30 cm<sup>-3</sup>. The latter value comes from the filling factor and an $`rms`$ electron density of 10 cm<sup>-3</sup>, derived from the H$`\alpha `$ image. Also, an electron density of 30 cm<sup>-3</sup> predicts a \[S ii\]6717/6731 ratio which is in better agreement with the observed value.
### 3.2 The radiation field
The radiation field used as input to CLOUDY is the spectral energy distribution generated by the evolutionary synthesis code developed by Leitherer and colaborators (Leitherer et al. 1999). The code includes the new set of stellar evolutionary models of the Geneva group (Schaller et al. 1992; Schaerer et al. 1993a,b; Charbonnel et al. 1993; Meynet et al. 1994), and the stellar atmospheres grid compiled by Lejeune et al. (1997), supplemented by the expanding spherical non-LTE models of Schmutz, Leitherer & Gruenwald (1992). The latter stellar atmosphere models are applied to stars with very strong stellar winds. The spectral energy distribution was generated using the Z=0.008 (Z$``$/2.5) metallicity tracks, assuming that stars and gas have the same metallicity. We assume two different star formation scenarios: instantaneous burst and continuous star formation at a constant rate. Different models are also computed with different assumptions about the IMF. The slope is Salpeter ($`\alpha `$= 2.35), flatter ($`\alpha `$= 1.5), or steeper ($`\alpha `$= 3.3). The upper mass limit cut-off is set to 30, 60, 80 or 120 M$``$. The lower mass limit is fixed to 1 M$``$. However, photoionization models are not influenced by the lower limit mass cut-off if it is below 10 M$``$. The spectral energy distribution is normalized to an ionizing photon luminosity of $`\mathrm{log}\mathrm{Q}=51.54`$ s<sup>-1</sup>. This value is derived from the total H$`\alpha `$ flux measured in the image, after correcting by the mean extinction, $`A_V=0.5`$ mag, derived by Churchwell & Goss (1999). We will check below the compatibility of this normalization with the mass of the stellar cluster derived from the ultraviolet continuum luminosity.
### 3.3 The observational constraints
The emission lines were measured in our scanned spectrum of the nebula. The observed ratios have been corrected by the mean extinction, c(H$`\beta `$)=0.22, derived by Churchwell & Goss (1999). This value is in agreement with the reddening derived from the H$`\gamma `$/H$`\beta `$ ratio measured in the scanned spectrum, when this ratio is corrected by the underlying stellar absorption. This extinction is lower than the values derived by Díaz et al. (1987), c(H$`\beta `$)=0.3-0.4; however, the extinction changes across NGC 604 (Maíz-Apellániz 1999). The de-reddened emission line ratios to be fitted by the photoionization models are included in Table 3. This table also includes the emission line ratios derived by Díaz et al. (1987) in several parts of the nebula, and those from an integrated spectrum in Maíz-Apellániz (1999). The range of values for each of these emission line ratios defines the tolerance that we accept for the difference between the observed and the predicted values from the models.
Following Stasińska & Schaerer (1999), we try to fit the strength of the emission lines, and also emission-line ratios indicative of the electron temperature (\[O iii\]4363/5007), electron density (\[S ii\]6717/6732) and ionization structure (\[O iii\]5007/\[O ii\]3727, \[S iii\]9069/\[S ii\]6716+6732). Another important observational constraint is the radius of the nebula. The H$`\alpha `$ image shows that the flux extends out to 200 pc of the center of the region; however, more than 90, 75 and 50 percent of the total flux is within the inner 140 pc, 100 pc and 60 pc radii, respectively (Fig. 3). Our determination of the core radius of the region is 80 pc.
### 3.4 Model results
We start by fixing the filling factor and the electron density by means of fitting the line ratios \[S ii\]6717+6731/H$`\beta `$ and \[S ii\]6717/6731, and the radius of the nebula. \[S ii\]6717+6731/H$`\beta `$ is a good calibrator of the ionization parameter, because this ratio does not depend much on the IMF assumptions, on the star formation law or on the evolutionary state of the cluster (González Delgado et al. 1999a). Instead, this ratio depends on the filling factor, and for fixed geometry, ionizing photon luminosity and electron density, it depends on the ionization parameter<sup>5</sup><sup>5</sup>5Nonetheless, we should always be cautious when using sulfur dependent diagnostics, given the uncertainties about the atomic parameters and its not so well known ionization; also this calibration relies on the assumption of an ionization bounded nebula.. The observed ratio, \[S ii\]6717+6731/H$`\beta 0.4`$, indicates a filling factor $`\varphi 0.1`$ (Fig. 6). Therefore, a large fraction of the volume of the region is filled with ionized gas; however, this gas is very tenuous because the $`rms`$ electron density measured in the H$`\alpha `$ image is $`10`$ cm<sup>-3</sup>, and the electron density in the S<sup>+</sup> zone, derived from the definition of filling factor ($`\varphi =(N_\mathrm{e}(rms)/Ne)^2`$), is $`N_\mathrm{e}30`$ cm<sup>-3</sup>. The observed ratio \[S ii\]6717/6731=1.4, indicates that the electron density is $`100`$ cm<sup>-3</sup>. Models with a filling factor of 0.1 and $`N_\mathrm{e}=30`$ cm<sup>-3</sup> predict well the observed \[S ii\]6717/6731 ratio. If the electron density were as high as $`N_\mathrm{e}=100`$ cm<sup>-3</sup>, then the filling factor would have to be lower, $`\varphi `$=0.01, and in this case the values of \[S ii\]6717+6731/H$`\beta `$ and \[S ii\]6717/6731 predicted by the models are 0.9 and 1.3, respectively. These predictions are larger and smaller than the observed values, respectively. Furthermore, models with $`\varphi =0.1`$ and $`N_\mathrm{e}=30`$ cm<sup>-3</sup> predict a Strömgren radius $``$110 pc, in agreement with the radius derived from the H$`\alpha `$ surface brightness photometry. Models with $`N_\mathrm{e}=100`$ cm<sup>-3</sup> predict a radius a factor two smaller than models with $`N_\mathrm{e}=30`$ cm<sup>-3</sup>. Therefore, we conclude that a filling factor $`\varphi 0.1`$ and a density $`N_\mathrm{e}=30`$ cm<sup>-3</sup> fit well the ratios \[S ii\]6717+6731/H$`\beta `$ and \[S ii\]6717/6731, and the size of the region.
The emission-line ratios \[O iii\]5007/H$`\beta `$, \[O ii\]3727/H$`\beta `$, \[O i\]6300/H$`\beta `$, and \[N ii\]6584/H$`\beta `$ have been plotted as a function of the cluster age, using a filling factor $`\varphi `$= 0.1 and an electron density $`N_\mathrm{e}=30`$ cm<sup>-3</sup>. To check whether the models also fit line ratios indicative of the structure of the region, we have plotted the $`\eta `$ parameter, defined as (\[O iii\]5007+4959/\[O ii\]3727)/(\[S iii\]9069+9532\]/\[S ii\]6717+6731), and the electron temperature predicted by the models as a function of the age. The parameter $`\eta `$ is a measure of the softness of the radiation field (Vílchez & Pagel 1988). The models assume two different star formation laws, an instantaneous burst or continuous star formation at a constant rate, and different assumptions about the slope ($`\alpha `$=2.35, 1.50, 3.0) and upper mass limit cut-off of the IMF ($`M_{\mathrm{up}}=`$120, 80, 60, 30 M$``$).
Fig. 7 compares the observed emission line ratios with the prediction of continuous and instantaneous burst, for a Salperter IMF and upper mass limit cut-off of 80 M$``$. The observed ratios are in better agreement with burst than with continuous star formation models (that predict high excitation lines and \[O i\]6300/H$`\beta `$ larger than observed). The continuous star formation scenario can be made compatible with the observed emission-line ratios if stars more massive than 50 M$``$ are not formed in the cluster; however, these models predict significantly fewer numbers of Wolf-Rayet stars in the cluster. Thus, continuous star formation models cannot fit the emission-line ratios and predict the existence of Wolf-Rayet stars simultaneously.
Fig. 8 compares the observed emission line ratios with the prediction of an instantaneous burst, for different assumptions of the IMF upper mass limit cut-off. Models with $`M_{\mathrm{up}}=30`$ M$``$ predict high excitation lines (low excitation lines) ratios which are much lower (higher) than the observed values. Thus, similar results are obtained if $`M_{\mathrm{up}}50`$ M$``$. Burst models with $`M_{\mathrm{up}}60`$ M$``$ fit the emission-line ratios if the age of the cluster is 2.5–3 Myr or 4.5–4.8 Myr. However, if $`M_{\mathrm{up}}=60`$ M$``$, Wolf-Rayet stars only appear in the cluster if it is 4–5 Myr old. Models with $`M_{\mathrm{up}}=120`$ M$``$ and 80 M$``$ give similar results after the first 2 Myr, and the same behaviour with age. The electron temperature, $`\eta `$ and the emission-line ratios, all indicate that massive stars ($`M_{\mathrm{up}}80`$ M$``$) have to be present in the cluster, and that the cluster is $`3`$ Myr or $`4.5`$ Myr old. Further constraints on the upper mass limit cut-off will be imposed by the relavite number of Wolf-Rayet with respect to O stars in the cluster.
Burst models have been computed also for different values of the IMF slope ($`\alpha `$=2.35, 3.0 and 1.5). The behaviour of the emission line ratios with the age is very similar for these three values of the IMF slope (Fig. 9). The observed ratios indicate an age $`3`$ Myr or 4.5 Myr, and cannot discriminate between the different values of the slope. Further constraints on the IMF slope will come from the WR/O ratio. This point will be discussed further in section 6.
## 4 Modeling the H Balmer and He i absorption lines
The optical continuum of an RH ii is dominated by early-type stars. The spectra of O and B stars are characterized by strong H Balmer and He i absorption lines, with very weak metallic lines formed in the photosphere of these stars (Walborn & Fitzpatrick 1990). The H Balmer and He i recombination nebular emission lines are superposed on the corresponding photospheric lines. However, the higher order terms of the Balmer series and some of the He i lines can be detected in absorption and the lower terms of the Balmer series can show absorption wings. The detection of these absorption features depends on the spectral and spatial resolution of the observations, on the spatial distribution of the stellar cluster with respect to the nebular emission, and on the evolutionary state of the stellar cluster. These photospheric features have been detected in a spectrum of NGC 604 (Fig. 5) because it corresponds to the zone of the nebular hole where the core of the central cluster is located, so it maximizes the contrast of the stellar with respect to the nebular contributions (see Fig. 2). In this section we will constrain the evolutionary state of NGC 604 and the IMF using the profile of the higher-order terms of the H Balmer series and the strength of some of the He i lines in absorption.
### 4.1 Description of the models
González Delgado et al. (1999b) have computed evolutionary stellar population synthesis models that predict the photospheric absorption H Balmer and He i lines, between 3700 and 5000 Å, for a single-metallicity stellar population. The models, which are optimized for galaxies with active star formation, synthesize the profiles of the H Balmer series (H$`\beta `$, H$`\gamma `$, H$`\delta `$, H8, H9, H10, H11, H12 and H13) and the He i absorption lines (He i $`\lambda `$4922, He i $`\lambda `$4471, He i $`\lambda `$4388, He i $`\lambda `$4144, He i $`\lambda `$4121, He i $`\lambda `$4026, He i $`\lambda `$4009 and He i $`\lambda `$3819), with a spectral sampling of 0.3 Å pixel<sup>-1</sup>, for a burst and for continuous star formation at a constant rate. They use a stellar library that includes NLTE absorption profiles for stars hotter than 25000 K, and LTE profiles for lower temperatures. The temperature and gravity coverage is $`4000T_{\mathrm{eff}}50000`$ K, and $`0.0logg5.0`$, respectively (González Delgado & Leitherer 1999). The models assume that stars evolve from the main sequence following the evolutionary tracks of the Geneva group (Schaller et al. 1992; Schaerer et al. 1993a,b; Charbonnel et al. 1993; Meynet et al. 1994). The strength of the Balmer and He i lines is sensitive to the age after the first 3 Myr of evolution, and sensitive to the IMF if the age is younger than 3–4 Myr. Models assume that stars have a metallicity $`Z=0.008`$, and make different assumptions about the IMF slope ($`\alpha `$=1.5, 2.35 and 3.0), and upper mass limit cut-off ($`M_{\mathrm{up}}=80`$ M$``$ and 30 M$``$). For the purpose of this paper, the range of ages computed spans from 0 to 10 Myr.
### 4.2 Model Results
We use the equivalent width of He i $`\lambda `$4388, $`\lambda `$4026, and $`\lambda `$3819, and the wing absorption profiles of H$`\beta `$, H$`\gamma `$, H$`\delta `$, H8, H9 and H10 to constrain the age, IMF and star formation law. He i $`\lambda `$4388, $`\lambda `$4026, and $`\lambda `$3819 are detected in absorption, because the corresponding nebular emission lines have equivalent widths that are at least a factor 3 weaker than the nebular emission in He i $`\lambda `$4471 (González Delgado et al. 1999b). Note that He i $`\lambda `$4471 is only partially filled with the nebular emission; therefore, the nebular emission of the lines He i $`\lambda `$4388, He i $`\lambda `$4026 and He i $`\lambda `$3819 is $``$0, and the equivalent widths of these absorption features represent well the strength of the stellar continuum radiation of the cluster. The equivalent widths of these lines, He i $`\lambda `$4388, He i $`\lambda `$4026 and He i $`\lambda `$3819, measured in the spectrum of Fig. 5 are 0.32, 0.55 and 0.26 Å, respectively. Their strength is compatible with continuous star formation and with burst models 3 Myr old (Fig. 10). However, He i lines can constrain the IMF. Models with an IMF steeper than Salpeter or with $`M_{\mathrm{up}}=30`$ M$``$ predict He i lines stronger than observed. The strengths of the lines are compatible with a 3 to 4 Myr old burst formed following a Salpeter or slightly flatter IMF ($`\alpha `$= 1.5, Fig. 10). The absorption wings of the Balmer lines are also compatible with burst models 3 Myr old, with Salpeter or slightly flatter IMF ($`\alpha `$= 1.5, Fig. 11).
## 5 Modelling the ultraviolet stellar lines
The ultraviolet light from an RH ii is dominated by O stars. These hot stars develop strong wind stellar lines due to the radiation pressure in ultraviolet resonance lines (Morton 1967). As a result, all the strong ultraviolet lines (e.g. O vi $`\lambda `$1034, N v $`\lambda `$1240, Si iv $`\lambda `$1400, C iv $`\lambda `$1550 and N iv $`\lambda `$1720) show a blueshifted absorption (about 2000–3000 km s<sup>-1</sup>) or a PCygni profile. The shape of the profile reflects the stellar mass-loss rate, which is related to the stellar luminosity, and thus to the stellar mass. Most of the ultraviolet spectra of RH ii are dominated by absorption features (Rosa et al. 1984; Vacca et al. 1995; Mas-Hesse & Kunth 1999), without any nebular emission, which is very similar to those of starburst galaxies<sup>6</sup><sup>6</sup>6One important exception is knot A in the RH ii NGC 2363 (Drissen et al. 2000), so young that the effect of winds is not yet reflected in its uv spectrum.. Most of these lines are formed in the stellar winds of the massive stars that belong to the starburst. The profile of these lines reflects the stellar massive content in the starburst; therefore, they depend on the IMF and star formation law (Leitherer, Robert & Heckman 1995). In this section, we constrain the evolutionary state and the IMF of NGC 604 by means of fitting the profiles of the ultraviolet stellar lines Si iv and C iv.
### 5.1 Description of the models
Evolutionary stellar population models have been computed with the code Starburst 99 (Leitherer et al. 1999). The code uses a stellar library of IUE ultraviolet spectra of O, B and Wolf-Rayet stars (Robert, Leitherer & Heckman 1993). The spectral resolution of the O and Wolf-Rayet stars observed is 0.75 Å. These stars are located in the solar neighborhood. However, the evolutionary models computed here assume that the stars evolve from the main sequence following the evolutionary track at $`Z=0.008`$. Models are computed for instantaneous bursts between 0 and 10 Myr, and for continuous star formation lasting 10 Myr, and different assumptions about the slope ($`\alpha `$=2.35, 3.0 and 1.5) and upper mass limit cut-off ($`M_{\mathrm{up}}=`$120, 80, 60 and 40 M$``$) of the IMF. The model spectra are smoothed to the IUE spectral resolution of the observations of NGC 604, which is 6 Å.
### 5.2 Model results
The strongest wind stellar features in the spectra of NGC 604 are N v $`\lambda `$1240, Si iv $`\lambda `$1400, C iv $`\lambda `$1550, He ii $`\lambda `$1640, and N iv $`\lambda `$1720. These features are present in all the spectra of NGC 604; however, for this analysis we use the average spectrum of those at PA=149$`\mathrm{°}`$, 136$`\mathrm{°}`$ and 159$`\mathrm{°}`$, and the spectrum at PA=32$`\mathrm{°}`$, because they have the best signal to noise ratio. The evolutionary state, the star formation law and the IMF in NGC 604 is constrained with the profiles of these lines.
He ii shows a broad emission profile in Wolf-Rayet and O3–O5 supergiant stars; N iv shows a PCygni profile in Wolf-Rayet stars; N v and C iv show strong PCygni profiles in all O stars, and Si iv only in O supergiants. Thus, the profile of these lines in the integrated spectrum of a starburst depends strongly on the stellar content and age of the stellar cluster. In particular, Si iv shows a strong PCygni profile if the cluster formed in an instantaneous burst and its age is between 3 and 5 Myr, because within this age interval the ultraviolet light is dominated by O blue supergiants. On the other hand, the profile of C iv depends strongly on the IMF. It shows a strong PCygni profile if stars more massive that 60 M$``$ are formed in the cluster, and if the slope of the IMF is flatter than $`\alpha `$=3.0.
The spectrum of NGC 604 indicates that its stellar population must be dominated by massive, young O stars; thus, the stellar cluster must be young and of short duration. Continuous star formation models can be ruled out because they show Si iv weaker than observed (Fig. 12). In these models the line is diluted because the fraction of O supergiants with respect to the total number of O stars is lower in the continuous star formation than in the instantaneous burst models. The strength of Si iv in NGC 604 indicates that it is an instantaneous burst. The age of the burst has to be between 3 to 5 Myr, because bursts younger than 3 Myr or older than 5 Myr have very few O supergiants, and thus very weak Si iv (Fig. 13). On the other hand, from the strength of C iv we can exclude instantaneous bursts with $`M_{\mathrm{up}}60`$ M$``$ (Fig. 14), and IMF steeper than $`\alpha `$=3.0 (Fig. 15). Therefore, we conclude that continuous star formation and instantaneous burst models with few very massive stars can be ruled out. The wind lines are compatible with an instantaneous burst, formed following a Salpeter or slightly flatter IMF, with upper mass limit cut-off higher than 60 M$``$. A 3 Myr instantaneous burst with Salpeter and $`M_{\mathrm{up}}80`$ M$``$ fits well the profile of the wind stellar lines.
## 6 The massive stellar content of NGC 604
We derive the massive stellar population in NGC 604 from the ultraviolet continuum luminosity, comparing the observations with the predictions of synthetic models. The number of Wolf-Rayet stars are derived from the luminosity of the He ii $`\lambda `$1640 line. The stellar content derived in this way is compared with that derived from HST photometry of stars in NGC 604. However, we need first to estimate the extinction.
### 6.1 Extinction estimates
Leitherer & Heckman (1995) have shown that the ultraviolet continuum arising from a young starburst ($`\mathrm{age}10\mathrm{Myr}`$) has a spectral index, $`\beta `$ (F$`{}_{\lambda }{}^{}\lambda ^\beta `$), which is independent of the IMF, star formation law and metallicity. Thus, any deviation from the predicted value, $`2.5`$ for a 3 Myr instantaneous burst, could be attributed to reddening.
The ultraviolet continuum flux distribution of the IUE spectra of NGC 604 shows evidence of reddening, because after correcting by Galactic extinction, $`E(BV)=0.03`$ (McClure & Racine 1969), the spectra are flatter than the spectral energy distribution predicted by the evolutionary synthesis models. Massey & Hutchings (1983), analysing IUE spectra of H ii regions in M33, conclude that the M33 extinction curve is significantly different from that of the Galaxy, because the spectra of the H ii regions show very weak 2200 Å interstellar absorption dip, resembling those of the LMC and SMC. Therefore, we use the LMC curve to match the observed spectra with the spectral energy distribution of a 3 Myr old instantaneous burst. The $`E(BV)`$ derived is 0.1 for the spectrum at PA=32$`\mathrm{°}`$, and 0.12 for the average spectrum of those at PA=136$`\mathrm{°}`$, 159$`\mathrm{°}`$ and 149$`\mathrm{°}`$. These values are in agreement with those derived by Massey & Hutchings (1983). However, the extinction can be 0.03 higher than the values above if it is derived comparing the observed spectra with the spectral energy distribution predicted including only the stellar contribution. This is due to the nebular contribution being dominated by the two-photon continuum emission that peaks at $``$1500 Å. If the nebular continuum is not included, the spectral energy distribution is steeper and a higher extinction is required to match the observed and predicted ultraviolet flux distribution. The reason to do this new estimation is because the nebula is more extended than the IUE aperture; thus, only a fraction of the nebular continuum contributes to the IUE aperture.
We have estimated also the extinction using the Calzetti, Kinney & Storchi-Bergmann (1994) extinction law. This is an empirical extinction curve derived for starbursts, which is very similar to the MW extinction law but it does not show the 2200 Å bump. Using this curve, we estimate $`E(BV)=0.2`$ to match the observed spectra and the spectral energy distribution predicted by the models.
To distinguish between these two results, we compare the spectral energy distribution predicted by the evolutionary models with that of NGC 604. We build the spectral energy distribution of NGC 604 from the ultraviolet to near-infrared using the SWP+LWR IUE spectrum at PA=110$`\mathrm{°}`$ plus that at PA=149$`\mathrm{°}`$, and the average optical spectra scanning the inner 11$`\times `$60 arcsec (at PA=120$`\mathrm{°}`$), and 14$`\times `$60 arcsec (at PA=60$`\mathrm{°}`$). We use the $`B`$ (F439W) and $`I`$ (F814W) HST+WFPC2 images to flux calibrate the optical spectra, because the mode in which the scanning was performed precludes an accurate absolute flux calibration based only on the ground-based observations. The resulting optical and ultraviolet spectra are normalized dividing by the flux at 2900 Å. Then, the spectra are dereddened by $`E(BV)=0.1`$ using the LMC extinction law, and by $`E(BV)=0.2`$ using the Calzetti et al. (1994) extinction law. Even though apparently the spectra dereddened with Calzetti’s law seem to match better the spectral energy distribution, a change in the scaling calibration factor of only 10 percent (larger than that used) gives a better match with the spectrum dereddened by the LMC law. Thus, given the uncertainties in the flux calibration of the ground-based optical spectra, we cannot really distinguish between these two extinction laws.
### 6.2 The number of O stars and W/O ratio
We assume that the total ultraviolet flux of the region is approximately the sum of the fluxes at PA=110$`\mathrm{°}`$ and at PA=149$`\mathrm{°}`$. After correcting by Galactic extinction, the flux at 1500 Å is $`5.1\times 10^{13}`$ erg s<sup>-1</sup> cm<sup>-2</sup> Å<sup>-1</sup>. The intrinsic luminosity<sup>7</sup><sup>7</sup>7Note that this luminosity and the quantities derived in Table 4 should be almost a factor two lower if the total ultraviolet flux is equal to the flux of only one of the IUE apertures, or to that measured in the F170W HST image, $`L_{1700}7\times 10^{37}`$ erg s<sup>-1</sup> Å<sup>-1</sup>. is 10<sup>38.16</sup> erg s<sup>-1</sup> Å<sup>-1</sup> (10<sup>38.32</sup> erg s<sup>-1</sup> Å<sup>-1</sup>) if the flux is corrected by $`E(BV)=0.12`$ (0.2) using the LMC (Calzetti et al. 1994) extinction law.
The mass of the cluster, the ionizing photon luminosity and the number of O stars predicted by the evolutionary synthesis models are given in Table 4. These quantities indicate that very massive stars must be present in the cluster and that it is very young. In fact, instantaneous bursts older than 3.5 Myr predict an ionizing photon luminosity which is at least a factor 3 lower than the value derived from the H$`\alpha `$ flux, which is $`logQ=51.54`$ (s<sup>-1</sup>). In contrast, a $``$3 Myr instantaneous burst (with $`M_{\mathrm{up}}80\mathrm{M}`$) reproduces well the ionizing photon luminosity<sup>8</sup><sup>8</sup>8However the values of $`Q`$ predicted are still a factor 1.5 lower than the photon luminosity derived from the H$`\alpha `$ flux, suggesting that other sources in addition to the central cluster contribute to the ionization of the gas, and/or that the extinction is higher than $`E(BV)=0.1`$. Note that if log $`L_{1500}=38.32`$ ergs<sup>-1</sup> Å<sup>-1</sup>, then all the quantities derived in Table 4 should be a factor $`1.5`$ larger.). These models predict a mass of the stellar cluster that ranges between 0.1 and $`2\times 10^5`$ M$``$, depending on the assumption about the upper mass limit cut-off and slope of the IMF. The number of O stars ranges between 150 and 215 (if $`M_{\mathrm{up}}80\mathrm{M}`$), which is in agreement with the number of O stars (186) reported by Hunter et al.(1996).
The slope and upper mass limit cut-off of the IMF can also be constrained with the observed Wolf-Rayet over O ratio, WR/O. HST images of the NGC 604 cluster have detected 14 Wolf-Rayet or Of candidates (Drissen et al. 1993). Thus, WR/O$``$0.075. This ratio is very high, suggesting that the cluster is young (3–3.5 Myr), and that it must contain very massive stars. Instantaneous burst models with the IMF slope steeper than Salpeter or $`M_{\mathrm{up}}80`$ M$``$ predict a ratio much lower than 0.075 (Fig. 17). Models with IMF flatter than Salpeter or with upper mass limit cut-off as high as 120 M$``$ reproduce well the observed value.
### 6.3 Deposition of kinetic energy in the interstellar medium
Massive stars not only interact with the interstellar medium via their radiation, but also depositing kinetic energy via their stellar winds. The energy released during their lifetime is comparable to that deposited by a supernova event (Leitherer, Robert & Drissen 1992). The evolutionary state and the stellar content derived above for NGC 604 suggests that the massive stars in NGC 604 can release enough kinetic energy to form many of the filaments and expanding shell structures observed in the region.
Yang et al. (1996) have reported the properties of five shells; these have expansion velocities that range between 40 and 125 km s<sup>-1</sup>, and sizes between 35 and 125 pc (see also Sabalisck et al. 1995). The wind power required to form these bubbles ranges between 2$`\times 10^{37}`$ erg s<sup>-1</sup> and 9$`\times 10^{38}`$ erg s<sup>-1</sup>, with a sum total of $`10^{39.25}`$ erg s<sup>-1</sup>.
We estimate the wind power produced by the stellar cluster in NGC 604 using Starburst99 (Table 4). For an instantaneous burst, the wind power ranges between 3$`\times 10^{37}`$ erg s<sup>-1</sup> and 4$`\times 10^{38}`$ erg s<sup>-1</sup>, depending on the assumptions about the slope and massive cut-off of the IMF. Thus, the energy relased by winds in NGC 604 is not enough to explain the formation of the five shells; however, it can explain the formation of the central hole (shell number 3) and shell number 4 (see fig. 2 in Yang et al. 1996), that encircles the core of NGC 604 cluster (cf. Fig. 2). The wind power required by these two bubbles is 5$`\times 10^{37}`$ erg s<sup>-1</sup> and 2$`\times 10^{37}`$ erg s<sup>-1</sup>, respectively.
## Summary and conclusions
The main goal of this work is to constrain the evolutionary state of the giant H ii region NGC 604. For this purpose, we have analyzed the integrated ultraviolet spectra taken by IUE, and optical ground-based spectra of the region. The data are interpreted using evolutionary synthesis models optimized for star forming regions. These data are complemented with ultraviolet and H$`\alpha `$ images taken by HST with the WFPC2.
The ultraviolet image shows that the ionizing cluster is spatially spread in the inner 20$`\times `$20 arcsec, with the core of the cluster within a central nebular hole shell structure. The optical spectrum, as it is well known, is dominated by nebular emission lines from the surrounding photoionized medium. In contrast, the ultraviolet spectrum of the region is dominated by absorption lines formed in the stellar winds of massive stars. However, other photospheric stellar lines (the high order terms of the Balmer series and He i lines) are detected at optical wavelengths near the Balmer jump. The spatial distribution of the stellar cluster with respect to the nebular emission, allows to detect these lines in absorption in a spectrum corresponding to the inner 4 arcsec of the region, where the stellar light is maximized with respect to the nebular contribution.
The evolutionary state and the massive stellar content of the region is derived in a self-consistent way using evolutionary synthesis and photoionization models, applied to the ultraviolet resonance wind stellar lines, to the nebular emission-lines and to the higher-order terms of the Balmer series and He i lines in absorption. The three techniques applied suggest that the central ionizing cluster of NGC 604 is very young, 3 Myr old, with no evidence for an age spread. The overall properties suggest that the massive stars in the cluster were formed following a Salpeter or flatter IMF, with presence of stars more massive than 80 M$``$.
Particular results from the modelling of the nebular emission lines include:
\- The nebula is well described by a sphere of inner radius 20 pc and outer radius 110 pc (determined by the ionization front). This value is in agreement with the value derived from the surface brightness photometry of the H$`\alpha `$ emission.
\- The sphere is partially filled (filling factor $``$0.1) with ionized gas that is very tenuous, having electron density of $``$30 cm<sup>-3</sup>.
\- \[O i\]6300 and \[O iii\]4363 emissions are well accounted for by photoionization.
Particular results from the modelling of the ultraviolet continuum are:
\- The extinction affecting the stellar cluster is little, $`E(BV)=0.1`$, if the LMC extinction law is used to derive the value, or $`E(BV)=0.2`$ if the Calzetti et al. (1994) extinction law is used. The 3 Myr instantaneous burst spectral energy distribution is well matched by the spectrum of NGC 604 corrected with any of these two extinction laws. The value of the extinction is similar to the average extinction, $`A_V=0.5`$, derived for the gas by Churchwell & Goss (1999).
\- The massive cluster provides a number of high energy photons that is enough to photoionize the whole nebula. Thus, within the limits of the integrated models that we develop for a uniform geometry, our results are compatible with most of the ionizing radiation being reprocessed in the nebula, with no significant escape of ionizing photons.
\- The wind power provided by the massive stars of the cluster is enough to form the central hole structure where the core stellar cluster is located. However, it cannot provide all the wind power required for the formation of the (at least) five shell structures seen in the H$`\alpha `$ image.
\- The number of massive stars estimated is in agreement with that derived from photometric studies based on the detection of individual stars.
## Acknowledgments
It is a pleasure to thank Gary Ferland and Claus Leitherer for kindly making their codes available, and Jaime Perea for the use of SIPL. Claus suggested that something should be done with those IUE spectra. We have benefited from estimulating and helpful discussions with members of the GEFE collaboration, in particular with Angeles Díaz, Jesús Maíz-Apellániz, Divakara Mayya, Guillermo Tenorio-Tagle, Elena Terlevich, Roberto Terlevich and José M. Vílchez. We are also indebted to Grazyna Stasińska for his detailed comments from a thorough reading of the paper. The IUE helpdesk have always been solicituous with our queries, in particular we acknowledge Pedro Rodríguez Pascual and Rosario González Riestra for support with the IUE spectra. HST images were retrieved from the ST-ECF HST archive.
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# 1 Introduction
## 1 Introduction
Chiral fermions in even dimensions can give rise to an anomaly . Although there is no chirality in odd dimensions, there is still a similar phenomenon for fermions. It is well-known that due to radiative effects at one fermion loop, a Chern-Simons action for the gauge bosons can be generated in odd dimensions . This phenomenon has wide applications in particle and condensed matter physics .
Recently, noncommutative field theory has shown up as effective description of string theory in a certain background . The noncommutativity takes the form
$$[x^\mu ,x^\nu ]=i\theta ^{\mu \nu },$$
(1)
where $`\theta ^{\mu \nu }`$ is a antisymmetric real constant matrix and is of dimension length squared. In the dual language, the algebra of functions is described by the Moyal product
$$(fg)(x)=e^{i\frac{\theta ^{\mu \nu }}{2}\frac{}{\xi ^\mu }\frac{}{\zeta ^\nu }}f(x+\xi )g(x+\zeta )|_{\xi =\zeta =0},$$
(2)
which is associative, noncommutative and satisfies
$$\overline{(fg)}=\overline{g}\overline{f}$$
(3)
under complex conjugation. We note also that under integration
$$fg=gf=fg,$$
(4)
which is a consequence of momentum conservation. Using this $``$-product, field theory on a noncommutative spacetime can be easily formulated. One simply needs to replace the usual multiplication of functions by the Moyal product. Pioneering analysis of this kind of noncommutative field theory was performed by Filk . Aspects of noncommutative field theories was further developed in -.
In this paper we analyze gauge theory with fermions defined on a noncommutative odd dimensional spacetime and determine the corresponding induced Chern-Simons action. Since the action for the noncommutative theory is a smooth deformation of the classical action, one may natively expect to get back the usual commutative description in the limit $`\theta 0`$ . We find that this is not always the case when one is in the quantum regime: for the case of a Majorana spinor, there is a jump in the induced Chern-Simons action. Singularities in $`\theta `$ have also been displayed in the scalar theories and in QED .
## 2 Induced Chern-Simons in odd dimensions
We will begin with noncommutative QED in (2+1)-dimensions with a 2-components massless fermion. We will take $`\theta _{12}=\theta 0`$, and the action is given by
$$S=d^3x(\frac{1}{4}F_{\mu \nu }F^{\mu \nu }+i\overline{\psi }D/\psi ),$$
(5)
where $`F_{\mu \nu }=_\mu A_\nu _\nu A_\mu +ig[A_\mu ,A_\nu ]_{}`$ and the covariant derivative (and hence the fermion coupling) is given by $`D_\mu \psi =_\mu \psi +igA_\mu \psi `$ for a Dirac spinor and $`D_\mu \psi =_\mu \psi +ig[A_\mu ,\psi ]_{}`$ for a Majorana spinor. These couplings reproduce the correct corresponding commutative limit, in particular a Majorana spinor is neutral in the commutative case.
The $`\gamma `$-matrices are given by the Pauli matrices and satisfies
$$\gamma ^\mu \gamma ^\nu =g^{\mu \nu }iϵ^{\mu \nu \lambda }\gamma _\lambda ,$$
(6)
where $`g^{\mu \nu }=(,+,+)`$ and $`ϵ^{012}=+1`$. The action is invariant under the gauge transformation
$$\delta A_\mu =_\mu \alpha ig[\alpha ,A_\mu ]_{},$$
(7)
and $`\delta \psi =ig\alpha \psi `$ or $`\delta \psi =ig[\alpha ,\psi ]_{}`$ for a Dirac or Majorana spinor respectively.
The one fermion-loop effective action is given by
$`S_{eff}[A,m]=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{d^3p}{(2\pi )^3}A_\mu (p)A_\nu (p)(i\mathrm{\Gamma }^{\mu \nu }(p))}+`$ (8)
$`+{\displaystyle \frac{1}{3}}{\displaystyle \frac{d^3p_1}{(2\pi )^3}\frac{d^3p_2}{(2\pi )^3}A_\mu (p_1)A_\nu (p_2)A_\lambda (p_1p_2)(i\mathrm{\Gamma }^{\mu \nu \lambda }(p_1,p_2))}`$
for $`m=0`$. The 2-point and 3-point functions $`\mathrm{\Gamma }_{\mu \nu }(p)`$ and $`\mathrm{\Gamma }_{\mu \nu \lambda }(p)`$ will be analyzed now.
### 2.1 Dirac fermions
We first consider the case of a Dirac spinor. It is crucial to observe that for Dirac spinors, one only gets planar diagrams from doing the Wick contractions. Hence the 2-point and 3-point functions are
$$i\mathrm{\Gamma }_D^{\mu \nu }(p)=g^2\frac{d^3k}{(2\pi )^3}\mathrm{tr}[\gamma ^\mu \frac{k/p/m}{(kp)^2+m^2}\gamma ^\nu \frac{k/m}{k^2+m^2}],$$
(9)
$$i\mathrm{\Gamma }_D^{\mu \nu \lambda }(p_1,p_2)=g^3\frac{d^3k}{(2\pi )^3}\frac{\mathrm{tr}[\gamma ^\mu (k/m)\gamma ^\nu (k/+p_2/m)\gamma ^\lambda (k/p_1/m)]}{(k^2+m^2)((k+p_2)^2+m^2)((kp_1)^2+m^2)}e^{\frac{i}{2}p_1\theta p_2}.$$
(10)
The only difference from the commutative case is the phase factor in (10) which depends only on the external momenta.
The effective action (8) is ultraviolet divergent and needs to be regularized. This can be achieved by the standard Pauli-Villars method:
$$S_{eff}^{reg}[A]=S_{eff}[A,m=0]\underset{m\mathrm{}}{lim}S_{eff}[A,m].$$
(11)
As in the commutative case, apart from the wanted terms that cure the divergences of (8), there are terms that remain even after sending $`m\mathrm{}`$ and these terms give rises to the induced Chern-Simons Lagrangian. One can verify that in the large $`m`$ limit:
$`\underset{m\mathrm{}}{lim}i\mathrm{\Gamma }_{\mu \nu }^D(p)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Lambda }}{3\pi ^2}}g_{\mu \nu }{\displaystyle \frac{i}{4\pi }}{\displaystyle \frac{m}{|m|}}ϵ_{\mu \nu \lambda }p^\lambda ,`$ (12)
$`\underset{m\mathrm{}}{lim}i\mathrm{\Gamma }_{\mu \nu \lambda }^D(p_1,p_2)`$ $`=`$ $`{\displaystyle \frac{i}{4\pi }}{\displaystyle \frac{m}{|m|}}ϵ_{\mu \nu \lambda }e^{\frac{i}{2}p_1\theta p_2}.`$ (13)
Substituting back to (8), we obtain the following fermion-loop induced term in $`S_{eff}^{reg}`$,
$$S_{ind}=\pm \frac{1}{2}S_{CS}$$
(14)
with
$$S_{CS}=\frac{1}{4\pi }d^3xϵ^{\mu \nu \lambda }(g^2A_\mu _\nu A_\lambda +\frac{2i}{3}g^3A_\mu A_\nu A_\lambda ),$$
(15)
or in terms of the gauge field $`𝒜_\mu =igA_\mu `$
$$S_{CS}=\frac{1}{4\pi }d^3xϵ^{\mu \nu \lambda }(𝒜_\mu _\nu 𝒜_\lambda +\frac{2}{3}𝒜_\mu 𝒜_\nu 𝒜_\lambda )$$
(16)
is the noncommutative Chern-Simons action. $`S_{CS}`$ is local invariant under (7), while under a finite gauge transformation,
$$𝒜_\mu h_{}^1𝒜_\mu h+h_{}^1_\mu h$$
(17)
$`S_{CS}`$ changes as <sup>1</sup><sup>1</sup>1 We have dropped a total derivative piece $`d(Adhh_{}^1)`$ which on a manifold with boundary will give rises to an anomaly for the boundary theory. Descent relations and the algebraic structures of chiral anomaly still hold generally for a noncommutative gauge theory . See also the third and fourth references of for recent discussion.
$$S_{CS}S_{CS}2\pi w$$
(18)
where
$$S_{WZW}:=\frac{1}{24\pi ^2}_Bd^3xϵ^{\mu \nu \lambda }(h_{}^1_\mu hh_{}^1_\nu hh_{}^1_\lambda h)$$
(19)
is the noncommutative WZ term over $`B`$ (or WZW action over $`B`$) and $`w:=S_{WZW}`$ for $`B=S^3`$ is the “winding number”. Here $`h_{}^1`$ is the inverse of $`h`$ with respect to the Moyal product:
$$hh_{}^1=h_{}^1h=1.$$
(20)
For any $`h=e^{i\alpha }`$ of $`U(1)`$, it is $`h_{}^1=e^{i\alpha }`$. It is easy to check that $`w`$ is invariant under an infinitesimal transformation $`\delta h=\lambda h`$ and that under a finite transformation,
$$w(h^{}h)=w(h^{})+w(h).$$
(21)
Now we claim that $`w`$ is zero for the Abelian case. Consider
$$I[\alpha ]=ϵ^{\mu \nu \lambda }d^3xe^{i\alpha }_\mu e^{i\alpha }e^{i\alpha }_\nu e^{i\alpha }e^{i\alpha }_\lambda e^{i\alpha }.$$
(22)
It is clear from (21) that $`I[\frac{m}{n}\alpha ]=\frac{m}{n}I[\alpha ]`$ for any integers $`m,n`$, therefore $`I[s\alpha ]=sI[\alpha ]`$ for any real constant $`s`$, and hence
$$I[\alpha ]=\frac{}{s}I[s\alpha ].$$
(23)
But we already knew that $`I`$ is invariant under arbitrary infinitesimal transformation, hence our claim. One can also get the same result by noticing that since the LHS of (23) is independent of $`s`$, one can evaluates the RHS of (23) at the particular value $`s=0`$ and obtains the desired result. That $`w=0`$ for $`U(1)`$ may be expected intuitively since $`S^3`$ is too big to fit in $`S^1`$. As a result, the 2-dimensional noncommutative WZW action is well defined.
For the case of non-Abelian $`U(N)`$ gauge fields, we just have to replace (8) by
$`S_{eff}[A,m]={\displaystyle \frac{1}{2}}\mathrm{tr}(T^aT^b){\displaystyle \frac{d^3p}{(2\pi )^3}A_\mu ^a(p)A_\nu ^b(p)(i\mathrm{\Gamma }^{\mu \nu }(p))}`$
$`+{\displaystyle \frac{1}{3}}\mathrm{tr}(T^aT^bT^c){\displaystyle \frac{d^3p_1}{(2\pi )^3}\frac{d^3p_2}{(2\pi )^3}A_\mu ^a(p_1)A_\nu ^b(p_2)A_\lambda ^c(p_1p_2)(i\mathrm{\Gamma }^{\mu \nu \lambda }(p_1,p_2))}`$ (24)
with the same 2-point and 3-point functions. Thus one again obtains (14), now with
$$S_{CS}=\frac{1}{4\pi }d^3xϵ^{\mu \nu \lambda }\mathrm{tr}(𝒜_\mu _\nu 𝒜_\lambda +\frac{2}{3}𝒜_\mu 𝒜_\nu 𝒜_\lambda ),$$
(25)
where $`𝒜_\mu =igA_\mu ^aT^a`$.
Again, $`S_{CS}`$ is local gauge invariant, while under a finite gauge transformation (17) with $`h`$ in $`U(N)`$, $`S_{CS}`$ changes as
$$S_{CS}S_{CS}2\pi w$$
(26)
where
$$S_{WZW}:=\frac{1}{24\pi ^2}ϵ^{\mu \nu \lambda }_B\mathrm{tr}(h_{}^1_\mu hh_{}^1_\nu hh_{}^1_\lambda h)$$
(27)
is the noncommutative WZ term over $`B`$ and $`w:=S_{WZW}`$ for $`B=S^3`$. It is easy to check that $`w`$ is invariant under a local $`U(N)`$ gauge transformation $`\delta h=\lambda h`$ and we have again (21) under a finite transformation. It is not clear whether $`w`$ is an integer or not when integrated on $`S^3`$. One can shows that this number is independent of $`\theta `$, at least to the first order in $`\theta `$. Since $`w[h^n]=nw[h]`$ and $`w`$ is invariant under small changes of the map $`h`$, these suggest that $`w`$ may again serve as some sort of homotopy invariant and that the $`\theta `$-dependence factorize. If $`\theta `$ dependence does not disappear and $`w`$ is not an integer, then one will need to arrange the fermions content so that the global anomaly cancel. We leave these interesting issues for future investigation.
### 2.2 Majorana fermions
We next consider the case of a Majorana spinor. The situation differs in that now nonplanar diagram can also contribute. Similar considerations have also been made in . The Feynman rules have been worked out in and we will not repeat them here. We note that the phase factor $`e^{iq_1\frac{\theta }{2}q_2}`$ for a Dirac spinor coupled to the gauge field (with momentum $`q_1,q_2,q_3`$ coming into the vertex) is now replaced by $`2i\mathrm{sin}(q_1\frac{\theta }{2}q_2)`$ due to the commutator nature of the coupling for the Majorana spinors. In particular the 2-point and 3-point functions are now given by
$$i\mathrm{\Gamma }_M^{\mu \nu }(p)=4g^2\frac{d^3k}{(2\pi )^3}\mathrm{tr}[\gamma ^\mu \frac{k/p/m}{(kp)^2+m^2}\gamma ^\nu \frac{k/m}{k^2+m^2}]\mathrm{sin}^2(\frac{\stackrel{~}{p}k}{2}),$$
(28)
$`i\mathrm{\Gamma }_M^{\mu \nu \lambda }(p_1,p_2)=8ig^3{\displaystyle \frac{d^3k}{(2\pi )^3}\frac{\mathrm{tr}[\gamma ^\mu (k/m)\gamma ^\nu (k/+p_2/m)\gamma ^\lambda (k/p_1/m)]}{(k^2+m^2)((k+p_2)^2+m^2)((kp_1)^2+m^2)}}`$
$`\mathrm{sin}({\displaystyle \frac{\stackrel{~}{p}_1k}{2}})\mathrm{sin}({\displaystyle \frac{\stackrel{~}{p}_2k}{2}})\mathrm{sin}({\displaystyle \frac{\stackrel{~}{p}_3(k+p_2)}{2}}),`$ (29)
where for convenience we have denoted $`\stackrel{~}{p}=p\theta `$ and momentum conservation
$$p_1+p_2+p_3=0$$
(30)
has to be used. Since
$$\mathrm{sin}^2(\frac{\stackrel{~}{p}k}{2})=\frac{1}{2}(1\mathrm{cos}\stackrel{~}{p}k),$$
(31)
one can reduce (28) to a sum corresponding to planar and nonplanar contributions
$$\mathrm{\Gamma }_M^{\mu \nu }=2\mathrm{\Gamma }_D^{\mu \nu }+\text{“nonplanar”},$$
(32)
where “nonplanar” are expressions of the form
$$\frac{d^3k}{(2\pi )^3}\frac{f(k)}{((kp)^2+m^2)(k^2+m^2)}e^{ik\stackrel{~}{p}}$$
(33)
and $`f(k)=k^\mu k^\nu `$, $`k^\mu `$ or 1. It is easy to show that the “nonplanar” terms all vanish in the large $`m`$ limit. It is enough to calculate for the case of $`f=1`$ as the others can be obtained by differentiating with respect to $`\stackrel{~}{p}`$. Introducing Schwinger parameters, one obtains
$`{\displaystyle \frac{d^3k}{(2\pi )^3}\frac{1}{((kp)^2+m^2)(k^2+m^2)}e^{ik\stackrel{~}{p}}}`$ $`=`$ $`{\displaystyle \frac{1}{(2\sqrt{\pi })^3}}{\displaystyle _0^{\mathrm{}}}𝑑T{\displaystyle \frac{1}{T^{1/2}}}e^{m^2T\frac{\stackrel{~}{p}^2}{4T}}{\displaystyle _0^1}𝑑xe^{Tx(1x)p^2}`$ (34)
$`<`$ $`{\displaystyle _0^{\mathrm{}}}dT{\displaystyle \frac{1}{T^{1/2}}}e^{m^2T\frac{\stackrel{~}{p}^2}{4T}}{\displaystyle \frac{1}{m}}e^{|m||\stackrel{~}{p}|}`$
and hence
$$\underset{m\mathrm{}}{lim}\text{“nonplanar”}=0$$
(35)
as long as $`|\stackrel{~}{p}|=|\theta |\sqrt{p_1^2+p_2^2}0`$, which is of measure zero in the integration of (8). Similarly one gets using (30)
$$i\mathrm{sin}(\frac{\stackrel{~}{p}_1k}{2})\mathrm{sin}(\frac{\stackrel{~}{p}_2k}{2})\mathrm{sin}(\frac{\stackrel{~}{p}_3(k+p_2)}{2})=\frac{i}{4}\mathrm{sin}(\frac{p_1\theta p_2}{2})+\text{phases involving internal momenta}.$$
(36)
Therefore again we can write (2.2) as
$$\mathrm{\Gamma }_M^{\mu \nu \lambda }=2\mathrm{\Gamma }_D^{\mu \nu \lambda }+\mathrm{`}\mathrm{`}\text{nonplanar}^{\prime \prime }$$
(37)
and similarly show that the “nonplanar” terms vanish in the large $`m`$ limit. Substituting back to (8) and remember that there is now an extra factor of $`1/2`$ since we are considering Majorana spinor and hence $`S_{eff}=1/2\mathrm{t}\mathrm{r}\mathrm{log}(iD/+A/)`$. We obtain finally the induced term
$$S_{ind}=\pm \frac{1}{2}S_{CS}.$$
(38)
We note that the factor of 1/2 and 1/4 in (31) and (36) can be easily understood. There are one planar and one nonplanar diagram contributing to the 2-point function, and one planar and three nonplanar diagrams contributing to the 3-point function (given a definite cyclic order of the external momentum). We also note that in the commutative case, Majorana spinors are neutral and there is no Chern-Simons action induced from them. In the noncommutative case, a coupling can be written down that reduces to zero in the commutative limit. Both planar and nonplanar diagrams contribute to the effective actions, but the nonplanar diagrams are suppressed and only the planar diagrams survive the large $`m`$ limit. Therefore upon removing the Pauli-Villars regulator, we are left with the induced Chern-Simons action (38).
For the $`U(N)`$ case, the coupling is replaced by
$$ig(f^{abc}\mathrm{cos}(q_1\frac{\theta }{2}q_2)+d^{abc}\mathrm{sin}(q_1\frac{\theta }{2}q_2))$$
(39)
where we have adopted the normalization $`\mathrm{tr}T^aT^b=\delta ^{ab}/2`$ and $`T^aT^b=\frac{1}{2}(if^{abc}T^c+d^{abc}T^c)`$. The $`U(1)`$ generator is $`T^0=\frac{1}{\sqrt{2N}}1`$ and it is $`f^{000}=0`$ and $`d^{000}=\sqrt{2/N}`$. The previous $`U(1)`$ case is recovered by $`gd^{000}2g`$. Using the identities
$`f^{ab^{}a^{}}f^{bc^{}b^{}}f^{ca^{}c^{}}=f^{ab^{}a^{}}d^{bc^{}b^{}}d^{ca^{}c^{}}={\displaystyle \frac{N}{2}}f^{abc},`$ (40)
$`d^{ab^{}a^{}}d^{bc^{}b^{}}d^{ca^{}c^{}}={\displaystyle \frac{N}{2}}d^{abc}+\delta ^{ab}\mathrm{tr}T^c+\delta ^{bc}\mathrm{tr}T^a+\delta ^{ca}\mathrm{tr}T^b`$ (41)
$`f^{ab^{}a^{}}f^{bc^{}b^{}}d^{ca^{}c^{}}={\displaystyle \frac{N}{2}}d^{abc}\delta ^{ab}\mathrm{tr}T^c+\delta ^{bc}\mathrm{tr}T^a+\delta ^{ca}\mathrm{tr}T^b`$ (42)
where $`T^a`$ are in the defining representation. One easily obtains the induced Chern-Simons term (38), with now $`S_{CS}`$ (25) defined in the adjoint representation.
Finally we remark that all of the above can be generalized straightforwardly to higher dimensions (provided that massive Majorana spinors exist in that dimensions ) and one still obtains the same shift with the corresponding higher dimensional Chern-Simons action.
## 3 Discussions
In this paper, we investigated the induced Chern-Simons action due to Dirac and Majorana fermion coupled to Abelian $`U(1)`$ or non-Abelian $`U(N)`$ gauge fields. The surprising result is that for the Majorana spinor case, the induced term does not go to zero as $`\theta 0`$ and displays a discontinuity. Since there is a finite difference no matter how small is the noncommutativity, this kind of discontinuity in a physical quantity may be useful to searching of experimental signals of noncommutativity and Lorentz symmetry breaking in nature.
Similar $`\theta `$-singularities have also been discovered in scalar field theories and QED . These examples show that although classically noncommutative physics is a smooth deformation of the commutative description, there can be important new quantum mechanical terms that do not vanish in the commutative limit. The commutative limit and the classical limit generally do not commute. The technical reason is very simple, generally one cannot exchange the order of taking the $`\theta 0`$ limit and the integration.
One may then get the opposite impression that noncommutative quantum physics always do not have a smooth commutative limit. The example of the induced Chern-Simons term for Dirac spinors shows that this is not true. So long as only planar diagrams contribute, there is a smooth limit. Take another example, the computation of the chiral anomaly for Weyl fermion. Again only planar diagrams contribute and so the loop effects get modified only by a phase factor which depend on the external momentum. Therefore it can be Fourier transformed back easily and gives the usual expression of the anomaly, except that now the products are replaced by the Moyal products. There is again a smooth limit because only planar diagrams contribute.
The opposite situation would be more interesting: is there a case where only nonplanar diagrams contribute? From the string point of view, we need a string theory with only nonplanar worldsheets. One needs to identify a limit so that only nonplanar worldsheets survive. Consistency may be an important issue, but it may also be possible that the field theory limit is not sensitive to the details of the string consistency issues.
It is known that in the commutative case, the Chern-Simons action can also be induced from the gauge boson loops. This bosonic shift has been shown explicitly in the $`(2+1)`$-dimensional case by Witten using a saddle point analysis. The situation is however less clear if one perform a perturbative calculation and the shift depends on the regularization scheme. A natural regularization is to use string theory. It would be interesting to see what string regularization has to say about the shift for both the commutative and noncommutative case.
Another interesting question concerns the nonrenormalization of our result. In the commutative case, it has been shown through explicit calculations that the fermion induced Chern-Simons term is not renormalized at two loops , for both the Abelian and non-Abelian case. Beyond this result, we have the Coleman-Hill theorem for the Abelian case which shows that there is no contribution to the induced Chern-Simons action beyond one loop. Due to the topological origin of the non-Abelian Chern-Simons action and its connection with the chiral anomaly, it is also expected that it should not be renormalized beyond one-loop. Thus the situation is more or less clear. It would be interesting to investigate the status of nonrenormalization theorem for induced Chern-Simons action, as well as for chiral anomaly.
Induced Chern-Simons term in odd dimensions were originally obtained for fermions coupled to gauge fields in a flat spacetime. This has been extended to full generality for arbitrary curved backgrounds and any odd dimensions and is shown to be related to the Atiyah-Patodi-Singer index theorem. The induced Chern-Simons action in $`(2n+1)`$-dimensions is given (up to a normalization factor) by the secondary characteristic class $`Q(A,\omega )`$ satisfying
$$dQ(A,\omega )=\widehat{A}(R)ch(F)|_{2n+2},$$
(43)
where $`\omega `$ is the gravitational connection. In view of the results on the mathematical structure of anomalies in noncommutative gauge theory, it seem appropriate to investigate the deformation theory for elliptic operators due to a Moyal product and to study the possible “topological” meaning of the Chern-Simons action and chiral anomaly and to establish the corresponding index theorems. Better understanding of the properties of the noncommutative WZW action and its topological meaning is also highly desired.
It would also be interesting to investigate aspects of AdS/CFT correspondence in relation to chiral anomaly and Chern-Simons action. This will provide another channel to understand better this so far rather poorly understood correspondence. We plan to come back to these issues in the future.
Acknowledgments
I would like to thank A. Bilal, J.-P. Derendinger, Rodolfo Russo and Bruno Zumino for many useful discussions and comments. This work was partially supported by the Swiss National Science Foundation, by the European Union under TMR contract ERBFMRX-CT96-0045 and by the Swiss Office for Education and Science.
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# CDM models with a steplike initial power spectrum
## 1 INTRODUCTION
The theory of galaxy formation based on gravitational instability describes how primordially generated fluctuations grow into galaxies and clusters of galaxies due to self-gravity of matter. The initial field of density fluctuations $`\delta (𝐱,t)`$ can be decomposed into its Fourier components $`\delta _𝐤(t)`$ and expressed in terms of the power spectrum $`P(k)=|\delta _𝐤|^2`$.
As the study of the large-scale structure in the Universe has pushed to ever larger scales, several data samples have suggested the presence of a peak in the power spectrum at the wavenumber $`k0.05h`$ Mpc<sup>-1</sup> (or at the wavelength $`\lambda 120130h^1`$ Mpc). Einasto et al. (1997) and Retzlaff et al. (1998) studied the spatial distribution of the Abell-ACO clusters and found that the power spectrum of the clusters has a well-defined peak at the wavenumber $`k=0.052h`$ Mpc<sup>-1</sup>. A similar peak in the one-dimensional power spectrum of a deep pencil-beam survey was detected by Broadhurst et al. (1990), and in the two-dimensional power spectrum of the Las Campanas redshift survey by Landy et al. (1996). The power spectrum of the spatial distribution of APM galaxies also has a feature on the same scale (Caztanaga and Baugh 1998). Independent evidence for the presence of a preferred scale in the Universe at about $`130h^1`$ Mpc comes from an analysis of high-redshift galaxies. Broadhurst and Jaffe (1999) studied the distribution of the Lyman-break galaxies at redshift $`z3`$ and found a $`5\sigma `$ excess of pairs separated by $`\mathrm{\Delta }z=0.22\pm 0.02`$, equivalent to $`130h^1`$ Mpc for flat universe with the density parameter $`\mathrm{\Omega }_0=0.4\pm 0.1`$.
The power spectrum of density fluctuations depends on the physical processes in the early universe. The peak in the power spectrum of clusters at the wavenumber $`k0.05h`$ Mpc<sup>-1</sup> may be generated during the era of radiation domination or earlier. Standard cosmological models based on collisionless dark matter \[e.g. cold dark matter (CDM)\] and adiabatic fluctuations, when combined with power-law initial power spectra, predict smooth power spectra of density fluctuations at $`z10^3`$. The baryonic acoustic oscillations in adiabatic models may explain the observed power spectrum only if currently favored determinations of cosmological parameters are in substantial error (e.g., if the density parameter $`\mathrm{\Omega }_0<0.2h`$; Eisenstein et al. 1998).
One possible explanation for the observed power spectrum of clusters is an inflationary model with a scalar field whose potential $`V(\phi )`$ has a local steplike feature in the first derivative. This feature can be produced by a fast phase transition in a physical field different from the inflaton field. An exact analytical expression for the scalar (density) perturbations generated in this inflationary model was found by Starobinsky (1992). The initial power spectrum of density fluctuations in this model can be expressed as
$$P_{in}(k)\frac{kS(k,p)}{p^2},$$
$`(1)`$
where the function $`S(k,p)`$ can be written as
$$S(k,p)=1\frac{3(p1)}{y}\left[f_1(y)\mathrm{sin}2y+\frac{2}{y}\mathrm{cos}2y\right]+$$
$$+\frac{9(p1)^2f_2(y)}{2y^2}\left[f_2(y)+f_1(y)\mathrm{cos}2y\frac{2}{y}\mathrm{sin}2y\right].$$
$`(2)`$
Here, the function $`y=k/k_0`$, $`f_1(y)=1y^2`$ and $`f_2(y)=1+y^2`$. The initial power spectrum in this model depends on two parameters $`k_0`$ and $`p`$. The parameter $`k_0`$ determines the location of the step and the parameter $`p`$ \- the shape of the initial spectrum. For $`p=1`$, we recover the scale-invariant Harrison-Zel’dovich spectrum ($`S(k,1)1`$). At present, the initial spectrum (1,2) is probably the only example of an initial power spectrum with the desired properties, for which a closed analytical form exists.
Fig. 1 shows the function $`S(k,p)/p^2`$ for different values of the parameter $`p`$. For the models with $`p<1`$ and $`p>1`$, the step parameter was chosen to be $`k_0=0.016h`$ Mpc<sup>-1</sup> and $`k_0=0.03h`$ Mpc<sup>-1</sup>, respectively. In this case, in the models with $`p<1`$, the power spectrum has a well-defined maximum at the wavenumber $`k0.05h`$ Mpc<sup>-1</sup> and a second maximum at $`k0.1h`$ Mpc<sup>-1</sup>. In the models with $`p>1`$, the picture is inverted. The power spectrum has a flat upper plateau at the wavenumbers $`k<0.05h`$ Mpc<sup>-1</sup>, a sharp decrease on smaller scales ($`k=0.050.1h`$ Mpc<sup>-1</sup>) and a secondary maximum at $`k0.15h`$ Mpc<sup>-1</sup>.
Lesgourgues, Polarski & Starobinsky (1998, hereafter LPS) compared the CDM models with a steplike initial spectrum (1,2) with observational data. They studied the rms mass fluctuation on an $`8h^1`$ Mpc scale, $`\sigma _8`$; the rms bulk velocity of galaxies and the cosmic microwave background (CMB) anisotropies on different angular scales. In this paper we continue this research and investigate the properties of galaxy clusters in these models. We examine the mass function, the peculiar velocities and the power spectrum of clusters for different values of the density parameter $`\mathrm{\Omega }_0`$, the normalized Hubble constant $`h`$ and the spectral parameter $`p`$. The results are compared with observations. We also investigate the rms bulk velocity and $`\sigma _8`$ in the models, where the mass function, the peculiar velocities and the power spectrum of clusters are consistent with the observed data.
In their study LPS assumed that the parameter $`\sigma _8`$ lies in the interval $`\sigma _8=(0.57\pm 0.06)\mathrm{\Omega }_0^{0.56}`$. This interval was derived by White, Efstathiou & Frenk (1993) by analysing the mass function of clusters. For $`\mathrm{\Omega }_0=0.3`$ this gives $`\sigma _8=1.01.25`$. In this paper we examine the observed values of the mass function of galaxy clusters in more detail and obtain lower values for the parameter $`\sigma _8`$. For $`\mathrm{\Omega }_0=0.3`$ we find that $`\sigma _8=0.81.05`$. Therefore, the allowed values for the parameter $`p`$ in the ($`\mathrm{\Omega }_0`$, $`h`$) plane that we obtain are different from those found by LPS.
We examine flat cosmological models with the density parameter $`\mathrm{\Omega }_0=0.20.5`$ and the normalized Hubble constant $`h=0.50.8`$. These parameters are in agreement with measurements of the density parameter (e.g. Bahcall et al. 1999) and with measurements of the Hubble constant using various various distance indicators (e.g. Tammann 1998). To restore the spatial flatness in the low-density models, we assume a contribution from a cosmological constant: $`\mathrm{\Omega }_\mathrm{\Lambda }=1\mathrm{\Omega }_0`$. The Hubble constant is written as $`H_0=100h`$ km s<sup>-1</sup>Mpc<sup>-1</sup>. The transfer function $`T(k)`$ is computed with the fast Boltzmann code CMBFAST developed by Seljak & Zaldarriaga (1996). The models are normalized using the COBE normalization derived by Bunn & White (1997). We assume that the initial density fluctuation field in the Universe is a Gaussian field. In this case, the power spectrum provides a complete statistical description of the field.
This paper is organized as follows. In Section 2 we study the mass function of clusters of galaxies and compare the results with observations. In Section 3 we examine the peculiar velocities of galaxy clusters. In Section 4 we investigate the power spectrum of clusters. Discussion and summary are presented in Section 5.
## 2 THE MASS FUNCTION OF CLUSTERS OF GALAXIES
To study the mass function of clusters we use the Press-Schechter (1974, PS) approximation. The PS mass function has been compared with N-body simulations (Efstathiou et al. 1988; White, Efstathiou & Frenk 1993; Lacey & Cole 1994; Eke, Cole & Frenk 1996; Borgani et al. 1997a) and has been shown to provide an accurate description of the abundance of virialized cluster-size halos. In the PS approximation the number density of clusters with the mass between $`M`$ and $`M+dM`$ is given by
$$n(M)dM=\sqrt{\frac{2}{\pi }}\frac{\rho _b}{M}\frac{\delta _t}{\sigma ^2(M)}\frac{d\sigma (M)}{dM}\mathrm{exp}\left[\frac{\delta _t^2}{2\sigma ^2(M)}\right]dM.$$
$`(3)`$
Here $`\rho _b`$ is the mean background density and $`\delta _t`$ is the linear theory overdensity for a uniform spherical fluctuation which is now collapsing; $`\delta _t=1.686`$ for $`\mathrm{\Omega }_0=1`$, with a weak dependence on $`\mathrm{\Omega }_0`$ for flat and open models (Eke et al. 1996; Kitayama & Suto 1996). In this paper we will use values of $`\delta _t`$ found by Eke et al. (1996) for flat models. The function $`\sigma (M)`$ is the rms linear density fluctuation at the mass scale $`M`$. We will use the top-hat window function to describe halos. For the top-hat window, the mass $`M`$ is related to the window radius $`R`$ as $`M=4\pi \rho _bR^3/3`$. In this case, the number density of clusters of mass larger than $`M`$ can be expressed as
$$n_{cl}(>M)=_M^{\mathrm{}}n(M^{})𝑑M^{}=$$
$$=\frac{3}{(2\pi )^{3/2}}_R^{\mathrm{}}\frac{\delta _t}{\sigma ^2(r)}\frac{d\sigma (r)}{dr}\mathrm{exp}\left[\frac{\delta _t^2}{2\sigma ^2(r)}\right]\frac{dr}{r^3}.$$
$`(4)`$
Fig. 2 shows the cluster mass function in the flat models with $`\mathrm{\Omega }_0=0.3`$ and $`h=0.65`$. The threshold density $`\delta _t=1.675`$. We investigated the cluster masses within a $`1.5h^1`$ Mpc radius sphere around the cluster center. This mass $`M_{1.5}`$, is related to the window radius $`R`$ as
$$R=8.43\mathrm{\Omega }_0^{\frac{0.2\alpha }{3\alpha }}\left[\frac{M_{1.5}}{6.99\times 10^{14}\mathrm{\Omega }_0h^1M_{}}\right]^{\frac{1}{3\alpha }}(h^1\mathrm{Mpc}).$$
$`(5)`$
Here the parameter $`\alpha `$ describes the cluster mass profile, $`M(r)r^\alpha `$, at radii $`r1.5h^1`$ Mpc. Numerical simulations and observations of clusters indicate that the parameter $`\alpha 0.60.7`$ for most of clusters (Navarro, Frenk & White 1995; Carlberg, Yee & Ellingson 1997). In this paper we use a value $`\alpha =0.65`$.
Fig. 2 shows also the mass function of clusters of galaxies derived by Bahcall and Cen (1993, BC) and by Girardi et al. (1998, G98). BC used both optical and X-ray observed properties of clusters to determine the mass function of clusters. The function was extended towards the faint end using small groups of galaxies. G98 determined the mass function of clusters by using virial mass estimates for 152 nearby Abell-ACO clusters including the ENACS data (Katgert et al. 1998). The mass function derived by G98 is somewhat larger than the mass function derived by BC, the difference being larger at larger masses (see Fig. 2).
Let us consider the amplitude of the mass function of galaxy clusters at $`M_{1.5}=410^{14}h^1M_{}`$. For this mass, the cluster abundances derived by BC and G98 are $`n(>M)=(2.0\pm 1.1)10^6h^3`$ Mpc<sup>-3</sup> and $`n(>M)=(6.3\pm 1.2)10^6h^3`$ Mpc<sup>-3</sup>, respectively. By analysing X-ray properties of clusters, White, Efstathiou & Frenk (1993) found that the number density of clusters with the mass $`M_{1.5}4.210^{14}h^1M_{}`$ is $`n(>M)=410^6h^3`$ Mpc<sup>-3</sup>.
We derived the limits for the parameter $`p`$, assuming that the mass function of galaxy clusters at $`M_{1.5}=410^{14}h^1M_{}`$ is in the range $`(26.5)10^6h^3`$ Mpc<sup>-3</sup>. Fig. 3a shows the results for the models with $`h=0.65`$ for different $`\mathrm{\Omega }_0`$ and Fig. 3b for the models with $`\mathrm{\Omega }_0=0.3`$ for different $`h`$. For the model with $`\mathrm{\Omega }_0=0.3`$ and $`h=0.65`$, we find that $`p=0.791.0`$. Fig. 2 demonstrates the mass function of clusters in this model for $`p=0.79`$ and $`p=1.0`$.
LPS used high values of the parameter $`\sigma _8`$ and found that one of the best-fit models is the model with parameters $`\mathrm{\Omega }_0=0.3`$, $`h=0.7`$, $`p=0.8`$. We find that the number density of clusters in this model is substantially higher than observed; for the mass $`M_{1.5}=410^{14}h^1M_{}`$, the cluster abundance $`n(>M)=9.8\times 10^6h^3`$ Mpc<sup>-3</sup>. For $`\mathrm{\Omega }_0=0.3`$ and $`h=0.7`$, the mass function of clusters is consistent with the observed data, if $`p=0.881.13`$.
## 3 PECULIAR VELOCITIES OF CLUSTERS OF GALAXIES
The observed rms peculiar velocity of galaxy clusters has been studied in several papers (e.g. Bahcall, Gramann & Cen 1994, Bahcall and Oh 1996, Borgani et al. 1997b, Watkins 1997). In this paper we use the results obtained by Watkins (1997). He developed a likelihood method for estimating the rms peculiar velocity of clusters from line-of-sight velocity measurements and their associated errors. This method was applied to two observed samples of cluster peculiar velocities: a sample known as the SCI sample (Giovanelli et al. 1997) and a subsample of the Mark III catalog (Willick et al. 1997). Watkins (1997) found that the rms one-dimensional cluster peculiar velocity is $`265_{75}^{+106}`$ km s<sup>-1</sup>, which corresponds to the three-dimensional rms velocity $`459_{130}^{+184}`$ km s<sup>-1</sup>.
To investigate the peculiar velocities of clusters in our models, we use the linear theory predictions for peculiar velocities of peaks in the Gaussian field. The linear rms velocity fluctuation on a given scale $`R`$ at the present epoch can be expressed as
$$\sigma _v(R)=H_0f(\mathrm{\Omega }_0)\sigma _1(R),$$
$`(6)`$
where $`f(\mathrm{\Omega }_0)\mathrm{\Omega }_0^{0.56}`$ is the linear velocity growth factor in the flat models and $`\sigma _j`$ is defined for any integer $`j`$ by
$$\sigma _j^2=\frac{1}{2\pi ^2}P(k)W^2(kR)k^{2j+2}𝑑k.$$
$`(7)`$
Bardeen et al. (1986) showed that the rms peculiar velocity at peaks of the smoothed density field differs systematically from $`\sigma _v(R)`$, and can be expressed as
$$\sigma _p(R)=\sigma _v(R)\sqrt{1\sigma _0^4/\sigma _1^2\sigma _1^2}.$$
$`(8)`$
Suhhonenko & Gramann (1999) examined the linear theory predictions for the peculiar velocities of peaks and compared these to the peculiar velocities of clusters in N-body simulations. The N-body clusters were determined as peaks of the density field smoothed on the scale $`R1.5h^1`$ Mpc. The numerical results showed that the rms peculiar velocity of small clusters is similar to the linear theory expectations, while the rms peculiar velocity of rich clusters is higher than that predicted in the linear theory. The rms peculiar velocity of clusters with a mean cluster separation $`d_{cl}=30h^1`$ Mpc was $`18`$ per cent higher than that predicted by the linear theory. We assume that the observed cluster sample studied by Watkins (1997) corresponds to the model clusters with a separation $`d_{cl}30h^1`$ Mpc ($`n_{cl}3.710^5h^3`$ Mpc<sup>-3</sup>) and determine the rms peculiar velocity of the clusters, $`v_{cl}`$, as
$$v_{cl}=1.18\sigma _p(R),$$
$`(9)`$
where the radius $`R=1.5h^1`$ Mpc.
Fig. 3 shows the limits for $`p`$ in different models obtained on the basis of Watkins (1997) results. Fig. 3a shows the results in the $`h=0.65`$ models for different $`\mathrm{\Omega }_0`$ and Fig. 3b in the $`\mathrm{\Omega }_0=0.3`$ models for different $`h`$. In the region studied in Fig. 3b, the rms peculiar velocity is larger than 329 km s<sup>-1</sup>. For the model with $`\mathrm{\Omega }_0=0.3`$ and $`h=0.65`$, we find that the initial parameter $`p>0.87`$.
For high values of $`\mathrm{\Omega }_0`$ and $`h`$, peculiar velocities are larger than observed for any values of the parameter $`p`$. The velocities are sensitive to the amplitude of the large-scale fluctuations at wavenumbers $`k<0.1h`$ Mpc<sup>-1</sup>. By increasing $`p`$, if $`p>1`$, the power spectrum on these wavenumbers and, therefore, the velocities remain almost unchanged (see Fig. 1). Fig. 3a shows that in the $`h=0.65`$ models with $`\mathrm{\Omega }_0>0.35`$, the peculiar velocities are not consistent with observations. For the $`\mathrm{\Omega }_0=0.4`$ model with $`h=0.5`$ and $`h=0.6`$, we find that peculiar velocities are consistent with observations, if $`p>0.86`$ and $`p>1.26`$, respectively.
If we compare the observational constraints obtained by studying the mass function and peculiar velocities of clusters of galaxies, we see from Fig. 3 that the mass function and the peculiar velocities of clusters are consistent with the observed data only in a small interval of the parameter $`p`$. For the model with $`\mathrm{\Omega }_0=0.3`$ and $`h=0.65`$, we find that $`p=0.871.0`$. For the $`\mathrm{\Omega }_0=0.3`$ models with $`h=0.5`$ and $`h=0.6`$, we find that $`p=0.570.64`$ and $`p=0.760.88`$, respectively.
## 4 THE POWER SPECTRUM OF CLUSTERS OF GALAXIES
Let us now consider the power spectrum of clusters in these models where both the mass function and the peculiar velocities are consistent with observations. To investigate the power spectrum of clusters, we can also use the PS approach. The power spectrum of clusters for a given number density in the PS approximation can be expressed as (Gramann & Suhhonenko 1999, hereafter GS)
$$P_{cl}(k)=b_{cl}^2P(k),$$
$`(9)`$
where the cluster bias parameter $`b_{cl}`$ is
$$b_{cl}=1\frac{3}{(2\pi )^{3/2}n_{cl}}_R^{\mathrm{}}\frac{1}{\sigma ^2(r)}\frac{d\sigma (r)}{dr}(y^21)\mathrm{exp}(y^2)\frac{dr}{r^3}.$$
$`(10)`$
Here, the function $`y=\delta _t/\sigma (r)`$ and $`\sigma (r)`$ is the rms linear density fluctuation at the radius $`r`$. The cluster bias parameter $`b_{cl}`$ depends on the minimal mass $`M`$ (or the window radius $`R`$) of clusters and on the power spectrum of density fluctuations, $`P(k)`$, which determines the function $`\sigma (r)`$. For fixed $`P(k)`$ and $`n_{cl}`$, the minimal mass $`M`$ (or scale $`R`$) can be determined by inverting equation (4).
Observations provide the distribution of clusters in redshift space, which is distorted due to peculiar velocities of clusters. On large scales, where linear theory applies, the power spectrum of matter density fluctuations in redshift space is given by (Kaiser 1987):
$$P^s(k)=\left[1+\frac{2f(\mathrm{\Omega }_0)}{3}+\frac{f^2(\mathrm{\Omega }_0)}{5}\right]P(k).$$
$`(11)`$
Using the PS approximation (10), relation (11) takes the form
$$P_{cl}^s(k)=\left[1+\frac{2f(\mathrm{\Omega }_0)}{3b_{cl}}+\frac{f^2(\mathrm{\Omega }_0)}{5b_{cl}^2}\right]b_{cl}^2P(k).$$
$`(12)`$
Equation (12) determines the power spectrum of clusters for a given $`n_{cl}`$ in redshift space.
GS examined the power spectrum of clusters in the PS theory and in N-body simulations. They determined the power spectrum of clusters for mean separations $`d_{cl}=3040h^1`$ Mpc. The numerical results showed that at wavenumbers $`k<0.1h^1`$ Mpc, the power spectrum of clusters in the simulations is linearly enhanced with respect to the power spectrum of the matter distribution. However, the amplitude of the spectrum of clusters was somewhat lower than predicted by the approximation (12). It is possible that the linear approximation (12) overestimates the power spectrum of clusters due to dynamical effects that are not taken into account in this approximation. The power spectrum of clusters can be expressed as
$$P_{cl}^s(k)=F\left[1+\frac{2f(\mathrm{\Omega }_0)}{3b_{cl}}+\frac{f^2(\mathrm{\Omega }_0)}{5b_{cl}^2}\right]b_{cl}^2P(k),$$
$`(13)`$
where the factor $`F=0.70.8`$. The factor $`F`$ depends slightly on the model. GS examined the distribution of clusters in two cosmological models which start from the observed power spectra of the distribution of galaxies and clusters of galaxies. In the model (1), the initial linear power spectrum of density fluctuations was chosen in the form $`P(k)k^2`$ at wavelengths $`\lambda <120h^1`$ Mpc. In the model (2), GS assumed that the initial power spectrum contains a primordial feature at the wavelengths $`\lambda 3060h^1`$ Mpc. They found that $`F=0.8`$ and $`F=0.7`$ in the model (1) and model (2), respectively. In this paper we use a value $`F=0.75`$.
Fig. 4 shows the redshift-space power spectrum of the model clusters with a mean separation $`d_{cl}=34h^1`$ Mpc. For comparison, we show the power spectrum of the Abell-ACO clusters determined by Einasto et al. (1999). This spectrum represents the weighted mean of the power spectra determined by Einasto et al. (1997) and Retzlaff et al. (1998). Einasto et al. (1997) determined the power spectrum of the Abell-ACO clusters from the correlation function of clusters, while Retzlaff et al. (1998) estimated the power spectrum directly (see Einasto et al. 1999 for details). The power spectrum of the distribution of the Abell-ACO clusters peaks at the wavenumber $`k=0.052h`$ Mpc<sup>-1</sup>. The mean intercluster separation of the Abell-ACO clusters is $`d_{cl}34h^1`$ Mpc ($`n_{cl}2.510^5h^3`$ Mpc<sup>-3</sup>) (Einasto et al. 1997, Retzlaff et al. 1998).
Fig. 4 shows the power spectrum of clusters predicted in the $`\mathrm{\Omega }_0=0.3`$ models with ($`h=0.5`$, $`p=0.6`$), ($`h=0.6`$, $`p=0.8`$), ($`h=0.65`$, $`p=0.95`$) and ($`h=0.7`$, $`p=1.05`$). For these models, the mass function and the peculiar velocities are consistent with the observed data (see Fig. 3b). The models with $`p=0.6`$ and $`p=0.8`$ are in good agreement with the observed power spectrum, while the models with $`p=0.95`$ and $`p=1.05`$ are not consistent with the observed data. In the models with $`p>0.95`$, we do not see a peak in the power spectrum at the wavenumber $`k0.05h`$ Mpc<sup>-1</sup>.
We used also the $`\chi ^2`$ test to calculate the probability that the models fit the observed power spectrum. The observed power spectrum represents the weighted mean of the power spectra determined by Einasto et al. (1997) and Retzlaff et al. (1998) and it is not clear how many points of the power spectrum are independent. As a first step we used all the data points in Fig. 4. The models with $`p=0.6`$ and $`p=0.8`$ are consistent with the observed power spectrum at a confindence level higher than $`90`$%. For the models with $`p=0.95`$ and $`p=1.05`$, the probability that the models fit the observed power spectrum is less than $`30`$%.
We studied also the power spectrum of clusters in other models with different values of ($`\mathrm{\Omega }_0`$, $`h`$) and found similar results. The power spectrum of clusters is in good agreement with the observed power spectrum of the Abell-ACO clusters, if the initial parameter $`p`$ is in the range $`p=0.60.8`$. In the models with $`h=0.65`$, the allowed values for the parameter $`p`$ are in this range, if $`\mathrm{\Omega }_0=0.200.27`$ (see Fig. 3a). Similarly, the $`\mathrm{\Omega }_0=0.3`$ models are consistent with the observed mass function, peculiar velocities and the power spectrum of clusters, if $`h=0.500.63`$ (Fig. 3b).
In the $`\mathrm{\Omega }_0=0.4`$ model with $`h=0.5`$, the peculiar velocities are larger than observed, if $`p<0.86`$. For higher values of $`\mathrm{\Omega }_0`$ and $`h`$, this limit for $`p`$ increases. Therefore, for $`\mathrm{\Omega }_00.4`$ and $`h0.5`$, the observed power spectrum and peculiar velocities of clusters are not consistent with each other. Either the observed peculiar velocities are underestimated, or the observed peak in the power spectrum of clusters is overestimated. The peculiar velocities and the power spectrum of clusters are consistent with observations only if $`\mathrm{\Omega }_0<0.4`$.
## 5 DISCUSSION AND SUMMARY
In this paper, we have examined the properties of clusters of galaxies in the $`\mathrm{\Lambda }`$CDM models with a steplike initial power spectrum (1,2) that depends on two parameters $`k_0`$ and $`p`$. The parameter $`k_0`$ determines the location of the step and the parameter $`p`$ \- the shape of the initial spectrum. For the models with $`p<1`$ and $`p>1`$, the step parameter was chosen to be $`k_0=0.016h`$ Mpc<sup>-1</sup> and $`k_0=0.03h`$ Mpc<sup>-1</sup>, respectively. We investigated the mass function, peculiar velocities and the power spectrum of clusters in models with different values of $`\mathrm{\Omega }_0`$, $`h`$ and $`p`$.
We found that the mass function and the peculiar velocities of clusters are consistent with the observed data only in a small interval of the parameter $`p`$. For the model with $`\mathrm{\Omega }_0=0.3`$ and $`h=0.65`$, we find that $`p=0.871.0`$. For the $`\mathrm{\Omega }_0=0.3`$ models with $`h=0.5`$ and $`h=0.6`$, we find that $`p=0.570.64`$ and $`p=0.760.88`$, respectively.
The power spectrum of clusters in the $`\mathrm{\Lambda }`$CDM models with a steplike initial power spectrum is in good agreement with the observed power spectrum of the Abell-ACO clusters, if the initial parameter $`p`$ is in the range $`p=0.60.8`$. In the models with $`h=0.65`$, the allowed values for the parameter $`p`$ are in this range if $`\mathrm{\Omega }_0=0.200.27`$. The $`\mathrm{\Omega }_0=0.3`$ models are consistent with the observed mass function, peculiar velocities and the power spectrum of clusters, if $`h=0.500.63`$. The peculiar velocities and the power spectrum of clusters are consistent with observations only if $`\mathrm{\Omega }_0<0.4`$.
We also studied the rms bulk velocity in the models, where the mass function, the peculiar velocities and the power spectrum of clusters are consistent with observations. We choosed three models in the allowed parameter space and studied these models in more detail. Tabel 1 lists the parameters of the models studied. We have determined the number density of clusters with mass $`M_{1.5}=410^{14}h^1M_{}`$, $`n_{cl}`$; the rms mass fluctuation on an $`8h^1`$ Mpc scale, $`\sigma _8`$; the rms peculiar velocity of clusters, $`v_{cl}`$, and the rms bulk velocity for a radius $`r=50h^1`$ Mpc, $`V_{50}`$. The rms bulk velocity was determined by using eq. (6). The observed bulk velocities are determined in a sphere centered on the Local Group and represent a single measurement of the bulk flow on large scales. The observed bulk velocity derived from the Mark III catalog of peculiar velocities for $`r=50h^1`$ Mpc is $`375\pm 135`$ km s<sup>-1</sup> (Kolatt & Dekel 1997). In the models studied, the rms bulk velocity is $`270285`$ km s<sup>-1</sup>, which is consistent with the observed data.
Therefore, in many aspects the $`\mathrm{\Lambda }`$CDM models with a steplike initial spectrum fit the observed data. Further work is needed to study the properties of galaxies and clusters of galaxies in these models in more detail.
## ACKNOWLEDGEMENTS
We thank A. Starobinsky, J. Einasto and E. Saar for useful discussions. This work has been supported by the ESF grant 3601.
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# On the Radii of Close-in Giant Planets
## 1. Introduction
The recent indirect detections of extrasolar giant planets (EGPs) by Doppler spectroscopy have taught us that such planetary systems can be very much unlike our own. To date, nearly $``$30 EGPs have been discovered around stars with spectral types from M4 to F7 (e.g., Mayor & Queloz (1995); Marcy et al. 1999a ; Fischer et al. (1999)). The planets themselves have minimum masses (M<sub>p</sub>$`\mathrm{sin}(i)`$) between $``$0.42 M$`_\mathrm{J}`$and $``$10 M$`_\mathrm{J}`$, orbital semi-major axes from $``$0.042 A.U. to $``$3.8 A.U., and eccentricities as high as $``$0.71. Such variety vastly expands the parameter space within which both theorists and observers must operate in defining the physical character of extrasolar planetary systems.
The most interesting, unexpected, and problematic subclass of EGPs are those found within $``$0.1 A.U. of their primaries, 50–100 times closer than Jupiter is to our Sun. At such orbital distances, due to prodigious stellar irradiation alone, an EGP can have an effective temperature (T$`_{\mathrm{eff}}`$) greater than 600 K. Indeed, the EGPs HD187123b, HD209458b, $`\tau `$ Boo b, HD75289b, 51 Peg b, $`\upsilon `$ And b, and HD217107b likely all have T<sub>eff</sub>s above 1000 K. This is to be compared with T<sub>eff</sub>s for Jupiter and Saturn of 125 K and 95 K, respectively. The minimum masses of these close-in EGPs imply that they occupy the high-T$`_{\mathrm{eff}}`$/low-gravity corner of parameter space for which the compositions and composition profiles may be unique (Seager and Sasselov (1998); Goukenleuque et al. (1999)).
However, high stellar fluxes on a close-in EGP can have profound structural consequences for the planet. In particular, stellar insolation can be responsible for maintaining the planet’s radius at a value 20% to 80% larger than that of Jupiter itself (Guillot et al. (1996)). This prediction has recently been verified by the observation of the transit by HD209458b of its primary. The depth of the transit yields values for its radius that range from $``$1.27 R$`_\mathrm{J}`$(Charbonneau et al. 1999b ) to $``$1.7 R$`_\mathrm{J}`$(Henry et al. (1999)), with a best value near $``$1.4 R$`_\mathrm{J}`$(Mazeh et al. (2000)). Importantly, since HD209458b transits its primary, astronomers can derive $`\mathrm{sin}(i)`$, from which the planet’s mass can be directly determined. This is a major advance in the emerging study of extrasolar planets. In this paper, we focus on the HD209458 system and what broadly can be concluded theoretically from these new transit data. We do not provide detailed models from which one can extract bulk or atmospheric composition. Rather, we show that the large radius of HD209458b is a consequence of the retardation by stellar irradiation of the otherwise natural cooling of the convective core of the planet and that such a radius requires that the planet did not dwell for long, if at all, at large orbital distances after its formation. Hence, the large radius of a close-in EGP is not due to the thermal expansion of its atmosphere, but to the high residual entropy that remains throughout its bulk as a consequence of its early proximity to a luminous primary.
Recently, Cameron et al. (1999), using spectral deconvolution, claim to have seen $`\tau `$ Boo b in reflection. To investigate this exciting possibility, we calculate a range of theoretical radii for the $`\tau `$ Boötis planet. As we show, the Cameron et al. value of 1.6–1.8 R$`_\mathrm{J}`$for $`\tau `$ Boo b’s radius, if verified, is a challenge to the still embryonic theory of massive, close-in EGPs. Nevertheless, with the discovery of a transiting planet, HD209458b, with the maturation of the technique of spectral deconvolution, and with the anticipated development of adaptive optics and interferometry for the direct study of extrasolar planets (Angel (1994)), we are clearly entering a new phase in extrasolar planetary research.
## 2. A Summary of the Data
The transit of the F8V/G0V star HD209458 (at a distance of 47 parsecs) by HD209458b lasts $``$3 hours (out of a total period of 3.524 days) and has a depth of $``$1.5-2.0%. Its ingress and egress phases each last $``$25 minutes (Charbonneau et al. 1999b ; Henry et al. (1999)). The properties of the planet, in particular its orbital distance and its radius, scale with the properties of the star and the most recent study of HD209458 was conducted in the context of these transits by Mazeh et al. (2000). They conclude from $`log\mathrm{g}`$/T$`_{\mathrm{eff}}`$spectral-line fits and $`M_V`$/$`(BV)`$ photometric fits that M= $`1.1\pm 0.1`$ M, R = $`1.2\pm 0.1`$ R, T$`_{\mathrm{eff}}`$$``$6000 K, \[Fe/H\]$``$0.0, $`t=5.5\pm 1.5`$ Gyr, and L$`{}_{}{}^{}2.0`$ L$`_{}`$. Mazeh et al. then derive for the planet: R$`_\mathrm{p}`$= 1.40$`\pm 0.17`$ R$`_\mathrm{J}`$, M<sub>p</sub>= $`0.69\pm 0.05`$ M$`_\mathrm{J}`$($``$M$`{}_{}{}^{2/3}{}_{}{}^{}`$), $`i=86^{}.1\pm 1^{}.6`$, and $`a=0.047`$ A.U. ($``$M$`{}_{}{}^{1/3}{}_{}{}^{}`$). (All symbols have their standard meanings.)
Fuhrmann et al. (1998) provide parameters for the F7V star $`\tau `$ Boötis: M= $`1.42\pm 0.05`$ M, R = $`1.48\pm 0.05`$ R, T$`_{\mathrm{eff}}`$$``$6360 K, \[Fe/H\]$``$+0.27$`\pm `$0.08, $`t=1.0\pm 0.6`$ Gyr, and L$`{}_{}{}^{}3.2`$ L$`_{}`$. Its Hipparcos distance is $``$15.6 parsecs. With the Fuhrmann et al. mass for $`\tau `$ Boötis, Butler et al. (1997) would have obtained for its close-in EGP: M<sub>p</sub>$`\mathrm{sin}(i)`$ = 4.33 M$`_\mathrm{J}`$, $`a0.049`$ A.U., and $`P=3.313`$ days. Charbonneau et al. (1999a) quote an upper limit at $``$4900Å of $`5\times 10^5`$ to the fraction of the star’s light reflected off the planet. Cameron et al. (1999) claim to have detected in the blue-green region of the spectrum a reflected fraction of $`1.9\pm 0.4\times 10^4`$. From the semi-amplitude ($``$74 km s<sup>-1</sup>) of the Doppler shift of this reflected fraction, they derive an orbital inclination ($`i`$) of 29, which would yield a mass for $`\tau `$ Boo b of $``$9 M$`_\mathrm{J}`$. From their reflected fraction, an orbital planetary phase function, and a geometric albedo ($`A_g`$) of 0.55 (similar to that of Jupiter in the visible), Cameron et al. obtain a radius for $`\tau `$ Boo b of 1.6–1.8 R$`_\mathrm{J}`$. If $`A_g`$ were smaller, the inferred radius of the planet would be larger ($`A_g^{1/2}`$).
## 3. The Radii of Close-in EGPs
By whatever processes giant planets are initially assembled, they must start out significantly larger than they end up. In isolation, they would cool inexorably due to radiation from their surfaces and shrink accordingly, just as does a protostar or a brown dwarf. Early on, due to the negative effective specific heat of an object in hydrostatic equilibrium supported by ideal gas pressure, energy loss results in an increase in its central temperature. However, the density increases more quickly and, as a consequence, the specific entropy ($`S`$) monotonically decreases. Since EGPs are almost fully convective (even if under significant stellar irradiation; Guillot et al. (1996)), they are isentropes. This is an essential point. Given an EOS and a planetary mass, an EGP’s core entropy determines its radius (and its surface gravity). A large radius is a consequence of a large entropy.
As Zapolsky and Salpeter (1969) demonstrated for planets made of high-$`Z`$ material, any planet with a cold radius larger than $``$0.5 R$`_\mathrm{J}`$, must be made predominantly of hydrogen. Using the ANEOS equation of state tables (Thompson (1990)), we derive that an “olivine” (rock) or H<sub>2</sub>O (ice) planet with a mass of 0.69 M$`_\mathrm{J}`$has a radius of 0.31 R$`_\mathrm{J}`$or 0.45 R$`_\mathrm{J}`$, respectively. Importantly, these radii are 3–4 times smaller than observed for HD209458b and prove that HD209458b must be a hydrogen-rich gas giant; it cannot be a giant terrestrial planet or an ice giant such as Neptune or Uranus.
The rate with which an EGP shrinks is determined by the opacities in its outer radiative zone and the degree of stellar irradiation. If there were no radiative losses, the EGP would not shrink. In isolation, the energy loss rate and T$`_{\mathrm{eff}}`$are determined in the context of a self-consistent radiative/convective model and T$`_{\mathrm{eff}}`$is a function of only $`S`$ and M<sub>p</sub>. (Note that, for a given metallicity, the surface gravity of the EGP is a function of $`S`$ and M<sub>p</sub> alone.) The atmospheric flux in the outer skin determines the temperature profile at the boundary of the convective core and, hence, the rate of core entropy and radius decrease. According to theory (Saumon et al. (1996); Burrows et al. 1995,1997), isolated EGPs shrink rapidly. A 1-M$`_\mathrm{J}`$EGP in isolation contracts below 2.0 R$`_\mathrm{J}`$in less than $`10^6`$ years. The theory depends upon the atmospheric opacities and metallicity and can reproduce the current Jupiter (Hubbard et al. (1999)), but there still remain major uncertainties concerning grain and cloud formation (Lunine et al. (1989)), rainout (Burrows and Sharp (1999)), gas-phase opacities (Burrows et al. (1997)), and the depth of penetration of the stellar flux (Guillot et al. (2000)). Hence, while the basic theory of EGP evolution is firm, the details are not and significant ambiguities in the variation of R$`_\mathrm{p}`$and T$`_{\mathrm{eff}}`$with age persist.
It is by altering the temperature/pressure profile of the atmosphere of an EGP that stellar irradiation can retard the evolution of the EGP’s core entropy and, hence, R$`_\mathrm{p}`$. Essentially, irradiation flattens the temperature profile at the top of the convective zone, while at the same time moving the radiative/convective boundary inward. The consequent growth of the radiative zone is a central feature of the large-radius phenomenon (Guillot et al. (1996)). Since radiative fluxes are governed by the product of thermal diffusivities and temperature gradients and since the thermal diffusivity decreases with increasing pressure, the flux of energy out of the convective core and the rate of core entropy change are reduced. For close-in EGPs, T$`_{\mathrm{eff}}`$stabilizes early at large values, while the “effective” T$`_{\mathrm{eff}}`$of the core, where most of the heat and mass resides, is drastically lowered. The upshot is a retardation of the contraction of the planet. Large EGP radii are a consequence of such retardation, and not of the expansion of the outer envelope by stellar heating. This is easy to demonstrate by noting that the scale height of HD209458b’s atmosphere, under the assumption that T$`_{\mathrm{eff}}`$is between 1200 K and 1700 K, is only $``$1% of R$`_\mathrm{p}`$. Though the mapping between $`S`$ and R$`_\mathrm{p}`$is unaltered by irradiation, the mapping between $`S`$ and age can change significantly. Guillot et al. (1996) predicted this behavior for 51 Peg b, using a Bond albedo ($`A_B`$) of 0.35, obtaining R$`_\mathrm{p}`$’s of 1.1–1.3 R$`_\mathrm{J}`$after 8 Gyr. Scaling the results of that paper using the stellar flux on HD209458b and a mass of 0.69 M$`_\mathrm{J}`$, we obtain R$`_\mathrm{p}`$’s between 1.4 R$`_\mathrm{J}`$and 1.6 R$`_\mathrm{J}`$for ages between 10 and 3 Gyr, perfectly in line with the transit observations (§2). Figure 1 depicts two possible evolutionary trajectories for HD209458b, if born and fixed at 0.049 A.U. Also included on Figure 1 is a theoretical R$`_\mathrm{p}`$-$`t`$ trajectory for a 0.69-M$`_\mathrm{J}`$EGP in isolation. We have used for these calculations the formalisms of Guillot et al. (1995, 1996), Guillot and Morel (1995), and Burrows et al. (1997). Tidal and Roche effects have been ignored and are relevant only at very early ages ($`<`$$`10^4`$ yrs). Note, however, that despite its small orbital distance, HD209458b is beyond the Roche limit by a factor of $``$2.5 and is stable against loss both by thermal (Jeans) escape and non-thermal processes involving absorption of stellar UV flux (Trilling (1998)). A box of empirical ages and R$`_\mathrm{p}`$’s (Mazeh et al. 2000) is superposed.
As can be seen in Figure 1, there is a great difference between the theoretical radius of an isolated and an irradiated EGP. Importantly, since the scale-height effect is miniscule, the R$`_\mathrm{p}`$-$`t`$ trajectory of the isolated EGP (model I) immediately suggests that if HD209458b were allowed to dwell at large orbital distances ($``$0.5 A.U.) for more than a few $`\times 10^7`$ years, its observed radius could not be reproduced. As Figure 1 demonstrates, it is at such ages that the radius of an isolated 0.69-M$`_\mathrm{J}`$EGP falls below HD209458b’s observed radius. Note that for $``$0.5 A.U. the intrinsic luminosity of our isolated HD209458b model at such an age exceeds the amount of stellar light that would have been intercepted. Hence, stellar irradiation would have had little effect on this model. However, many orbital-distance/age histories can be contemplated and these will be the subject of a future work (Guillot et al. (2000)). To reiterate, HD209458b could not have cooled off and achieved a radius or an entropy like that of Jupiter and then moved in. If it migrated, it had to have moved in early in its life. The radius of a close-in EGP depends upon its history (and the history of the star); large radii require early proximity to a central star. This fact provides an upper limit to the timescale of planetary migration, if migration did indeed occur (Trilling et al. (1998)): conservatively, HD209458b dwelled less than a few$`\times 10^7`$ years at more than 0.5 A.U.
The actual evolution of a close-in EGP’s radius depends upon its Bond albedo (Marley et al. (1999); Seager and Sasselov (1998); Sudarsky, Burrows, and Pinto (1999)), the level of any clouds formed and their optical depths, the gas-phase abundances and opacities, the helium and metallicity fractions, the H-He EOS (Saumon, Chabrier, & Van Horn (1995); Saumon et al. (1999)), the depth of stellar flux penetration, and the primary star. In a later paper (Guillot et al. (2000)), we explore these effects and conduct a detailed parameter study. A range of trajectories similar to those depicted in Figure 1 are still possible and when data on the HD209458b transit and primary star are further refined, theorists may well be able to sharply constrain the character of HD209458b’s atmosphere and composition.
### 3.1. Theoretical Radius for $`\tau `$ Boo b
We provide in Figure 1 two representative theoretical R$`_\mathrm{p}`$-$`t`$ trajectories for $`\tau `$ Boo b, along with an error box constructed using Cameron et al. (1999) and Fuhrmann et al. (1998). Depending upon L and its Bond albedo, which might vary from $``$0.0 to perhaps 0.6, $`\tau `$ Boo b’s T$`_{\mathrm{eff}}`$is between 1350 K and 1750 K. These T<sub>eff</sub>s are higher than those corresponding to HD209458b and reflect $`\tau `$ Boötis’ higher L. Due to uncertainties in the relative position of the silicate and iron cloud decks and the region of neutral alkali metals, there are still uncertainties in the atmospheric composition and albedos of such an EGP (Marley et al. (1999); Sudarsky, Burrows, and Pinto (1999)). A high Bond albedo might be a consequence of the presence at altitude of reflecting clouds. Without these clouds, both the Bond albedo and the geometric albedo in the blue-green region of the planetary spectrum would be low, perhaps below 0.1. A low geometric albedo would put the measured radius of $`\tau `$ Boo b even higher than the range quoted by Cameron et al. . As the evolutionary models on Figure 1 suggest, though a lower Bond albedo results in a larger R$`_\mathrm{p}`$at a given age, the R$`_\mathrm{p}`$-$`t`$ error box for $`\tau `$ Boo b given in Figure 1 still seems out of reach. Note that a 9-M$`_\mathrm{J}`$$`\tau `$ Boo b model in isolation achieves a radius of 1.6 R$`_\mathrm{J}`$within a scant 5 Myr and that at 1 Gyr such a model has a radius of 1.1 R$`_\mathrm{J}`$.
The early radii of $`\tau `$ Boo b shown on Figure 1 are so much smaller than those of HD209458b because of the larger inferred mass of $`\tau `$ Boo b and the particular choice of opacity data used for these exploratory models. Other opacity databases would give quantitatively different radius-age trajectories early on, but would not alter the conclusion that the theoretical radius of a massive planet at late times ($`>`$ few$`\times 10^7`$ years) is significantly below the $`\tau `$ Boo b error box shown in Figure 1. Furthermore, under the assumptions that the Bond albedo is zero, that there is no outer radiative zone, and that the atmosphere and core of the planet have the same entropy, we can obtain a strong upper bound to the radius of a planet of a given age and mass (Guillot et al. 1996). For $`\tau `$ Boo b, this radius is 1.58 R$`_\mathrm{J}`$for 7 M$`_\mathrm{J}`$and 1.48 R$`_\mathrm{J}`$for 10 M$`_\mathrm{J}`$, both below the quoted radius range. Hence, we have difficulty fitting the Cameron et al. (1999) reflection data. In particular, despite significant irradiation, the large planetary mass measured by Cameron et al. (7-10 M$`_\mathrm{J}`$) results in rapid early contraction. If the Bond albedo is lower, the theoretical $`\tau `$ Boo b radius for an age near 1 Gyr would be larger, but the measured radius would also be larger.
## 4. The Effects of the Equation of State
Our cooling calculations use the hydrogen/helium EOS of Saumon, Chabrier, and Van Horn (1995, SCVH). Recent shock-compression experiments on deuterium (Holmes, Ross, and Nellis (1995); Da Silva et al. (1997); Collins et al. (1998)) show that this EOS underestimates the degree of molecular dissociation for pressures in the range $`0.1<P<2`$Mbar. However, by softening the repulsive part of the H<sub>2</sub>–H and H–H potentials, the EOS of SCVH can be made to reproduce all experimental results (Saumon et al. (1999)). Since adiabats of hot Jupiters computed with the modified EOS are systematically cooler and denser than those of SCVH by up to 6% in $`T`$ and 10% in $`\rho `$, such changes in the EOS will reduce the theoretical radius of HD209458b by only a few percent, while modifying the corresponding quantity for $`\tau `$ Boo b by a negligible amount.
## 5. Atmospheric Transmission and Refraction
Details of the transit lightcurve depend upon the distribution of slant optical depth, as determined by molecular opacity and cloud layers, and, to a lesser extent, on refractive redistribution of the stellar surface brightness by the planetary atmosphere.
We use a theory to compute refractive effects that is essentially identical to that of the standard theory for occultations of stars by planetary atmospheres (Hubbard, Yelle, and Lunine (1990)) and have calculated the stellar brightness distribution for the transit of a planet with a radius equal to 1.4 R<sub>J</sub> and a hot solar-composition atmosphere. The slant optical depth $`\tau `$ is computed for molecular opacity sources (with and without clouds) and at a variety of wavelengths. The effective radius of the planet, as determined by fitting the transit lightcurve, will depend upon the gas-phase opacity, the molecular composition, and the location of the optically-thick cloud layers (Seager and Sasselov (1999)). It will also be a diagnostic function of wavelength. Our preliminary calculations for cloud-free models indicate that $`|\mathrm{\Delta }R_\mathrm{p}/R_\mathrm{p}|`$ between 4500 Å and 6500 Å might be as much as 3%, smaller ($``$0.5%) if the alkali metals are important. Using a simple model for silicate cloud growth (Lunine et al. (1989); Marley et al. (1999)), we find a cloud base for HD209458b between $`10^2`$ and 0.4 bars and grain particle sizes between 1 and 100 microns. While for cloud models the wavelength dependence is muted, since the opacity varies strongly with particle size, there exists the possibility of remote sensing of cloud properties with high-precision measurements of upcoming transits (Hubbard et al. (2000)).
## 6. Conclusions
Our theoretical calculations allow us to draw several specific conclusions:
1. HD209458b is a real object, made predominantly of hydrogen.
2. HD209458b’s radius is a consequence of the retardation of contraction by stellar irradiation and is not due to atmospheric expansion by stellar heating.
3. A large radius such as that of HD209458b requires early proximity to its central star.
Curiously, given L, $`a`$, and R$`_\mathrm{p}`$, HD209458b’s total luminosity is $`1.5\times 10^4`$ L$`_{}`$, about twice that of a star with $``$100 times the mass at the very bottom of the stellar main sequence. Given the large inferred mass of $`\tau `$ Boo b, its large radius is less easy to explain theoretically. However, the inherent difficulties of close-in EGP modeling may yet be responsible for theoretical surprises of a qualitative nature.
We thank Tim Brown, David Charbonneau, Geoff Marcy, Michel Mayor, Sara Seager, Richard Freedman and Jim Liebert for many useful discussions and for the use of data in advance of publication. This work was supported in part by NASA grants NAG5-7499, NAG5-7073, NAG5-4214, NAG5-7211, NAG2-6007, NAG5-4988, and NAG5-4987, as well as by an NSF CAREER grant (AST-9624878) to M.S.M.
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# On the Tractable Counting of Theory Models and its Application to Belief Revision and Truth Maintenance
## Introduction
Knowledge compilation has been emerging recently as a new direction of research for dealing with the computational intractability of general propositional reasoning (??). According to this approach, the reasoning process is split into two phases: an off-line compilation phase and an on-line query-answering phase. In the off-line phase, the propositional theory is compiled into some target language, which is typically a tractable one. In the on-line phase, the compiled target is used to efficiently answer a (potentially) exponential number of queries. The main motivation behind knowledge compilation is to push as much of the computational overhead as possible into the off-line phase, in order to amortize that overhead over all on-line queries.
One of the key aspects of any compilation approach is the target language into which the propositional theory is compiled. Previous compilation approaches have proposed Horn theories, prime implicates/implicants, and ordered binary decision diagrams (OBDDs) as targets for such compilation (?????). A more recent compilation target language is decomposable negation normal form (DNNF) (???). DNNF is universal; supports a rich set of polynomial-time operations; is more space–efficient than OBDDs (?); and is very simplistic as far as its structure and algorithms are concerned. Propositional theories in DNNF are highly tractable (?):
1. Deciding the satisfiability of a DNNF can be done in linear time.
2. Conjoining a DNNF with a set of literals can be done in linear time.
3. Projecting a DNNF on some atoms can be done in linear time. Intuitively, to project a theory on a set of atoms is to compute the strongest sentence entailed by the theory on these atoms.
4. Computing the minimum–cardinality of a DNNF can be done in linear time. The cardinality of a model is the number of atoms that are set to false in the model. The minimum–cardinality of a theory is the minimum–cardinality of any of its models.
5. Minimizing a DNNF can be done in quadratic time. To minimize a theory is to produce another theory which models are exactly the minimum-cardinality models of the original theory.
6. Enumerating the models of a DNNF can be done in time linear in its size and quadratic in the number of its models.
This paper rests on two key contributions. First, we show that DNNF representations scale as well as OBDD representations (?) by presenting a linear-time algorithm for converting an OBDD representation of a propositional theory into an equivalent DNNF representation. Second, we identify a subclass of DNNF, which we call deterministic DNNF, d-DNNF, and present a new linear-time operation for counting its models under the assertion, retraction and flipping of literals. In particular, we show how to traverse a d-DNNF only twice and yet compute: its number of models under the assertion, retraction and flipping of each literal. This allows us to test in linear time: the consistency of the d-DNNF under the assertion, retraction and flipping of each literal, therefore, allowing us to implement linear-time, complete truth maintenance and belief revision systems.
What is interesting though is that two of our key complexity results with respect to DNNF compilations continue to hold with respect to the subclass of d-DNNF. This includes a structure-based algorithm which can compile any CNF into an equivalent DNNF in linear time given that the treewidth of the CNF is bounded (?). It also includes the newly proposed algorithm for converting any OBDD into an equivalent DNNF in linear time.
This paper is structured as follows. We first review DNNF, introduce the class of deterministic DNNF and then discuss the new operation for model counting. We follow that by discussing the application of this counting operation to truth maintenance, belief revision, and model-based diagnosis systems. We finally close with some concluding remarks. Proofs of all results are available in the full paper (?).
## Deterministic DNNF
A propositional sentence is in negation normal form (NNF) if it is constructed from literals using only the conjoin and disjoin operators. Figure 1 shows a sentence in NNF depicted as a rooted, directed acyclic graph where the children of each node are shown below it in the graph. Each leaf node represents a literal and each non-leaf node represents a conjunction or a disjunction. We allow $`\mathrm{𝗍𝗋𝗎𝖾}`$ and $`\neg \mathrm{𝖿𝖺𝗅𝗌𝖾}`$ to appear as leaves in a DNNF to denote a conjunction with no conjuncts. Similarly, we allow $`\mathrm{𝖿𝖺𝗅𝗌𝖾}`$ and $`\neg \mathrm{𝗍𝗋𝗎𝖾}`$ as leaves to represent a disjunction with no disjuncts. The size of an NNF is measured by the number of edges in its graphical representation. Our concern here is mainly with a subclass of NNFs:
###### Definition 1
(?) A decomposable negation normal form (DNNF) is a negation normal form satisfying decomposability property: for any conjunction $`_i\alpha _i`$ appearing in the form, no atom is shared by any pair of conjuncts in $`_i\alpha _i`$.
The NNF $`(AB)(\neg AC)`$ is not decomposable since atom $`A`$ is shared by the two conjuncts. But the NNF in Figure 1 is decomposable. It has ten conjunctions and the conjuncts of each share no atoms. Decomposability is the property which makes DNNF tractable.
One of the key operations on DNNF is that of conditioning:
###### Definition 2
(?) Let $`\alpha `$ be a propositional sentence and let $`\gamma `$ be an instantiation.<sup>1</sup><sup>1</sup>1An instantiation of a set of atoms is a conjunction of literals, one literal for each atom in the set. The conditioning of $`\alpha `$ on $`\gamma `$, written $`\alpha \gamma `$, is the sentence resulting from replacing each atom $`p`$ in $`\alpha `$ with $`\mathrm{𝗍𝗋𝗎𝖾}`$ if the +ve literal $`p`$ appears in $`\gamma `$ and with $`\mathrm{𝖿𝖺𝗅𝗌𝖾}`$ if the -ve literal $`\neg p`$ appears in $`\gamma `$.
For example, conditioning the DNNF $`(\neg A\neg B)(BC)`$ on instantiation $`BD`$ gives $`(\neg A\neg \mathrm{𝗍𝗋𝗎𝖾})(\mathrm{𝗍𝗋𝗎𝖾}C)`$ and conditioning it on $`\neg BD`$ gives $`(\neg A\neg \mathrm{𝖿𝖺𝗅𝗌𝖾})(\mathrm{𝖿𝖺𝗅𝗌𝖾}C)`$. Conditioning is a key operation because it allows us to conjoin a DNNF $`\mathrm{\Delta }`$ with some instantiation $`\alpha `$ (which may share atoms with $`\mathrm{\Delta }`$) while ensuring that the result is also a DNNF. Specifically, $`(\mathrm{\Delta }\alpha )\alpha `$ is a DNNF which is equivalent to $`\mathrm{\Delta }\alpha `$ and can be computed in time linear in the size of $`\mathrm{\Delta }`$.
We now introduce the class of deterministic DNNF:
###### Definition 3
A deterministic DNNF (d-DNNF) is a DNNF satisfying the following property: for any disjunction $`_i\alpha _i`$ appearing in the form, every pair of disjuncts in $`_i\alpha _i`$ are disjoint.
For example, $`(AB)C`$ is a DNNF but is not deterministic since the disjuncts $`AB`$ and $`C`$ are not disjoint. However, the DNNF $`(AB)(\neg AC)`$ is deterministic. Note that although every DNF is a DNNF, not every DNF is a d-DNNF.
The main value of d-DNNF is the ability to count its models in linear time and under the assertion, retraction and flipping of literals. Such operations will be discussed in the following section. In the remainder of this section we present two results that hold for DNNF, but continue to hold for d-DNNF. The first is a structure-based algorithm that we introduced in (?) for converting a CNF into a DNNF. The algorithm utilizes a decomposition tree, which is a binary tree the leaves of which correspond to the CNF clauses—see Figure 2. Each decomposition tree has a width and the complexity of the algorithm is exponential only in the width of used tree. The algorithm is given with a slight modification in Figure 3. The only difference between this version and the one in (?) is that we have $`\mathrm{cl2ddnnf}(\mathrm{𝑐𝑙𝑎𝑢𝑠𝑒}(N)\alpha )`$ instead of $`\mathrm{𝑐𝑙𝑎𝑢𝑠𝑒}(N)\alpha `$, therefore, converting a clause to a d-DNNF at the boundary condition.
###### Theorem 1
(?) Let $`N`$ be the root of decomposition tree $`T`$ used in Figure 3. Then $`\mathrm{cnf2ddnnf}(N,\alpha )`$ will return $`\mathrm{\Delta }\alpha `$ in DNNF, where $`\mathrm{\Delta }`$ contains the clauses attached to the leaves of $`T`$. Moreover, the time and space complexity of the algorithm is $`O(nw2^w)`$, where $`n`$ is the number of clauses in $`\mathrm{\Delta }`$ and $`w`$ is the width of decomposition tree $`T`$.
###### Theorem 2
The DNNF returned by the algorithm in Figure 3 is deterministic.
The class of CNF theories with bounded treewidth is defined in (?), where it is shown that, for this class of theories, one can construct in linear time a decomposition tree of bounded width. Therefore, one can compile a d-DNNF of linear size for this class of theories.
Binary decision diagrams (BDDs) are among the most successful representations of propositional theories (?). Two special classes of BDDs, OBDDs and FBDDs, are especially popular given (??):
the number of linear-time operations they support and
the number of real-world theories that admit efficient OBDD/FBDD representations.
We now present a linear-time algorithm for converting an FBDD into an equivalent d-DNNF, showing that d-DNNFs scale as well as FBDDs (which include OBDDs as a subclass). We start by the formal definitions of BDDs, OBDDs, and FBDDs.
###### Definition 4
A binary decision diagram (BDD) over a set of binary variables $`X=\{x_1,\mathrm{},x_n\}`$ is a directed acyclic graph with one root and at most two leaves labeled $`0`$ and $`1`$. Each non-leaf node $`m`$ is labeled by a variable $`\mathrm{𝑣𝑎𝑟}(m)X`$ and has two outgoing edges labeled $`0`$ and $`1`$, where $`\mathrm{𝑙𝑜𝑤}(m)`$ and $`\mathrm{ℎ𝑖𝑔ℎ}(m)`$ denote the nodes pointed to by these edges, respectively. The computation path for input $`(a_1,\mathrm{},a_n)`$, where $`a_i\{0,1\}`$, is defined as follows. One starts at the root. At inner node $`m`$, where $`\mathrm{𝑣𝑎𝑟}(m)=x_i`$, one moves to node $`\mathrm{𝑙𝑜𝑤}(m)`$ if $`a_i=0`$ and to node $`\mathrm{ℎ𝑖𝑔ℎ}(m)`$ otherwise. The BDD represents the Boolean function $`f`$ if the computation path for each input $`(a_1,\mathrm{},a_n)`$ leads to the leaf node labeled with $`f(a_1,\mathrm{},a_n)`$.
The size of a BDD is measured by the number of nodes it contains.
###### Definition 5
A binary decision diagram is called a free BDD (FBDD) if on each computation path each variable is tested at most once. A free BDD is called an ordered BDD (OBDD) if on each computation path the variables are tested in the same order.
OBDDs are a strict subclass of FBDDs (?) and have received much consideration in the verification literature, where they are used to test the equivalence between the specs of a Boolean function and its circuit implementation. OBDDs/FBDDs permit such a test to be performed in polynomial time. Their popularity stems from the existence of efficient OBDD/FBDD representations of many complex, real-world propositional theories. DNNF is more space-efficient than FBDDs (?), but this should not be surprising as FBDDs admit more linear-time operations than does DNNF. For example, the equivalence of two DNNFs cannot be decided in polynomial time while it can for FBDDs.
Figure 4 depicts a recursive algorithm for converting an FBDD into a DNNF, showing that DNNFs scale as well as FBDDs. The algorithm should be called on the root of given FBDD:
###### Theorem 3
(?) The algorithm of Figure 4 takes time linear in the size of given FBDD and returns a DNNF of the theory represented by the given FBDD.
###### Theorem 4
The DNNF returned by the algorithm of Figure 4 is deterministic.
This has major implications on our reported results regarding truth maintenance and belief revision systems, as it proves, constructively, that we can build efficient truth maintenance and belief revision systems for any propositional theory which has an efficient FBDD representation. Figure 5 depicts an FBDD and its corresponding d-DNNF as generated by the algorithm of Figure 4.
## Counting Models of d-DNNF
We now turn to an operation on d-DNNF which is of major significance to a number of AI applications, including belief revision, truth maintenance and diagnosis. Specifically, given a d-DNNF $`\mathrm{\Delta }`$ and a set of literals $`𝐒`$, we describe two traversal operations each taking linear time. By the end of the first traversal, we will be able to count the models of $`\mathrm{\Delta }𝐒`$. By the end of the second traversal, we will be able to count the models of:
1. $`\mathrm{\Delta }𝐒\{l\}`$ for every literal $`l𝐒`$;
2. $`\mathrm{\Delta }𝐒\{l\}`$ for every literal $`l𝐒`$;
3. $`\mathrm{\Delta }𝐒\{l\}\{\neg l\}`$ for every literal $`l𝐒`$.
That is, once we traverse the d-DNNF twice, we will be able to obtain each of these counts using constant-time, lookup operations. As we shall see in the following section, these counts are all we need to implement an interesting number of AI applications.
The traversal will not take place on the d-DNNF itself, but on a secondary structure that we call the counting graph. Without loss of generality, we will assume from here on that the d-DNNF is smooth:
###### Definition 6
A DNNF is smooth iff
1. every literal and its negation appear in the DNNF;
2. for any disjunction $`_i\alpha _i`$ in the DNNF, we have $`\mathrm{𝑎𝑡𝑜𝑚𝑠}(_i\alpha _i)=\mathrm{𝑎𝑡𝑜𝑚𝑠}(\alpha _i)`$ for every $`\alpha _i`$.
The d-DNNF in Figure 1 is smooth as it satisfies these two conditions. We can easily make a DNNF smooth using two operations:
1. If the negation of literal $`l`$ does not appear in the DNNF, replace the occurrence of $`l`$ with $`l(\neg l\mathrm{𝖿𝖺𝗅𝗌𝖾})`$;
2. For each disjunction $`_i\alpha _i`$, replace the disjunct $`\alpha _i`$ with $`\alpha _i_{A\mathrm{\Sigma }}(A\neg A)`$, where $`\mathrm{\Sigma }`$ are the atoms appearing in $`_i\alpha _i`$ but not in $`\alpha _i`$.
These operations preserve both the decomposability and determinism of a DNNF. They may increase the size of given DNNF but only by a factor of $`O(n)`$, where $`n`$ is the number of atoms in the DNNF. This increase is quite minimal in practice though. Note that the d-DNNFs generated by the algorithm of Figure 4 satisfy the first condition; and the d-DNNFs generated by the algorithm of Figure 3 satisfy the second condition as long as $`\mathrm{cl2ddnnf}(\mathrm{𝑐𝑙𝑎𝑢𝑠𝑒}(N)\alpha )`$ satisfies some simple conditions; see (?).
The counting graph of a d-DNNF is a function of many variables represented as a rooted DAG.
###### Definition 7
The counting graph of a smooth d-DNNF is a labeled, rooted DAG. It contains a node labeled with $`l`$ for each literal $`l`$, a node labeled with $`+`$ for each or-node, and a node labeled with $``$ for each and-node in the d-DNNF. There is an edge between two nodes in the counting graph iff there is an edge between their corresponding nodes in the d-DNNF.
Figure 6 depicts the counting graph of the d-DNNF in Figure 1. The size of a counting graph is therefore equal to the size of its corresponding d-DNNF. We will see now how such a graph can be used to perform the counting operations we are interested in.
###### Definition 8
The value of a node $`N`$ in a counting graph under a set of literals $`𝐒`$ is defined as follows:
* $`\mathrm{val}(N)=0`$ if $`N`$ is labeled with literal $`l`$ and $`\neg l𝐒`$;
* $`\mathrm{val}(N)=1`$ if $`N`$ is labeled with literal $`l`$ and $`\neg l𝐒`$;
* $`\mathrm{val}(N)=_i\mathrm{val}(N_i)`$ if $`N`$ is labeled with $``$, where $`N_i`$ are the children of $`N`$;
* $`\mathrm{val}(N)=_i\mathrm{val}(N_i)`$ if $`N`$ is labeled with $`+`$, where $`N_i`$ are the children of $`N`$.
The value of a counting graph $`G`$ under literals $`𝐒`$, written $`G(𝐒)`$, is the value of its root under $`𝐒`$.
Here’s our first counting result.
###### Theorem 5
Let $`\mathrm{\Delta }`$ be a smooth d-DNNF, $`𝐒`$ be a set of literals, and let $`G`$ be the counting graph of $`\mathrm{\Delta }`$. The value of $`G`$ under $`𝐒`$ is the number of models of $`\mathrm{\Delta }𝐒`$:
$$G(𝐒)=\mathrm{𝑀𝑜𝑑𝑒𝑙𝑠}\mathrm{\#}(\mathrm{\Delta }𝐒).$$
Note that $`G(𝐒)>0`$ iff $`\mathrm{\Delta }𝐒`$ is consistent. Therefore, by traversing the counting graph once we can test the consistency of d-DNNF $`\mathrm{\Delta }`$ conjoined with any set of literals $`𝐒`$. Figure 6 depicts the counting graph of d-DNNF $`\mathrm{\Delta }`$ in Figure 1, evaluated under the literals $`𝐒=A,\neg B`$. This indicates that $`\mathrm{\Delta }\{A,\neg B\}`$ has two models.
We now present the central result in this paper. First, we note that when viewing a counting graph $`G`$ as a function of many variables, we will use $`V_l`$ to denote the variable (node) which corresponds to literal $`l`$. Second, we can talk about the partial derivative of $`G(𝐒)`$ with respect to any of its variables $`V_l`$, $`G(𝐒)/V_l`$. Due to the decomposability of d-DNNF, the function $`G(𝐒)`$ is linear in each of its variables. Therefore, the change to the count $`G(𝐒)`$ as a result of adding, removing or flipping a literal in $`𝐒`$ can be obtained from the partial derivatives, without having to re-evaluate the counting graph $`G`$. This leads to the following consequential result:
###### Theorem 6
Let $`\mathrm{\Delta }`$ be a smooth d-DNNF, $`𝐒`$ be a set of literals, and let $`G`$ be the counting graph of $`\mathrm{\Delta }`$. We have:
When $`l,\neg l𝐒`$:
$$\mathrm{𝑀𝑜𝑑𝑒𝑙𝑠}\mathrm{\#}(\mathrm{\Delta }𝐒\{l\})=G(𝐒)/V_l.$$
When $`l𝐒`$:
$$\mathrm{𝑀𝑜𝑑𝑒𝑙𝑠}\mathrm{\#}(\mathrm{\Delta }𝐒\{l\})=G(𝐒)/V_l+G(𝐒)/V_{\neg l}.$$
When $`l𝐒`$:
$`\mathrm{𝑀𝑜𝑑𝑒𝑙𝑠}\mathrm{\#}(\mathrm{\Delta }𝐒\{l\}\{\neg l\})`$
$`=`$ $`G(𝐒)G(𝐒)/V_l+G(𝐒)/V_{\neg l}.`$
Therefore, if we can compute the partial derivative of $`G(𝐒)`$ with respect to each of its variables, then we can count the models of $`\mathrm{\Delta }𝐒`$ under the assertion of new literals not in $`𝐒`$, and under the retraction or flipping of literals in $`𝐒`$. Figure 7 depicts the value of each of these partial derivatives for the d-DNNF in Figure 1. The counting graph is evaluated under literals $`𝐒=A,\neg B,C`$ and the partial derivatives are shown below each variable. According to these derivatives, we have:
$`\mathrm{𝑀𝑜𝑑𝑒𝑙𝑠}\mathrm{\#}(\mathrm{\Delta }𝐒\{D\})=1`$ and $`\mathrm{𝑀𝑜𝑑𝑒𝑙𝑠}\mathrm{\#}(\mathrm{\Delta }𝐒\{\neg D\})=0`$. This immediately tells us that $`\mathrm{\Delta }𝐒D`$.
$`\mathrm{𝑀𝑜𝑑𝑒𝑙𝑠}\mathrm{\#}(\mathrm{\Delta }𝐒\{A\})=1+1=2`$; $`\mathrm{𝑀𝑜𝑑𝑒𝑙𝑠}\mathrm{\#}(\mathrm{\Delta }𝐒\{\neg B\})=1+1=2`$; and $`\mathrm{𝑀𝑜𝑑𝑒𝑙𝑠}\mathrm{\#}(\mathrm{\Delta }𝐒\{C\})=1+1=2`$. Therefore, retracting any literal in $`𝐒`$ increases the number of models to $`2`$.
$`\mathrm{𝑀𝑜𝑑𝑒𝑙𝑠}\mathrm{\#}(\mathrm{\Delta }𝐒\{A\}\{\neg A\})=11+1=1`$; $`\mathrm{𝑀𝑜𝑑𝑒𝑙𝑠}\mathrm{\#}(\mathrm{\Delta }𝐒\{\neg B\}\{B\})=11+1=1`$; $`\mathrm{𝑀𝑜𝑑𝑒𝑙𝑠}\mathrm{\#}(\mathrm{\Delta }𝐒\{C\}\{\neg C\})=11+1=1`$. Therefore, flipping any literal in $`𝐒`$ will not change the number of models (although it does change the model itself).
There is one missing link now: How do we compute the partial derivatives of a counting graph with respect to each of its variables? This actually turns out to be easy due to results in (??) which show how to evaluate and simultaneously compute all partial derivatives of a function by simply traversing its computation graph in linear time. Although (?) casts such computation in terms of summing weights of paths in such a graph, we present a more direct implementation here. In particular, if we let $`\mathrm{pd}(N)`$ denote the partial derivative of $`G(𝐒)`$ with respect to a node $`N`$ in the counting graph, then $`\mathrm{pd}(N)`$ is the summation of contributions made by parents $`M`$ of $`N`$:
$`\mathrm{pd}(N)`$ $`=`$ $`\{\begin{array}{cc}1,\hfill & \text{N is the root node;}\hfill \\ {\displaystyle \underset{M}{}}\mathrm{cpd}(M,N),\hfill & \text{otherwise;}\hfill \end{array}`$
where the contribution of parent $`M`$ to its child $`N`$ is computed as follows:
$$\mathrm{cpd}(M,N)=\{\begin{array}{cc}\mathrm{pd}(M),\hfill & M\text{ is a }+\text{node;}\hfill \\ \mathrm{pd}(M)\underset{KN}{}\mathrm{val}(K),\hfill & M\text{ is a }\text{node;}\hfill \end{array}$$
where $`K`$ is a child of $`M`$. This computation can be performed by first traversing the counting graph once to evaluate it, assigning $`\mathrm{val}`$ to each node $`N`$, and then traversing it a second time, assigning $`\mathrm{pd}`$ for each node $`N`$. We are then mainly interested in $`\mathrm{val}(N)`$ where $`N`$ is the root node, and $`\mathrm{pd}(N)`$ where $`N`$ is a leaf node.
Therefore, both the value of a counting graph under some literals $`𝐒`$ and the values of each of its partial derivatives under $`𝐒`$ can be computed by traversing the graph twice. Once to compute the values, and another to compute the partial derivatives. Note that such traversal needs to be redone once the set of literals $`𝐒`$ changes.
We close this section by pointing out that partial differentiation turns out to play a key role in probabilistic reasoning as well. Specifically, we present a comprehensive framework for probabilistic reasoning in (?) based on compiling a Bayesian network into a polynomial and then reducing a large number of probabilistic queries into the computation of partial derivatives of the compiled polynomial.
## Complete, Linear-Time Truth Maintenance
We now turn to some applications of the results in the previous section. That is, what can we do once we are able to count models under the conditions stated above?
We first consider truth maintenance systems and show how our results allow us to implement complete truth maintenance systems which take linear time on two important classes of propositional theories: those with bounded treewidth, and those admitting a linear FBDD representation. For each class of such theories, we can compile a smooth d-DNNF $`\mathrm{\Delta }`$ in linear time and then use it for truth maintenance as follows.<sup>2</sup><sup>2</sup>2We are assuming that smoothing a d-DNNF does not increase its size by more than a constant factor.
A truth maintenance system takes a set of clauses $`\mathrm{\Gamma }`$ and a set of literals $`𝐒`$ and tries to determine for each literal $`l`$ whether $`\mathrm{\Gamma }𝐒l`$. The most common truth maintenance system is the one based on closing $`\mathrm{\Gamma }𝐒`$ under unit resolution (?). Such a system takes linear time, but is incomplete. Given that the set of literals in $`𝐒`$ changes to $`𝐒^{}`$, the goal of a truth maintenance system is then to update the truth of each literal under the new “context” $`𝐒^{}`$. Sometimes, clauses in $`\mathrm{\Gamma }`$ are retracted and/or asserted. A truth maintenance system is expected to update the truth of literals under such clausal changes too.
Our model-counting results allow us to implement a complete truth maintenance system as follows. We compile the theory $`\mathrm{\Gamma }`$ into a smooth d-DNNF $`\mathrm{\Delta }`$ and construct the counting graph $`G`$ of $`\mathrm{\Delta }`$. Given any set of literal $`𝐒`$, we evaluate $`G`$ under $`𝐒`$ and compute its partial derivatives also under $`𝐒`$. This can be done in time linear in the size of $`\mathrm{\Delta }`$. We are now ready to answer all queries expected from a truth maintenance system by simple, constant-time, look-up operations:
Literal $`l`$ is entailed by $`\mathrm{\Delta }𝐒`$ iff $`\mathrm{\Delta }𝐒\{\neg l\}`$ has no models: $`G(𝐒)/V_{\neg l}=0`$.
Retracting literal $`l`$ from $`𝐒`$ will render $`\mathrm{\Delta }𝐒`$ consistent iff $`\mathrm{\Delta }𝐒\{l\}`$ has at least one model: $`G(𝐒)/V_l+G(𝐒)/V_{\neg l}>0`$.
Flipping literal $`l`$ in $`𝐒`$ will render $`\mathrm{\Delta }𝐒`$ consistent iff $`\mathrm{\Delta }𝐒\{l\}\{\neg l\}`$ has at least one model: $`G(𝐒)G(𝐒)/V_l+G(𝐒)/V_{\neg l}>0`$. <sup>3</sup><sup>3</sup>3Note that the flipping of literals is outside the scope of classical truth maintenance systems in the sense that they must retract $`l`$ and then assert $`\neg l`$, taking linear time, to perform the above operation.
If we want to reason about the assertion/retraction of clauses in theory $`\mathrm{\Gamma }`$, we can replace each clause $`\alpha `$ in $`\mathrm{\Gamma }`$ by $`C_\alpha \alpha `$, where $`C_\alpha `$ is a new atom that represents the truth of clause $`\alpha `$. We then compile the extended theory $`\mathrm{\Gamma }`$ into d-DNNF $`\mathrm{\Delta }`$. To assert all clauses initially, we have to include all atoms $`C_\alpha `$ in the set of literals $`𝐒`$. The assertion/retraction of clauses can then be emulated by the assertion/retraction of atoms $`C_\alpha `$. For example, in case of a contradiction, we can ask whether removing a clause $`\alpha `$ will resolve the contradiction by asking whether $`\mathrm{\Delta }𝐒\{C_\alpha \}`$ has more than one model:
$$G(𝐒)/V_{C_\alpha }+G(𝐒)/V_{\neg C_\alpha }>0.$$
.
## Complete, Linear-Time Belief Revision
We now turn to a second major application of model counting on d-DNNF: the implementation of a very common class of belief revision systems, which is adopted in model-based diagnosis and in certain forms of default reasoning. The problem here is as follows. We have a set of special atoms $`𝐃=\{d_1,\mathrm{},d_n\}`$ in the theory $`\mathrm{\Gamma }`$ which represent defaults. Typically, we assume that all of these defaults are true, allowing us to draw some default conclusions. We then receive some evidence $`𝐒`$ (a set of literals) which is inconsistent with $`\mathrm{\Gamma }𝐃`$. We therefore know that not all defaults are true and some must be retracted—that is, some $`d_i`$s will have to be replaced by $`\neg d_i`$ in $`𝐃`$. Our goal then is to identify a set of literals $`𝐃^{}`$ such that
1. $`d_i𝐃^{}`$ or $`\neg d_i𝐃^{}`$ for all $`i`$;
2. $`\mathrm{\Gamma }𝐃^{}𝐒`$ is consistent;
3. the number of negative literals in $`𝐃^{}`$ is minimized;
and then then report the truth of every literal under the new set of defaults $`𝐃^{}`$. Note that there may be more than one set of defaults $`𝐃^{}`$ that satisfies the above properties. In such a case, a literal holds after the revision process only if it holds under $`\mathrm{\Gamma }𝐃^{}𝐒`$ for every $`𝐃^{}`$.
How can we do this? As we shall see now, if $`\mathrm{\Gamma }`$ is a smooth d-DNNF, then all of this can be done in time linear in the size of $`\mathrm{\Gamma }`$!
This works exactly as in the previous section, except that we have to minimize the d-DNNF first, a process which eliminates some of the d-DNNF models (?). To define this minimization process more precisely, we need the following definition first:
###### Definition 9
If $`\mathrm{\Sigma }`$ is a set of atoms, then the $`\mathrm{\Sigma }`$-cardinality of a model is the number of atoms in $`\mathrm{\Sigma }`$ that the model sets to false.
To $`\mathrm{\Sigma }`$-minimize a theory $`\mathrm{\Gamma }`$ is to convert it into another theory whose models are exactly the models of $`\mathrm{\Gamma }`$ having a minimum $`\mathrm{\Sigma }`$-cardinality. Consider the d-DNNF $`\mathrm{\Gamma }`$ in Figure 1 for an example and suppose that $`\mathrm{\Sigma }=\{A,B,C,D\}`$; that is, we want to minimize the d-DNNF with respect to each of its atoms. This theory has eight models, each having an odd cardinality (one or three). If we $`\mathrm{\Sigma }`$-minimize this d-DNNF, we obtain another with four models, shown in Figure 9.
Given these definitions, we can re-phrase the problem of belief revision (stated above) as follows. Let $`\mathrm{\Sigma }`$ be a set of atoms representing defaults, and let $`\mathrm{\Gamma }`$ be a smooth d-DNNF. Given observation $`𝐒`$, $`\mathrm{\Sigma }`$-minimize the theory $`\mathrm{\Gamma }𝐒`$ to yield $`\mathrm{\Delta }`$ and then report on the truth of each literal under the minimized theory $`\mathrm{\Delta }`$.
As it turns out, one can minimize a smooth d-DNNF in linear time, to yield another smooth d-DNNF to which we can apply the techniques of the previous section. We now describe the process of minimizing a DNNF which is described in more details in (?). We do this in a two-step process:
1. We assign a cardinality to every node in the d-DNNF as follows:
1. each literal whose atom is not in $`\mathrm{\Sigma }`$ gets cardinality zero;
2. each positive literal whose atom is in $`\mathrm{\Sigma }`$ gets cardinality zero;
3. each negative literals whose atom is in $`\mathrm{\Sigma }`$ gets cardinality one;
4. the cardinality of a disjunction is the min of its disjuncts’ cardinalities;
5. the cardinality of a conjunction is the summation of its conjuncts’ cardinalities.
2. For each or-node $`N`$ and its child $`M`$, we delete the edge connecting $`N`$ and $`M`$ if the cardinality of $`N`$ is smaller than the cardinality of $`M`$.
Figure 8 depicts the result of assigning cardinalities to the d-DNNF of Figure 1, and Figure 9 depicts the result of deleting some of its edges. This is the minimized d-DNNF and it has four models:
$`\neg ABCD`$;
$`A\neg BCD`$;
$`AB\neg CD`$;
$`ABC\neg D`$.
Once we have minimized the d-DNNF, we can apply the results of the previous section to obtain the answers we want.
As an example, Figure 10 depicts the counting graph of the minimized d-DNNF $`\mathrm{\Delta }`$ in Figure 9, with its value and partial derivatives computed under the observation $`𝐒=\{\neg A\}`$. From these partial derivatives and Theorem 6, we immediately get:
$`\mathrm{𝑀𝑜𝑑𝑒𝑙𝑠}\mathrm{\#}(\mathrm{\Delta }𝐒\{\neg B\})=0`$; $`\mathrm{𝑀𝑜𝑑𝑒𝑙𝑠}\mathrm{\#}(\mathrm{\Delta }𝐒\{\neg C\})=0`$; $`\mathrm{𝑀𝑜𝑑𝑒𝑙𝑠}\mathrm{\#}(\mathrm{\Delta }𝐒\{\neg D\})=0`$. That is, $`B,C`$ and $`D`$ are all entailed by $`\mathrm{\Delta }𝐒`$.
$`\mathrm{𝑀𝑜𝑑𝑒𝑙𝑠}\mathrm{\#}(\mathrm{\Delta }𝐒\{\neg A\})=1+3=4`$. That is, we have four models if we retract $`\neg A`$.
$`\mathrm{𝑀𝑜𝑑𝑒𝑙𝑠}\mathrm{\#}(\mathrm{\Delta }𝐒\{\neg A\}\{A\})=11+3=3`$. That is, we have three models if we flip $`\neg A`$ to $`A`$.
## Predicting the Behavior of Broken Devices
The above results have direct application to model-based diagnosis, where $`\mathrm{\Delta }`$ is the device description, $`𝐒`$ is the device observation and $`𝐃`$ contains the health modes $`\mathrm{𝑜𝑘}_1,\mathrm{},\mathrm{𝑜𝑘}_n`$. Initially, we assume that all device components are working normally, but then find some observation $`𝐒`$ such that $`\mathrm{\Delta }𝐃=\{\mathrm{𝑜𝑘}_1,\mathrm{},\mathrm{𝑜𝑘}_n\}𝐒`$ is inconsistent.
To regain consistency we must postulate that some of the components are not healthy, therefore, flipping some of the $`\mathrm{𝑜𝑘}_i`$s into $`\neg \mathrm{𝑜𝑘}_i`$ in the set $`𝐃`$. Assuming a smallest number of faults, we want to minimize the number of unhealthy components needed to regain consistency. A set $`𝐃^{}`$ such that:
1. $`\mathrm{𝑜𝑘}_i𝐃^{}`$ or $`\neg \mathrm{𝑜𝑘}_i𝐃^{}`$ for all $`i`$;
2. $`\mathrm{\Delta }𝐃^{}𝐒`$ is consistent;
3. the number of negative literals in $`𝐃^{}`$ is minimized;
is called a minimum-cardinality diagnosis and one goal of model-based diagnosis to enumerate such diagnoses (??).
Another practical problem, however, which has received much less attention in model-based diagnosis is the following: Assuming a smallest number of faults, what is the truth value of every literal appearing in the device description $`\mathrm{\Delta }`$? That is, we do not want to know what the minimum-cardinality diagnoses are. Instead — and under the assumption that one of them has materialized — we want to predict the behavior of the given faulty device. But this is exactly the problem we have treated in the previous section.
Therefore, our results allow us to predict the value of each device port (literal $`l`$) in a broken device, assuming that the number of broken components is minimal, in time linear in the size of device description $`\mathrm{\Delta }`$, as long as $`\mathrm{\Delta }`$ is a smooth d-DNNF. We are unaware of any similar complexity result for model-based diagnosis.
## Conclusion
We have identified two classes of propositional theories, those which have a bounded treewidth, and those which have a linear-sized FBDD representation. We have shown that each of these classes of theories can be converted in linear time into a tractable form that we called deterministic DNNF, d-DNNF. We have also defined linear-time, model-counting operations on d-DNNF, allowing us to implement (a) linear-time, complete truth maintenance systems and (b) linear-time, complete belief revision systems for the two identified classes of propositional theories. Our results also have major implications on the practice of model-based diagnosis as they allow us to efficiently predict the behavior of a broken device, assuming a smallest number of broken components.
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# Single chargino production with R-parity lepton number violation in electron-electron and muon-muon collisions
## 1 Introduction
R-parity is a discrete symmetry defined by assigning to every field the number $`R=(1)^{3B+L+2S}`$ ($`B(L)`$ \- baryon (lepton) number, $`S`$ \- spin of the particle) . If it is conserved then baryon and lepton number violating transitions are forbidden. In that case, the theory guarantees both proton stability and lepton universality. However, in supersymmetric extensions of the Standard Model, gauge invariance and renormalizability, the two main principles of any gauge theory, do not assure R-parity conservation. At present, wide phenomenological investigations of R-parity violating processes have been undertaken (for reviews see e.g. ).
Here we will explore the possibility of discovering the lepton number violating process of single chargino production at future lepton-lepton colliders (see Fig.1(I) for an electron-electron collision). To our knowledge this process has not yet been discussed, though lepton number violating charginos pair production in electron-electron collisions (Fig.1(II)) has been considered .
Let us start with electron-electron collisions. The analysis of the muon option is analogous and will be shortly discussed whenever needed. As can be seen from Fig.1(I), the cross section for single chargino production is proportional to $`\lambda ^2g^2`$ where $`\lambda `$ and $`g`$ are couplings involved in the following Lagrangians ($`\lambda _{abc}=\lambda _{acb}`$, $`a,b,c`$ are family indices):
$``$ $`=`$ $`g\overline{\stackrel{~}{\chi }}_i^cV_{i1}(1\gamma _5)e\stackrel{~}{K}_{em}\stackrel{~}{\nu }_m^{}+h.c.,`$ (1)
$`_\mathit{}_p`$ $`=`$ $`\lambda _{abc}\{\stackrel{~}{\nu }_{aL}\overline{l}_{cR}l_{bL}(ab)\}+h.c.`$ (2)
These Lagrangians are written in physical basis. The matrix $`\stackrel{~}{K}_{em}`$ in Eq. (1) comes from the sneutrino mass matrix diagonalization. If R-parity is violated, we have to take into account the mixing between the sneutrinos $`\stackrel{~}{\nu }_e`$, $`\stackrel{~}{\nu }_\mu `$, $`\stackrel{~}{\nu }_\tau `$ and the neutral Higgs bosons $`H_1^0`$, $`H_2^0`$. We shall, however, assume that this mixing is negligible and does not affect the results, at least at the stage of chargino production. In what follows we shall also assume that the exchange of the lightest (electron) sneutrino dominates (which is equivalent to some hierarchy assumption in the sneutrino sector) and neglect the contribution of the heavier $`\stackrel{~}{\nu }`$’s. We therefore set ($`e`$ stands for electron) $`\stackrel{~}{K}_{em}=\delta _{em}`$ in Eq. (1). For more complicated cases where the interplay between sneutrino masses in propagators and appropriate elements of the $`\stackrel{~}{K}`$ matrix matters we refer to .
The second mixing matrix, namely $`V_{i1}`$ in Eq. (1) is connected with the chargino sector and describes the weights of the wino component of the chargino fields . Since this is the only component of the charginos that couples to the electron and the sneutrino (the charged higgsino coupling is neglected in the limit of zero electron mass) we set for simplicity $`V_{i1}=1`$. This is further justified by the analysis (in the parameter region $`|\mu |100`$ GeV, $`M_2100`$ GeV for both small and large $`\mathrm{tan}\beta `$, with $`\mu ,M_2`$ being the higgsino and gaugino $`SU(2)`$ mass parameters, respectively, and $`\beta `$ a ratio of two vacuum expectation values involved in MSSM). In general the results should be multiplied by $`V_{i1}^2`$. Furthermore, with R-parity violation, additional couplings between leptons, gauginos and higgsinos $`(e,\mu ,\tau ,\stackrel{~}{W}^{},\stackrel{~}{H}^{})_L`$ exist, but are known to be smaller than the gauge ones .
## 2 Single chargino production and decays: results
In Fig. 2 we gather the cross sections for single chargino production at future electron-electron colliders with c.m. energies $`\sqrt{s}=500`$ GeV and $`\sqrt{s}=1`$ TeV as functions of the chargino mass for different sneutrino masses<sup>1</sup><sup>1</sup>1Results (see Appendix) assume a $`P_{}^e=100\%`$ electron beam polarization. In reality we can expect that $`P_{}^e=90\%`$ can be achieved. Then the cross sections must be multiplied by a factor $`\frac{1}{4}(1P_{}^{e_1})(1P_{}^{e_2})0.9`$.. For the R-parity violating coupling, we have used the most conservative available upper limit $`\lambda _{112(3)}\lambda =0.05`$ , independently of the $`\stackrel{~}{\nu }_e`$ mass (in the case of muon-muon collisions the $`\lambda _{212(3)}`$ couplings would be involved). For sneutrino masses larger than 100 GeV this limit becomes weaker .
As can be deduced from Fig.2, with a planned annual luminosity of some 50 fb<sup>-1</sup> yr<sup>-1</sup> and with a discovery limit at a level of 10 events per year ($`\sigma =0.2`$ fb), the process is detectable for a wide range of sparticle masses.
With the R-parity violating production process (I) we are already definitely out of the SM physics. It is therefore interesting to investigate the possible detector signals. With R-parity non-conservation, the collider phenomenology is quite different from the MSSM case and depends especially on the nature of the LSP (Lightest Supersymmetric Particle). In the MSSM, the stable LSP must be charge and color neutral for cosmological reasons . With R-parity violation there are no hints about the unstable LSP. It can be among others a sneutrino, gluino or even a chargino .
Here we give an example of nonstandard phenomenology but restrict ourselves to a scenario in which charginos decay uniquely (via sneutrino exchange) to charged leptons. Final leptonic signals with lepton number violation and without missing energy could be detected, an interesting situation from the point of view of nonstandard physics, as there is no SM background (see further discussion). These two conditions (charged leptons without missing energy in the final state) require the chargino to be the second lightest supersymmetric particle (NLSP) with sneutrino the LSP. This situation is schematically summarized in Fig.3. If the chargino were the LSP its lifetime should be long enough so that it would be seen in the detector. In other cases (i.e. when the chargino is neither NLSP nor LSP) the chargino would also have cascade decays to final jet states . Then, the situation would be more complicated but at least we can expect that for kinematical reasons a decay to the R-parity lepton violating LSP sneutrino would still be important and the final signal with four charged leptons could be observed (work in progress).
Let us discuss the scenario with NLSP chargino. First, we should notice that, to get substantial chargino production (e.g. $`e^{}e^{}\mu ^{}\stackrel{~}{\chi }^{}`$ in Fig.2), we are interested in the situation where at least one, let us say $`\lambda _{112}`$ coupling is large. Then the decay of chargino to three charged leptons must be observed in the detector, as it can undergo uniquely through the same large $`\lambda _{112}`$ coupling (the only possible decay channel, see Fig.3).
In Fig.4, we show the final results for the angular distribution of the final positron (Fig.3) for two different energies ($`\sqrt{s}=500(1000)`$ GeV). We have taken $`m_{\stackrel{~}{\chi }}=205`$ GeV and $`m_{\stackrel{~}{\nu }}=200`$ GeV. Results have been obtained using the VEGAS procedure. Four particles in the final state give us an 8 dimensional integration. We have also applied the narrow width approximation where $`\mathrm{\Gamma }_{\stackrel{~}{\chi }}<<m_{\stackrel{~}{\chi }}`$ (see Appendix for details). The solid line describes results based on Eq. (A.21), when interferences between production and decay of charginos with $`\stackrel{~}{\lambda }=\pm 1/2`$ are taken into account. The dashed line describes results with factorization assumed , which means the following replacement in Eq. (A.21) is done
$$\underset{\stackrel{~}{\lambda }}{}|M(;,\stackrel{~}{\lambda })T\left(\stackrel{~}{\lambda }\right)|^2\frac{1}{2}\underset{\stackrel{~}{\lambda }}{}|M(;,\stackrel{~}{\lambda })|^2\underset{\stackrel{~}{\lambda }}{}|T\left(\stackrel{~}{\lambda }\right)|^2$$
(3)
We can see that spin correlations do not change the results substantially ($`2`$ % for considered c.m energies and the chargino mass). It is important that the positron angular distributions are not so strongly peaked in the beam directions, even for $`\sqrt{s}=1`$ TeV collider energy. With assumed cuts $`(|\mathrm{cos}\mathrm{\Theta }_+|0.95)`$ enough events will be detected to investigate the process.
The only SM process, which gives a four charged lepton signal without missing energy is $`e^{}e^{}e^{}e^{}Z`$. With a possible Z boson decay to the lepton-antilepton pair, it does not coincide with the process under investigation ($`e^{}e^{}2\mu ^{}e^{}e^+`$). That means that we do not have to bother about the SM background contamination. However, this cross section is large enough ($`1`$ pb for $`0.5\sqrt{s}2`$ TeV energies) to cover some other scenarios. As an example let us assume that not only $`\lambda _{112}`$ but also $`\lambda _{121}`$ is not negligible and change the second coupling in the chargino decay channel (Fig.3) from $`\lambda _{112}`$ to $`\lambda _{121}`$. That means that we have now $`(e^{}e^{})e^{}e^{}\mu ^{}\mu ^+`$ in the final state, and this scenario will be dominated by the SM process given above with the Z decay to the muon antimuon pair. In this way we can find that $`\mu ^{}\mu ^{}e^{}e^+(\tau ^+),\mu ^{}\tau ^{}e^{}e^+(\mu ^+)`$ and $`\tau ^{}\tau ^{}e^{}e^+`$ charged lepton signals are testable (meaning sensitivity to the products $`\lambda _{112}\lambda _{112}(\lambda _{113})`$, $`\lambda _{112}\lambda _{113}(\lambda _{123})`$ and $`\lambda _{113}\lambda _{113}`$, respectively).
Finally, our results can be easily applied to the muon-muon collider where another set of R-parity lepton violating couplings can be tested, namely: $`\lambda _{221}\lambda _{212}`$ ($`e^{}e^{}\mu ^{}\mu ^+`$ in the final state), $`\lambda _{221}\lambda _{213}`$ ($`e^{}e^{}\mu ^{}\tau ^+`$) and $`\lambda _{221}\lambda _{231}`$ ($`e^{}\mu ^{}e^+\tau ^{}`$).
## 3 Conclusions
Present experimental limits on R-parity violating couplings do not exclude large and detectable lepton number violating signals in lepton-lepton collisions. We discuss such a possibility in conjunction with single chargino production and its subsequent leptonic decay. If at least one $`\lambda `$ value is large enough - in our discussion mainly $`\lambda _{112}`$ \- single chargino production in electron-electron collisions will be observable (Fig.2). If the chargino is NLSP and sneutrino the LSP, a unique lepton number violating signal of four charged leptons without missing energy could be observed.
## Acknowledgments
We would like to thank H. Fraas, C. Blöchinger, S. Hesselbach (Würzburg University) and G. Moortgat–Pick (DESY, Hamburg) for helpful discussions and valuable remarks. This work was supported by the Polish Committee for Scientific Research under Grant No. 2P03B05418. J.G. would like to thank the Alexander von Humboldt-Stiftung for fellowship.
## Appendix A Appendix
The helicity amplitude for single chargino production $`e^{}e^{}\mu ^{}(\tau ^{})\stackrel{~}{\chi }^{}`$ is given by:
$`M(\sigma _1,\sigma _2;\lambda ,\stackrel{~}{\lambda })`$ $`=`$ $`g\lambda _{112}[\overline{u}(p,\lambda )(1\gamma _5)u(k_1,\sigma _1){\displaystyle \frac{1}{tm_{\stackrel{~}{\nu }_e}^2}}\overline{u}(\stackrel{~}{p},\stackrel{~}{\lambda })(1\gamma _5)u(k_2,\sigma _2)`$ (4)
$``$ $`\overline{u}(\stackrel{~}{p},\stackrel{~}{\lambda })(1\gamma _5)u(k_1,\sigma _1){\displaystyle \frac{1}{um_{\stackrel{~}{\nu }_e}^2}}\overline{u}(p,\lambda )(1\gamma _5)u(k_2,\sigma _2)],`$
$`t(u)`$ $`=`$ $`m_{\stackrel{~}{\chi }}^2\sqrt{s}(\stackrel{~}{E}2\stackrel{~}{p}\mathrm{cos}\stackrel{~}{\mathrm{\Theta }})`$ (5)
where $`(\stackrel{~}{p},\stackrel{~}{\lambda })`$ denotes the momentum and helicity of the chargino, $`(k_{1(2)},\sigma _{1(2)})`$ are the corresponding quantities for the incoming electrons, $`p`$ and $`\lambda `$ denote momentum and helicity of the muon (tau) and $`(\stackrel{~}{\mathrm{\Theta }},\stackrel{~}{\varphi })`$ label the c.m. azimuthal and polar angles of a chargino with respect to the direction of the initial electron $`e_1^{}`$.
To work out the helicity amplitudes, we use the Weyl– van der Waarden spinor formalism in which the 4–spinors can be written:
$`u(p,\lambda )`$ $`=`$ $`\left(\begin{array}{c}\sqrt{Ep\lambda }\chi (p,\lambda )\\ \sqrt{E+p\lambda }\chi (p,\lambda )\end{array}\right),v=\left(\begin{array}{c}\lambda \sqrt{E+p\lambda }\chi (p,\lambda )\\ \lambda \sqrt{Ep\lambda }\chi (p,\lambda )\end{array}\right),`$ (6)
and the Weyl spinors are given by:
$`\chi (p,+1/2)`$ $`=`$ $`\left(\begin{array}{c}e^{i\varphi /2}\mathrm{cos}\theta /2\\ e^{i\varphi /2}\mathrm{sin}\theta /2\end{array}\right),\chi (p,1/2)=\left(\begin{array}{c}e^{i\varphi /2}\mathrm{sin}\theta /2\\ e^{i\varphi /2}\mathrm{cos}\theta /2\end{array}\right),`$ (7)
where $`\mathrm{\Theta },\varphi `$ denote the azimuthal and the polar angle of a particle with respect to the $`\widehat{z}`$ axis.
In the limit of zero mass of all charged leptons we have only two non-vanishing helicity amplitudes ($`\sigma _{1,2},\lambda =1/2,\stackrel{~}{\lambda }=\pm 1/2`$), namely:
$`M(,,)`$ $`=`$ $`g\sqrt{2s}\sqrt{(\stackrel{~}{E}\stackrel{~}{p})(\stackrel{~}{E}\stackrel{~}{p})}\lambda _{112}\left[{\displaystyle \frac{\mathrm{cos}^2\frac{\stackrel{~}{\mathrm{\Theta }}}{2}}{tm_{\stackrel{~}{\nu }_e}^2}}+{\displaystyle \frac{\mathrm{sin}^2\frac{\stackrel{~}{\mathrm{\Theta }}}{2}}{um_{\stackrel{~}{\nu }_e}^2}}\right]`$ (8)
$`M(,,+)`$ $`=`$ $`g\sqrt{2s}\sqrt{(\stackrel{~}{E}+\stackrel{~}{p})(\stackrel{~}{E}\stackrel{~}{p})}\mathrm{cos}{\displaystyle \frac{\stackrel{~}{\mathrm{\Theta }}}{2}}\mathrm{sin}{\displaystyle \frac{\stackrel{~}{\mathrm{\Theta }}}{2}}\lambda _{112}\left[{\displaystyle \frac{1}{tm_{\stackrel{~}{\nu }_e}^2}}{\displaystyle \frac{1}{um_{\stackrel{~}{\nu }_e}^2}}\right].`$
The cross section is:
$$d\sigma =\frac{1}{2s}dLips(s,p,\stackrel{~}{p})\underset{\stackrel{~}{\lambda }}{}\left|M(,;,\stackrel{~}{\lambda }_j)\right|^2$$
(10)
where
$`dLips(s,p,\stackrel{~}{p})`$ $`=`$ $`{\displaystyle \frac{\stackrel{~}{p}}{16\pi ^2p}}d\mathrm{cos}\stackrel{~}{\mathrm{\Theta }}d\stackrel{~}{\varphi }.`$ (11)
In Fig. 3 we study the R-parity violating chargino decays via sneutrino exchange in the t-channel: $`\stackrel{~}{\chi }e^{}e^+\mu ^{}`$. Analogously to the production the amplitude is:
$`T(\stackrel{~}{\lambda }_i)`$ $`=`$ $`g\lambda _{112}\left[\overline{u}(p_1,\sigma _1)(1\gamma _5)u(\stackrel{~}{p},\stackrel{~}{\lambda }){\displaystyle \frac{1}{tm_{\stackrel{~}{\nu }_e}^2+i\mathrm{\Gamma }_{\stackrel{~}{\nu }}m_{\stackrel{~}{\nu }_e}}}\overline{u}(p_2,\sigma _2)(1\gamma _5)v(p_+,\sigma _+)\right],`$
where ($`(p_{1(2)}^i,\sigma _{1(2)}^i)`$ denote the momenta and helicities of the electron and muon, $`(p_+,\sigma _+)`$ are analogous quantities for the final positron.
Using Eqs. (6,7) we get
$`T(\stackrel{~}{\lambda }=+1/2)`$ $`=`$ $`\mathrm{\Omega }_t(+)\left[e^{i/2(\varphi _1\stackrel{~}{\varphi })}\mathrm{sin}{\displaystyle \frac{\mathrm{\Theta }_1}{2}}\mathrm{cos}{\displaystyle \frac{\stackrel{~}{\mathrm{\Theta }}}{2}}+e^{i/2(\varphi _1\stackrel{~}{\varphi })}\mathrm{cos}{\displaystyle \frac{\mathrm{\Theta }_1}{2}}\mathrm{sin}{\displaystyle \frac{\stackrel{~}{\mathrm{\Theta }}}{2}}\right]`$ (13)
$`\times `$ $`\left[e^{i/2(\varphi _2\varphi _+)}\mathrm{sin}{\displaystyle \frac{\mathrm{\Theta }_2}{2}}\mathrm{cos}{\displaystyle \frac{\mathrm{\Theta }_+}{2}}+e^{i/2(\varphi _2\varphi _+)}\mathrm{cos}{\displaystyle \frac{\mathrm{\Theta }_2}{2}}\mathrm{sin}{\displaystyle \frac{\mathrm{\Theta }_+}{2}}\right]`$
$`T(\stackrel{~}{\lambda }=1/2)`$ $`=`$ $`\mathrm{\Omega }_t()\left[e^{i/2(\varphi _1\stackrel{~}{\varphi })}\mathrm{sin}{\displaystyle \frac{\mathrm{\Theta }_1}{2}}\mathrm{sin}{\displaystyle \frac{\stackrel{~}{\mathrm{\Theta }}}{2}}+e^{i/2(\varphi _1\stackrel{~}{\varphi })}\mathrm{cos}{\displaystyle \frac{\mathrm{\Theta }_1}{2}}\mathrm{cos}{\displaystyle \frac{\stackrel{~}{\mathrm{\Theta }}}{2}}\right]`$ (14)
$`\times `$ $`\left[e^{i/2(\varphi _2\varphi _+)}\mathrm{sin}{\displaystyle \frac{\mathrm{\Theta }_2}{2}}\mathrm{cos}{\displaystyle \frac{\mathrm{\Theta }_+}{2}}+e^{i/2(\varphi _2\varphi _+)}\mathrm{cos}{\displaystyle \frac{\mathrm{\Theta }_2}{2}}\mathrm{sin}{\displaystyle \frac{\mathrm{\Theta }_+}{2}}\right]`$
where $`(\mathrm{\Theta }_{+(1,2)},\varphi _{+(1,2)})`$ denote azimuthal and polar angles of the final positron (electron (1), muon (2)) which are defined with respect to the direction of the initial electron beam $`e_1`$ and $`\mathrm{\Omega }_t(\pm )`$ is given by
$$\mathrm{\Omega }_t(\pm )=g\sqrt{8E_+E_1E_2}\sqrt{\stackrel{~}{E}\pm \stackrel{~}{p}\stackrel{~}{\lambda }}\underset{m_{\stackrel{~}{\nu }_n}}{}\lambda _{n12}\frac{1}{tm_{\stackrel{~}{\nu }_n}^2+i\mathrm{\Gamma }_{\stackrel{~}{\nu }}m_{\stackrel{~}{\nu }_n}}.$$
(15)
The decay width can be written as
$$d\mathrm{\Gamma }=\frac{1}{2m_{\stackrel{~}{\chi }}}dLips(\stackrel{~}{m},p_1,p_2,p_+)\underset{\stackrel{~}{\lambda }}{}|T(\stackrel{~}{\lambda })|^2$$
(16)
where
$`dLips(\stackrel{~}{m},p_1,p_2,p_+)`$ $`=`$ $`{\displaystyle \frac{1}{\left(2\pi \right)^58}}{\displaystyle 𝑑E_+d\mathrm{cos}\mathrm{\Theta }_+𝑑\varphi _+d\mathrm{cos}\mathrm{\Theta }_2𝑑\varphi _2}.`$ (17)
Formulae Eq. (A.13) describe the 3–body–decay of a chargino with energy $`\stackrel{~}{E}`$ and angles $`\stackrel{~}{\mathrm{\Theta }},\stackrel{~}{\varphi }`$. The angles of the chargino and of the particles produced by chargino decay are defined with respect to the direction of the initial electron beam $`e_1^{}`$. Angles of the decaying particles are also defined with respect to the initial CM system of colliding electrons. We have left quantities connected with $`e^+`$ as independent parameters to be integrated over. From the 12 quantities describing the chargino 3–body–decay, four are eliminated by momentum conservation. These are chosen to be the angles $`\mathrm{\Theta }_1,\varphi _1`$ and the energies $`E_{1,2}`$, namely:
$`E_2`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{\stackrel{~}{m}_i2\stackrel{~}{E}E_++2\stackrel{~}{p}E_+\mathrm{cos}(\stackrel{~}{p},p_+)}{\stackrel{~}{E}E_+\stackrel{~}{p}\mathrm{cos}(\stackrel{~}{p},p_2)+2E_+\mathrm{cos}(p_2,p_+)}}`$ (18)
$`E_1`$ $`=`$ $`\stackrel{~}{E}E_2E_+`$ (19)
$`\mathrm{cos}\mathrm{\Theta }_1`$ $`=`$ $`{\displaystyle \frac{\stackrel{~}{p}\mathrm{cos}\stackrel{~}{\mathrm{\Theta }}E_2\mathrm{cos}\mathrm{\Theta }_2E_+\mathrm{cos}\mathrm{\Theta }_+}{E_1}}.`$ (20)
The angle $`\varphi _1`$ is fixed by two relations:
$`E_1\mathrm{sin}\mathrm{\Theta }_1\mathrm{cos}\varphi _1`$ $`=`$ $`\stackrel{~}{p}\mathrm{sin}\stackrel{~}{\mathrm{\Theta }}\mathrm{cos}\stackrel{~}{\varphi }E_2\mathrm{sin}\mathrm{\Theta }_2\mathrm{cos}\varphi _2E_+\mathrm{sin}\mathrm{\Theta }_+\mathrm{cos}\varphi _+,`$ (21)
$`E_1\mathrm{sin}\mathrm{\Theta }_1\mathrm{sin}\varphi _1`$ $`=`$ $`\stackrel{~}{p}\mathrm{sin}\stackrel{~}{\mathrm{\Theta }}\mathrm{sin}\stackrel{~}{\varphi }E_2\mathrm{sin}\mathrm{\Theta }_2\mathrm{sin}\varphi _2E_+\mathrm{sin}\mathrm{\Theta }_+\mathrm{sin}\varphi _+`$ (22)
We end up with the 8 parameters (these are given by Eq. (17) and $`\stackrel{~}{E},\stackrel{~}{\mathrm{\Theta }},\stackrel{~}{\varphi }`$).
For completeness, it is trivial to compute the sneutrino decay width. For $`m_{\stackrel{~}{\nu }}m_{\stackrel{~}{\chi }}`$ (the scenario discussed in the text, see Fig.4) only one decay channel to an $`e^{}\mu ^+`$ pair is open (for simplicity we assume that only one $`\lambda _{112}`$ dominates):
$$\mathrm{\Gamma }_{\stackrel{~}{\nu }}=\lambda _{112}^2m_{\stackrel{~}{\nu }_e}/8\pi .$$
(23)
Finally for the combined process of production and decay we obtain in the narrow width approximation:
$`d\sigma (e^{}e^{}2\mu ^{}e^{}e^+)`$ $`=`$ $`{\displaystyle \frac{1}{2s}}dLips(s,p,\stackrel{~}{p})dLips(\stackrel{~}{m},p_1,p_2,p_+){\displaystyle \frac{1}{2\pi }}`$ (24)
$`{\displaystyle \underset{\stackrel{~}{\lambda }}{}}|M(;,\stackrel{~}{\lambda })T\left(\stackrel{~}{\lambda }\right)|^2\left({\displaystyle \frac{\pi }{m\mathrm{\Gamma }}}\right).`$
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# Square billiard with a magnetic flux
## I Introduction
The behavior of mobile charged particles confined to some region and also subjected to a magnetic flux has interested physicists since the discovery of the Hall effect over a century ago. Quantum effects turned out to be subtle and surprising, as attested by Landau diamagnetism and the Aharonov-Bohm effect. About twenty years ago, with the advent of the quantum Hall effect and mesoscopic systems, the two dimensional case became prominent. At about the same time, the development of the subject of “quantum chaos” also focussed interest on such systems as being among the simplest of the “Gaussian unitary ensemble” universality class.
There is thus a long history of work on the confined quantum motion of charged particles in a magnetic flux. It is remarkable that all previous workers overlooked the fact that many fundamental cases, some of which have been extensively studied numerically, can be solved and classified analytically to good approximation. Moreover, the results have an interesting and suggestive complexity.
In this paper we obtain good approximate solutions to a couple of simply posed and well studied problems of this type. A preliminary version has appeared electronically. Rather than stress the generality of our method, we focus on a typical problem: a two dimensional charged particle confined to a square billiard in a perpendicular magnetic flux. Certain conditions on the flux are required to justify the approximations, and we also require the energy of the particle to be large. These approximate solutions are compared to numerical solutions. The two flux configurations considered are a uniform flux and an Aharonov-Bohm flux line. It is crucial that the square billiard is integrable. Systems other than the square to which our methods apply will be mentioned at the end.
There are many recent research papers in which the basic system studied is a square or rectangular billiard with a magnetic flux. A number of these are inspired by the experiments of Lévy et al., which measure the magnetic susceptibility of a collection of a considerable number of mesoscopic, two dimensional metallic systems, each approximately a square. The field is nearly uniform over the square in this case.
In the presence of a magnetic flux, there is the possibility that persistent currents exist. In other words, it is possible that the equilibrium state has a nontrivial current. Indeed, we find that eigenstates of the quantum system have interesting current densities. Although the wavefunction has a fairly simple representation, the currents can be quite complex. Some states are predominantly paramagnetic, others are predominantly diamagnetic. Still others may support both strong paramagnetic and strong diamagnetic currents which in total nearly cancel. However, these states strongly affected by the field are rather rare and most states have weak persistent currents.
Ideas from the field of quantum chaos have also motivated much work. Since the magnetic field intuitively has circular symmetry, which ‘conflicts’ with the symmetry of the square, one might expect chaos to ensue, as is the case with the Sinai billiard. Another theme of quantum chaos is that of energy level and wavefunction statistics. These statistics depend on whether time reversal symmetry (or other antiunitary symmetry) is in force. A natural way to break time reversal is by a magnetic flux. \[The square with a uniform flux still has an antiunitary symmetry, however.\]
Diffraction effects, in which a classical length shorter than the wavelength becomes important, are much studied in this context. This is obviously the case for the zero radius Aharonov-Bohm flux line. The sharp corners of the square also cause diffractive effects in the presence of a uniform magnetic field. These effects are much smaller than for the flux line, of course. We give estimates for the parameter range in which such diffraction becomes important, although we defer study of these effects.
Our main motivation however, is that we add to the store of solvable problems, and perhaps suggest some experiments. In the textbooks, there are relatively few such integrable problems, basically, only those which reduce to one dimension, or separate into several one-dimensional problems. The square without a magnetic flux is such a case in which the $`x`$ and $`y`$ motion separates.
The traditional way to widen the class of approximately solvable problems is perturbation theory and indeed, our approach is a form of quantum perturbation theory. We are able to study systems which are integrable except for a ‘classically weak’ perturbation. Of course, treating weak perturbations classically is challenging, because the long time behavior may be chaotic. However, quantum perturbation theory is better behaved and depends on the short time rather than the long time classical behavior.
On the other hand, perturbations small in this sense can give rise to very large quantum effects, especially on the wave functions. Moreover, even perturbations which change some length scale by an amount $`\delta L<<\lambda ,`$ where $`\lambda `$ is the wavelength, can have big effects. In terms of the standard nomenclature, a degenerate perturbation theory is required, and a number of unperturbed states are strongly mixed together to give the final result. Usually, this is done by diagonalization of a small matrix, but in our case the ‘matrix’ can be quite large. However, rather than just diagonalizing some matrix by a computer calculation, we obtain the result by an intuitively appealing Schrödinger differential equation.
Having solved, for the first time, this interesting class of problems, we discovered some other methods of obtaining the solution at the same level of approximation. We shall present these methods elsewhere.
## II Square in a uniform field
We begin with the case of uniform field. For given velocity $`v`$, the cyclotron radius is $`R_c=v/\omega _c=cp/eB`$ where $`\omega _c=eB/mc,`$ and $`p=mv.`$ The momentum $`p=\mathrm{}k=h/\lambda `$ is quantally related to wavenumber $`k`$ and wavelength $`\lambda `$. We define the classical small parameter $`ϵ=L/R_c=eBL/\mathrm{}ck=2\pi \varphi /\varphi _0kL`$. Here $`L`$ is the length of the side of the square, $`\varphi `$ is the magnetic flux $`BL^2`$ and $`\varphi _0`$is the flux quantum $`hc/e.`$ Small $`ϵ`$ allows us to approximate orbits within the square as straight lines, to first approximation. This is sometimes known as the Aharonov-Bohm regime, since the leading quantum effects come from the phase interference effects associated with the vector potential, and do not depend on the change of classical orbital motion caused by the Lorentz force. Many potential experiments are in this parameter range.
We choose units such that the dimensionless field is $`B2\pi \varphi /\varphi _0,`$ i.e., $`2\pi `$ times the number of flux quanta in the square. We take $`L,\mathrm{}`$ and $`2m`$to be unity so that
$$ϵB/k<<1.$$
(1)
The dimensionless wavenumber $`k`$ is the number of wavelengths in a side of the square, up to a factor of $`2\pi .`$ It satisfies
$$k>>1,$$
(2)
which is the basis for the quasiclassical approximation.
We shall see that the condition for standard quantum perturbation theory to work is
$$k\sqrt{ϵ}<<1,$$
(3)
or, in other words, $`\sqrt{kB}<<1.`$ This is not completely obvious, and in other contexts, it has been guessed incorrectly that the quantum ‘perturbation border’ is $`kϵ<<1,`$ i.e. $`B<<1,`$ which has the simple meaning that the number of flux quanta in the square is small. We, however, find that nothing much changes at the border $`B1.`$
There may also be a ‘high energy’ condition in the form of a requirement that $`kϵ^b<<1.`$ The exponent $`b`$ depends on the smoothness of the perturbation. We find that for the uniform field, $`b=2`$, while for the ideal flux line $`b=1.`$ Note that for fixed $`B`$ the energy can be arbitrarily high, but if instead $`ϵ`$ is kept fixed, there is a limitation on the energy.
This high energy condition is basically the requirement that diffraction effects not be too important. From a semiclassical perspective, diffraction effects occur where the classical system has a length scale as short or shorter than $`\lambda .`$ An ideal flux line obviously gives rises to diffraction effects. Billiards also have such a length scale of course, namely the distance it takes the confining potential to change from zero inside the billiard, to infinity outside the billiard. This can be taken into account by a ‘Maslov phase’ of $`\pi `$ at the boundary, however. There are also the sharp corners of the square. The square corners whose angular opening is $`\pi /N,`$ where $`N=2`$ is an integer, are a special case at which no diffraction occurs. With such an angle, the billiard can be extended by reflection, and the corner in effect disappears. However, in the presence of a magnetic field, this reflection technique does not work, and with sufficiently large field, orbits which hit directly into the corner eventually become important to the semiclassics.
## III Bogomolny’s Quasiclassical Surface of Section Method
Our approach utilizes the quasiclassical surface of section \[SS\] method of Bogomolny. Poincaré’s surface of section is a surface in classical phase space through which all interesting orbits repeatedly pass. For two dimensional systems, the surface of section is a two dimensional phase space. For a billiard, the Birkhoff surface of section is often chosen. Namely, the space part of the surface of section represents a point on the boundary at which the orbit bounces and is usually measured by the distance along the boundary of the billiard. The variable conjugate to this is the component of momentum parallel to the boundary at the moment of contact. However, many possible surfaces of section can be considered, and some are more convenient than others.
Bogolmony’s method is a generalization of the “boundary integral method”, applicable for billiards, and based on Birkhoff’s surface of section, to much more general systems and surfaces of section. The boundary integral method introduces an operator $`K(x,x^{},E)`$ and an integral equation $`\psi (x)=𝑑x^{}K(x,x^{};E)\psi (x^{})`$. This exact equation has nontrivial solutions only when $`E`$ is on the spectrum. The SS wavefunction $`\psi (x)`$ is the normal derivative of the full wavefunction, $`\psi (x)=\mathrm{\Psi }(𝐫)/n,`$ when $`𝐫`$ is at the boundary point $`x.`$ We should mention that only recently has the boundary integral method been extended to uniform magnetic fields in the case that $`ϵ`$ is of order unity.
Bogolmony’s operator $`T(x,x^{};E)`$ is basically the quasiclassical approximation to $`K.`$ It thus takes the particle crossing the SS at position $`x^{}`$ to its next crossing at position $`x,`$ all at energy $`E=k^2.`$ The quasiclassical approximation to the spectrum is determined by the existence of solutions of $`T\psi =\psi .`$ If only the spectrum is of interest, as it has been for many authors, the condition may be expressed as $`det\left[1T(E)\right]=0.`$
The operator $`T`$ is given quite generally by
$$T(x,x^{};E)=\left[\frac{1}{2\pi i}\left|\frac{^2S(x,x^{};E)}{xx^{}}\right|\right]^{\frac{1}{2}}\mathrm{exp}\left[iS(x,x^{};E)\right]$$
(4)
where $`S=_x^{}^x𝐩𝑑𝐫`$ is the action integral along the classical path from $`x^{}`$ to $`x`$. Note that, rather than giving position and momentum on the SS, positions at two sequential crossings of the SS are given. It is assumed, for notational convenience, that there is a unique orbit from $`x^{}`$ to $`x.`$ We also suppress the Maslov phase. Note that $`T`$ is semiclassically unitary.
The classical action $`S(x,x^{})`$ generates the surface of section map. Namely, the momenta conjugate to $`x,x^{}`$ are given by
$$p=\frac{S(x,x^{})}{x};p^{}=\frac{S(x,x^{})}{x^{}}.$$
(5)
Eq. (5) implicitly gives the surface of section map $`\{p,x\}=\{p^{},x^{}\}.`$
We attempt to simplify $`T`$ by astute choice of the surface of section. It seems simpler to use just one side of the square, rather than all four sides. It is even easier to use a method of images. Namely, we consider, instead of a unit square, $`x,y[\frac{1}{2},\frac{1}{2}][\frac{1}{2},\frac{1}{2}],`$ an infinite channel of width 2 obtained by reflecting the original square first about $`x=\frac{1}{2}`$ and then about $`y=\frac{1}{2},`$ and finally repeating the resulting $`2\times 2`$ square periodically to $`x=\pm \mathrm{}`$. The flux changes sign in neighboring squares. This geometry is shown in Fig. 1. There are a continuum of channel solutions. The solutions to the original square are a subset of these, which exist only at certain quantized energies. This quantization can be carried out in several ways, one of which is shown below.
The SS is taken as the axis $`y=\frac{1}{2}`$ which is identified with $`y=\frac{3}{2}.`$ Because the field is classically weak, the path used to calculate the action is approximated by a straight line. We immediately find
$$S(x,x^{})=k\sqrt{4+\left(xx^{}\right)^2}+\mathrm{\Phi }(x,x^{}),$$
(6)
i. e. the flux free result plus $`\mathrm{\Phi }=(e/c)𝐀𝑑𝐫`$, where the integral is done along the straight line path.
Our scheme finds solutions of $`T\psi =\psi `$ by a perturbation theory. Taking advantage of the fact that $`T`$ is unitary, we first solve
$$𝑑x^{}T(x,x^{};E)\psi (x^{})=e^{i\omega (k)}\psi (x),$$
(7)
treating $`k=\sqrt{E}`$ as a parameter, and then find the energies by solving $`\omega (k)=2\pi n.`$ Given $`\psi (x),`$ a quadrature which can be carried out quasiclassically yields the full wave function $`\mathrm{\Psi }(x,y).`$ Details are given in Appendix A.
## IV Resonances
Classical and quantum perturbation expansions in powers of $`ϵ`$ fail near resonances, necessitating modifications which introduce $`\sqrt{ϵ}`$. Classical resonances correspond to periodic orbits of the unperturbed system. Periodic orbits on the square correspond to straight line orbits in the channel from ($`x^{},\frac{1}{2}`$) to ($`x=x^{}+2p/q,`$ $`\frac{3}{2}`$). Here $`q`$ is a positive integer and $`p`$ is a positive or negative integer relatively prime to $`q.`$ Negative and positive $`p`$ are not equivalent if there is a magnetic flux. Between resonances, or near resonances with large $`p`$ and $`q,`$ ordinary perturbation theory works. See Ref. for a fuller discussion.
We now specialize to the $`(\pm 1,1)`$ resonances. These are the simplest resonances depending strongly on the field and, as we shall see, in some sense dominate the magnetic response.
We look for a solution of Eq. (7) of the form $`\psi (x)=e^{i\kappa x}u_m(x)`$ where $`\kappa =k\mathrm{cos}45^{}=k/\sqrt{2}`$ and $`u_m`$ varies much more slowly than the exponential. The reason for this choice is that the phase factor $`e^{i\kappa x}`$ makes the rapidly varying phases in the integral $`𝑑x^{}T\psi `$ to be stationary at $`x^{}=x2.`$ This corresponds to rectangular shaped periodic orbits of the original square whose sides make angles of $`45^{}`$ with the $`x`$ axis. Such orbits are shown in Fig. 1.
Because $`ϵ`$ is small, the phase $`\mathrm{\Phi }`$does not greatly influence the position of the stationary phase and it suffices to evaluate $`\mathrm{\Phi }(x,x^{})`$ at $`\mathrm{\Phi }(x,x2)=\mathrm{\Phi }(x+2,x).`$ \[The accuracy of this approximation depends on the smoothness of $`\mathrm{\Phi },`$ and failure of the approximation is related to the onset of diffraction effects mentioned earlier.\] $`\mathrm{\Phi }(x+2,x)`$ is obtained by integrating the vector potential about the closed rectangular loop, and somewhat remarkably is independent of gauge. The result is that $`u_m`$ satisfies the Schrödinger equation
$$u_m^{\prime \prime }+V(x)u_m=E_mu_m$$
(8)
where $`V(x)=k\mathrm{\Phi }(x+2,x)/`$ and $`=\sqrt{8}`$ is the length of the periodic orbits. Thus we convert the phase $`\mathrm{\Phi }`$ to a ‘potential’ $`V.`$ The transverse energy $`E_m`$ enters into the total energy of the eigenstate, according to Eq. (17) below, and $`m`$ is one of two quantum numbers classifying the states.
## V Transverse potential and periodic orbits
For a given resonance and surface of section, there is an effective potential $`V`$ which determines the functions $`u_m`$ and the energy $`E_m.`$ The resonance classically corresponds to a continuous set of nonisolated periodic orbits of the integrable problem. Before perturbation, each of these orbits has the same action. The potential $`V`$, to leading order, is proportional to the change of the action under perturbation, calculated along the unperturbed path. Each such path is labelled by the parameter $`x`$, where it crosses the surface of section. Higher order corrections may also be found.
Knowledge of the potential gives much qualitative insight into the problem. Its minima, \[if smooth\], are at stable periodic orbits, as a rule, and its maxima are at the unstable orbits. In that sense, it represents a classical island chain. Of course, it was known how to quantize states near the stable periodic orbits, if a harmonic expansion is allowed. However, the states that can be found with the aid of $`V`$ are much more general and in particular the states with energies $`E_m`$ near or even above the maxima of the potential can also be found.
In general, isolated unstable periodic orbits do not support wavefunctions, but rather ‘scar’ them. In other words, there appears some excess weight on the wavefunction near the unstable orbit. In the sense of Feynman’s path integral formulation, there are not enough classical paths ‘near’ the unstable orbit, to build a complete wavefunction. Here ‘near’ means that the paths are close to the periodic orbit in the sense of being well approximated by a quadratic expansion about the periodic orbit.
The same is true in the present case, and wavefunctions cannot be built just from orbits near an unstable periodic orbit. However, because of the small parameter $`ϵ`$, we can approximate well an entire shell of orbits in the Feynman integral, and express the result in terms of the potential $`V(x).`$ This shell can support many states which we find. There are states whose energies are near the maxima of $`V`$ and thus have extra weight near the unstable periodic orbits.
The interpretation of $`u_m`$ is that it gives the structure of the wavefunction ‘transverse’ to the resonant periodic orbits. Along the periodic orbits, the wavefunction varies rapidly, but transversely, it varies relatively slowly. The ‘longitudinal’ and transverse motions are weakly coupled, because $`V`$ and thus $`u_m`$ and $`E_m`$ depend on $`k,`$ but this is easy to take into account.
The concept of ‘transverse’ is a little murky in the quantum case, although there are cases, including the one under study where it can be made more precise. We shall not dwell on this further in this paper, however.
We also remark that the Schrödinger equation (8) requires boundary conditions, in order to pick out the physically interesting solutions. These boundary conditions come from the properties imposed on the solution by the physics of the problem, and are usually simplified by symmetries of the problem.
## VI Uniform field solution
### A Effective potential
For the uniform field, the potential is
$`V(x)`$ $`=`$ $`Bk\left(\frac{1}{2}2x^2\right)/;x[\frac{1}{2},\frac{1}{2}],`$ (9)
$`V(x)`$ $`=`$ $`+Bk\left[\frac{1}{2}2(x+1)^2\right]/;x[\frac{3}{2},\frac{1}{2}],`$ (10)
$`V(x)`$ $`=`$ $`V(x+2).`$ (11)
The factor $`\frac{1}{2}2x^2`$ is simply the area enclosed by a periodic $`(1,1)`$ resonant orbit in the shape of a rectangle which bounces from the bottom of the square at $`x.`$ In Appendix B we give the corresponding potential for other resonances.
This periodic potential consists of alternating positive and negative harmonic potential wells of depth $`Bk/2.`$ At the boundaries $`x=\pm \frac{1}{2},`$ the second derivative of the potential is discontinuous, a fact which leads to the mentioned diffraction effects at sufficiently large $`B^2/k.`$
For the $`(1,1)`$ resonance, whose orbits are time reversed $`(1,1)`$ orbits, $`V(x)`$ changes sign. This would not be true if $`V`$ had its origin in a time reversal invariant perturbation of the square, for example, a small change of shape. We can include the $`(1,1)`$ resonance in the present scheme by attributing the region $`1/2<x<3/2`$ to that resonance. This extension of the $`x`$ coordinate is thus similar to use of a ‘angle’ variable, with positive $`x`$-velocity $`v_x`$ for $`x[\frac{1}{2},\frac{1}{2}],`$ and negative $`v_x`$ for $`x[\frac{1}{2},\frac{3}{2}].`$
If $`\sqrt{Bk}=k\sqrt{ϵ\text{ }}`$ is small, the potential $`V(x)`$ can be treated perturbatively. On the other hand, for sufficiently large $`Bk/`$, Eq. (8) will have low lying tight binding harmonic oscillator type solutions centered at $`x=0,`$ (if $`B>0`$), with energies approximately given by
$$E_m=\frac{1}{2}Bk/+(m+\frac{1}{2})\sqrt{8Bk/}.$$
(12)
This formula holds for $`m<<\sqrt{Bk/}`$. The lowest wavefunction is approximately $`u_0(x)=e^{\sqrt{Bk/2}x^2}`$ which is arbitrarily narrow at large energy. These states are paramagnetic, as follows from the fact that $`E_m/B<0.`$ This will be seen more clearly below.
Eq. (8) is valid for larger $`m.`$ Although very simple analytic answers are not available, the problem is the well known one of a particle in a one dimensional periodic potential. We shall see below that we need only consider the boundary conditions $`u(x+2)=\pm u(x).`$ This simplification is a consequence of the symmetry of the square, and something slightly more complicated would be needed for the rectangle.
The solutions to Eq. (8) may be put into four classes, $`A,B,C,D`$. Class $`A`$ states are those with ‘low’ energies near the bottom of the well, $`E_m\frac{1}{2}Bk/`$. For these cases $`u(x)`$ has support only near $`x=0`$, $`\pm 2,\pm 4,\mathrm{}`$ These localized states are strongly paramagnetic, that is, the current circulates in the opposite direction from that of the particle in free space. In this case $`dE_m/dB<0.`$ \[We shall see that the transverse energy $`E_m`$ carries nearly all the field dependence of the total energy of the corresponding two dimensional eigenstates.\]
Class $`D`$ states have energies much greater than the maximum potential energy, i.e. $`E_m>>\frac{1}{2}Bk/`$. In this case, the magnetic field is a small perturbation, since the ‘potential energy’ $`V(x)`$ in Eq. (8) is small compared with the ‘kinetic energy’ given by $`u^{\prime \prime }.`$ These states are weakly diamagnetic. We shall not consider further this case. Of course, the approximation of expanding about the $`(1,1)`$ resonance eventually breaks down, and higher order resonances are eventually involved.
Class $`C`$ has total transverse energy near the top of the potential $`V(\pm 1),`$ that is, $`E_m\frac{1}{2}Bk/`$ and the states are strongly affected by the magnetic field. Very crudely, they are somewhat localized or ‘scarred’ near $`x=\pm 1,`$ since they spend more time in that region. This means that they are strongly influenced by the $`(1,1)`$ resonance. They are diamagnetic and $`dE_m/dB>0.`$
States of class $`B`$ form a transition region between the low-energy paramagnetic states, and the higher energy diamagnetic ones, i.e. near $`E_m0`$. These are states strongly affected by the field, but are such that $`dE_m/dB0.`$
### B Quantization
There are two states with identical energy in the repeated square scheme. These are $`\psi _I=e^{i\kappa x}u_m(x),`$ and $`\psi _{II}=e^{i\kappa x}u_m(x1).`$ \[Changing the sign of $`\kappa `$ is equivalent to changing the sign of the field, which in turn can be accomplished by replacing $`V(x)`$ by $`V(x+1).`$\]
Rather than finding the eigenvalues by imposing conditions directly on the $`\psi `$’s, as in the Appendix A, we produce the four two-dimensional solutions $`\mathrm{\Psi }`$ corresponding to $`\psi _{I,II}.`$ \[Each $`\psi (x)`$ gives two $`\mathrm{\Psi }(x,y)`$’s because $`y`$ and $`1y`$ in the strip represent the same point in the original square.\] We show elsewhere that these states can also be found directly by a Born-Oppenheimer approximation. One of these states may be written
$$\mathrm{\Psi }_0(x,y)=e^{i\kappa (x+y)}u_m(xy\frac{1}{2})$$
(13)
The remaining three states, $`1,2,3,`$ can be obtained by rotations, e.g. $`\mathrm{\Psi }_1(x,y)=\mathrm{\Psi }_0(x,y)=\mathrm{\Psi }_0(y,x),`$ etc. Here $`:(x,y)(y,x)`$ is the rotation by $`90^{}.`$ The gauge can be chosen so that the Hamiltonian is invariant under $``$. Therefore, the symmetry of an eigenstate can be labelled by $`r=0,1,2,3`$, where the eigenvalue of $``$ is $`i^r.`$
Thus, an eigenfunction with symmetry $`r`$ is given by
$$\mathrm{\Psi }_{(r)}(x,y)=\left(\underset{s=0}{\overset{3}{}}i^{rs}^s\right)e^{i\kappa (x+y)}u_m(xy\frac{1}{2}).$$
(14)
\[The sequence of rapidly varying phase factors is $`\{e^{i\kappa (x+y)},e^{i\kappa (x+y)},e^{i\kappa (x+y)},e^{i\kappa (xy)}\}.`$ These in turn are rapidly varying in the $`45^{}`$ directions of the sides of the periodic orbits.\] In general, a solution of Eq. (8) satisfies the boundary condition $`u_m(x+2)=e^{i\beta }u_m(x).`$ We need to find the allowed values for $`\beta `$ and $`\kappa `$ which will give the quantized energies. These conditions are obtained by requiring $`\mathrm{\Psi }_{(r)}(x,\frac{1}{2})`$ to vanish, corresponding to Dirichlet conditions in the original problem. If the wavefunction vanishes on the bottom, it will by symmetry vanish on the boundary of the square.
Clearly, the sum of the two terms \[$`s=0,3`$\] in Eq. (14) which are proportional to $`e^{+i\kappa x}`$ must vanish. This implies $`u_m(x)=i^{3r}e^{i\kappa }u_m(x).`$ The reflection symmetry $`V(x)=V(x)`$ allows us to take $`u_m(x)=(1)^mu_m(x)`$. In turn, this allows quantization of $`\kappa `$ in the form $`\kappa =n\pi /2,`$where $`n`$ is an integer satisfying certain conditions depending on $`r`$ and $`m.`$ Similarly the two terms proportional to $`e^{i\kappa x}`$ in Eq. (14) give the condition $`e^{i\beta }=(1)^r.`$ The relationship of $`n`$ to $`r`$ and $`m`$ is
$$nmod4=\left[2(1mmod2)+r\right]mod4.$$
(15)
It is straightforward to find for the eigenwavenumber
$$k_{n,m}=2\pi n/+E_m/k.$$
(16)
Eq. (16) should be solved iteratively. For example, the first approximation replaces the $`k`$ dependence of the term $`E_m/k`$ by $`2\pi n/.`$ Equivalently, the energy
$$E_{n,m}=4\pi ^2n^2/^2+2E_m.$$
(17)
Note that $`E_m`$ depends, relatively weakly, on $`n`$, since the $`k`$ in formula Eq.(12) should be replaced by $`2\pi n/`$. The dependence of the total energy on $`B`$comes through the term $`E_m.`$ Eqs. (16) and (17) hold for all symmetries and successive values of $`n`$ at fixed $`m`$ cycle through the representations of $``$. Note that, since $`E_m/k<<k`$ the wavelength is given approximately by $`/n,`$ i.e. the length of the classical orbits is an integer number of wavelengths.
Thus we have an expression for the energies of a class of states, namely the $`(\pm 1,1)`$ resonant states. They are labelled by integer $`n`$ which effectively give the number of wavelengths measured along the $`(1,1)`$ periodic orbits, and by a second integer $`m`$ which gives the number of ‘nodes’ ‘perpendicular’ to this orbit. The very low $`m`$ states could very well have been found by earlier methods, since they can be obtained by expansions about the stable periodic orbits. However, these remarkable states do not seem to have been noticed heretofore.
### C Orbital magnetism
The $`(1,1)`$ states just obtained dominate the magnetic orbital susceptibility in a parameter range appropriate to experiments. The susceptibility for the square is on a scale rather larger than the Landau diamagnetism. It is of course not necessary to find the states, or for that matter, their energies, to calculate the susceptibility. That is because the susceptibility depends only on the density of states smoothed over an energy width proportional to the temperature. The Gutzwiller or better, the perturbed Berry-Tabor trace formula is designed to give exactly that quantity in quasiclassical approximation. Nevertheless, it’s interesting and previously unremarked, that a small subset of states accounts for most of the magnetism.
We start by finding the orbital susceptibility $`\chi `$ of a system of noninteracting electrons in a grand canonical ensemble. This is given by $`\chi =/B`$ where the magnetization $`=\mathrm{\Omega }(T,\mu ,B)/B.`$ Here the grand potential is
$$\mathrm{\Omega }(T,\mu ,B)=k_BT\underset{a}{}\mathrm{ln}\left[1+e^{(E_a\mu )/k_BT}\right].$$
(18)
The temperature is $`T`$, $`k_B`$ is Boltzmann’s constant, and $`\mu =k_F^2`$ is the chemical potential. The dependence of $`\mathrm{\Omega }`$ on $`B`$ comes only because the eigenenergies $`E_a`$ depend on $`B.`$ The sum is over all eigenstates labelled by $`a`$.
We divide the states $`a`$ into those relatively few whose energies depend appreciably on the field and the rest. These field dependent states are exactly the $`(1,1)`$ states found above, plus possibly states classically associated with a few other low resonances, e.g. $`(1,3)`$. The reason for this is that the $`(1,1)`$ states enclose the maximum directed area. They also have the shortest length $``$ which we will see plays a role. The even shorter $`(0,1)`$ periodic orbits do not enclose any flux in the approximation of neglecting the curvature of the orbits although at higher fields they eventually become important.
Thus, we replace $`_a=_b+_{n,m}`$ and we can neglect the sum $`b`$ over field independent states. The second sum, over $`(1,1)`$ resonance states, has many fewer terms than the first in a given range of energy. Since $`\mu `$ is related to the number of particles, it is nearly independent of $`B.`$ It is possible to find $`\mu =k_F^2+\delta \mu (B)`$ and make a consistent expansion, and that is indeed necessary if an average over a large number of squares with canonical statistics is done. However, we just want to illustrate how the $`(1,1)`$ states dominate the susceptibility, and we will not consider this further average. Then, we may approximate
$$(T,\mu ,B)=\underset{n,m}{}\frac{E_{n,m}}{B}f_D\left[E_{n,m}(B)\right].$$
(19)
and $`f_D`$ is the Fermi-Dirac distribution function.
Using the Poisson sum formula, replace the sum on $`n`$ in Eq. (19) by an integral over $`k`$, and do the integral to obtain
$``$ $`=`$ $`{\displaystyle \frac{k_BT}{k_F}}{\displaystyle \underset{r,m,s=0}{\overset{\mathrm{}}{}}}\alpha _m\mathrm{exp}\left({\displaystyle \frac{\omega _rs}{2k_F}}\right)`$ (21)
$`\times \mathrm{sin}\left[s\left(k_F{\displaystyle \frac{E_m}{k_F}}\right)\right].`$
Here, $`\omega _r=\pi (2r+1)k_BT`$ is the Matsubara frequency and $`\alpha _m=E_m/B.`$ \[We have dropped the ‘leading’ term in the Poisson formula which totally neglects the discrete quantum nature of the states and which therefore cannot produce a magnetization.\] As an example, take $`k_BT`$ ten times the level spacing of all levels, i.e. $`k_BT=20\pi `$ in our units. Then, $`\omega _0/2300.`$ If $`k_F300600,`$ so that the square contains about 2-6$`\times `$10<sup>4</sup> electrons, the exponential suppression will not be too serious for $`r=0,s=1`$. However, larger $`r`$ or $`s`$ do not contribute much. \[In the trace formula approach, $`s`$ gives the number of repetitions of the primitive periodic orbit and the sum over $`r`$ is explicitly carried out.\]
Eq. (21) shows that states with larger $`,`$ i.e. smaller spacing, are suppressed, exactly as seen from the trace formula in terms of periodic orbits. It also shows that relatively large field dependence of the levels, $`\alpha _m,`$ is important. For the square, the $`(1,1)`$ states have the smallest $``$ and also the largest $`\alpha _m`$. The $`(2,1)`$ resonance does not couple to a small constant field.
It is also seen that if Eq. (21) is averaged over many squares of somewhat different sizes, because of the oscillations of the sine, the result is much reduced and it is necessary to go to higher order in $`\delta \mu .`$
In Fig. 2 we show $`E_m/\widehat{B}`$ as a function of $`E_m/\widehat{B}`$ for the (1,1) resonant states, where $`\widehat{B}=Bk/2`$. According to Eq. (21), if the sign of $`E_m/\widehat{B}`$ is negative the contribution of the corresponding state is paramagnetic. For an exact, two dimensional wavefunction $`\mathrm{\Psi }_{n,m}`$ with energy $`E_{n,m}`$, it is known that $`E_{n,m}/\widehat{B}`$ $`=(/k)d^2r[𝐫\times 𝐣(x,y)]_z`$ , i.e. the expectation value of the $`z`$-component of the magnetization density. Here the current is
$$𝐣(x,y)=2Re\mathrm{\Psi }^{}(x,y)\left(\frac{1}{i}𝐀(x,y)\right)\mathrm{\Psi }(x,y).$$
(22)
In our approximation, according to Eq. (17) above, $`E_{n,m}/\widehat{B}`$ is equal to $`2E_m/\widehat{B}.`$ It follows from Eqs. (8) and (11) above, that $`E_m/\widehat{B}=\widehat{V}(x)_m=u_m\left|\widehat{V}(x)\right|u_m`$ where $`\widehat{V}=V/\widehat{B}.`$ Since $`\widehat{V}<0`$ for $`x[\frac{1}{2},\frac{1}{2}],`$ and $`\widehat{V}>0`$ for $`x[\frac{1}{2},\frac{3}{2}],`$ etc. we see that the sign of the magnetic response of a given wavefunction depends on which region of $`\widehat{V}`$ dominates the expectation value. Of course, the classical periodic orbits in these two regions have the expected sense.
Finally, we may express $`E_m/\widehat{B}`$ quasiclassically as
$$\frac{E_m}{\widehat{B}}=\frac{𝑑x\widehat{V}(x)\left[\widehat{E}\widehat{V}(x)\right]^{\frac{1}{2}}}{𝑑x\left[\widehat{E}\widehat{V}(x)\right]^{\frac{1}{2}}}.$$
(23)
Here $`\widehat{E}=E_m/\widehat{B}`$ can be treated as a continuous variable, so that $`E_m/\widehat{B}`$ as a function of $`\widehat{E}`$ falls on a continuous curve which in this approximation is independent of $`\widehat{B}.`$ The integrals are between the turning points. For $`\widehat{E}`$ near the minimum $`\widehat{V}_{\mathrm{min}}`$ of $`\widehat{V}`$ , $`E_m/\widehat{B}\widehat{V}_{\mathrm{min}}.`$ For $`\widehat{E}`$ at the maximum of $`\widehat{V}`$, the integrals diverge at $`x=1`$, thus making $`E_m/\widehat{B}`$ positive, and also giving the cusp in Fig. 2. Of course, our simple quasiclassical approximation needs corrections in this case.
The quantization of the transverse motion relates $`\widehat{E}=E_m/\widehat{B}`$and $`m`$ according to the standard formula
$$𝑑x\sqrt{\widehat{E}\widehat{V}(x)}=\pi (m+\frac{1}{2})/\sqrt{\widehat{B}}.$$
(24)
The spacing of the quantum states therefore depends on $`\widehat{B}`$. The level $`m_{top}`$ such that $`E_{m_{top}}V(1)`$ satisfies, approximately, $`m_{top}0.6\sqrt{\widehat{B}}.`$ The states with $`k=142`$ and $`B=25`$ were chosen in part because with these parameters there is a state very near the diamagnetic maximum $`m_{top}=14,`$ and another state with $`m=10`$ for which $`E_m/\widehat{B}`$ is very small. On the other hand, with these parameters, while our approximate states are very good for small $`m,`$ the transverse momenta are becoming large enough (as compared with the longitudinal momenta) that our approximation deteriorates significantly at $`m10,`$ especially in regions of the square where the wavefunction is small.
### D Visual representations of eigenstates
Eigenstates of Hamiltonians not satisfying time reversal invariance are rarely shown graphically in the literature, except for trivial cases. In the invariant case, numerically produced picture galleries of eigenstates for systems such as the Bunimovich stadium, have led to a great deal of interest and insight, both theoretical and experimental.
‘Magnetic’ states are not real but are inherently complex. A complete graphical picture of a complex state would seem to require twice the number of pictures as that necessary for a real state satisfying time reversal invariance. In addition, the states themselves are gauge dependent, and only gauge invariant quantities are physically meaningful.
A difference of the magnetic case as compared with the time reversal invariant case is that generally $`\mathrm{\Psi }`$ or $`𝐣`$ do not vanish along nodal lines, but rather only at isolated nodal points. There may of course be symmetries or boundary conditions requiring these quantities to vanish along a line, but in general the representation of wavefunctions by their nodal patterns is not available in the absence of time reversal invariance.
Two gauge invariant quantities we choose to display are the absolute value squared of the wavefunction $`\left|\mathrm{\Psi }(x,y)\right|^2`$ and the current density. The current density is a two dimensional vector field that is divergence free, $`𝐣=0.`$ The one-dimensional surface of section states are also of interest.
We first show a picture of a class $`A`$ state, which is quite simple to represent. Fig. 3 shows $`\left|\mathrm{\Psi }(x,y)\right|^2,`$ for $`n=62,k2\pi 62/140,B=31.4,`$ $`\sqrt{Bk}66,`$ and $`m=0,`$ which implies the symmetry $`r=0.`$ For such a well localized $`u_m`$, each term in Eq. (14) dominates one side of the rectangular periodic orbit. For example, near $`x=y=\frac{1}{4}`$ only the first term, $`s=0,`$ in the sum (14) is appreciable. In this region $`\left|\mathrm{\Psi }(x,y)\right|`$ $`u_0(xy\frac{1}{2})`$ which is well approximated by a Gaussian.
There is interference near the square edges (e.g. near ($`0,\frac{1}{2}`$)), and two terms contribute appreciably. Near this point then
$`\mathrm{\Psi }(x,y)`$ (25)
$``$ $`\left|e^{\frac{1}{2}n\pi iy}u_0(xy\frac{1}{2})+e^{\frac{1}{2}n\pi iy}u_0(xy\frac{1}{2})\right|`$ (26)
$``$ $`\left|2u_0(x)\mathrm{cos}\frac{1}{2}n\pi y\right|`$ (27)
so that $`\left|\mathrm{\Psi }(0,\frac{1}{2}+\frac{1}{n})\right|`$ at its first maximum near $`(0,\frac{1}{2})`$ is about twice as large as $`\left|\mathrm{\Psi }(\frac{1}{4},\frac{1}{4})\right|`$. \[Note that $`nmod4=2.`$\] The current density is thus largest close to the middle of the square edges and it is of course nearly parallel to the edge there. In this case, the shape of $`\mathrm{\Psi }`$ close to an edge is given by $`u_m(x),`$ as can be seen in Fig. 3.
Fig. 4 shows the streamlines of the current in this state. The direction of the streamline gives the direction of the current flow, while the density of streamlines is proportional to the magnitude of the current density. That is, between any two neighboring streamlines, the same total current flows. The state in Fig. 4 is paramagnetic, that is, the current circulates in the opposite sense from that of a free particle in the field. The choice of which streamline to display is easily obtained in this case, since each line crosses a symmetry line like $`x=0,y[\frac{1}{2},0],`$ once and only once. The current has but one interior zero, at the center of the square.
Fig. 5. shows the state $`n=60,`$ $`B=31.4,`$ $`m=1`$, which has one ‘transverse node’. This node does not give rise to a nodal line in the total wavefunction, of course, although the wavefunction is small in regions corresponding to the node. The states have rotational symmetry, and we show a different representation in each corner of the unit square.
One remarkable feature of these states, $`m=0,1`$ is that they are localized very near the central paramagnetic diamond orbit. Theoretically, for large $`Bk`$ this localization can be as tight as one pleases. Although this is a result of our theory which uses a ‘potential’ function $`V(x),`$ it is clear that $`V(x)`$ is really something that arises from Aharonov-Bohm phases, and there is no classical localization based on energy considerations. Although a detailed analysis of this can be a little tricky, we find the same kind of effect in the AB flux line case, for large $`Bk`$ and small $`B.`$ This localization is thus like Anderson localization in the sense that, absent phase interference, localization would not exist. Of course, the random disorder aspects of Anderson localization are absent.
The next figures, Figs. 6-11, have $`n=62,B=25.`$ Figs. 6,7 show a state of class $`C`$ corresponding to $`m=14`$ which is energetically at the top of the periodic potential corresponding to the $`(1,1)`$ resonance. This state is diamagnetic with the current circulating in the opposite sense from the states with $`m=0,1.`$ The theoretical and numerical wavefunctions are shown.
The theoretical predictions in this case are considerably less good than for the states with small $`m.`$ First, the transverse wavenumber is no longer quite so small compared with the wavenumber along the path. Second, since according to Eq. (14), as many as four approximately found component states are added, there may be relatively large errors, especially in those parts of the square where destructive interference is important and the final wavefunction is small. Nevertheless, Eq. (14) is a quite good representation of the more accurate numerical results.
In Fig. 8 we show the state $`m=10`$, for which $`E_{10}/B`$ is very small. The streamlines of this state are rather striking. Note that there are large current loops which have opposite magnetic polarity. Again, the approximate wavefunction captures many features of the exact one, although it does not reproduce the finer details. In this case we show also, in Fig. 9, the transverse state $`u_{10}(x)`$ and the normal derivative of $`\mathrm{\Psi }_{62,10}`$ on the surface of section. Although for small $`m,`$ the normal derivative on the surface of section bears an understandable relation to $`u_m,`$ it is quite complicated for transverse energies this large.
Finally, in Fig. 10 we show streamlines for the sequence of states $`m=6,8,10,12.`$ Although there are systematic changes of pattern, we have not tried to rationalize these changes. We conclude that even though the wavefunction of Eq. (14) is fairly simple, it is difficult to foresee interference patterns when four terms are important.
These relatively complicated states are harder to represent adequately. In both theory and numerics the $`𝐣=0`$ character of the current is not exact. Numerically following a streamline, particularly in the neighborhood of a zero of the current is difficult. We therefore imposed the divergence free character by representing $`𝐣(x,y)=\times \widehat{𝐳}\chi (x,y).`$ We calculated $`\chi `$ as a symmetrized integral of $`j_x,`$ where $`j_x`$ was obtained either numerically or theoretically. The streamlines are then contour lines of $`\chi `$.
The diamagnetic states are not so spatially localized as the low $`m`$ states, although they continue to have a sort of localization in ‘momentum’ space, as we show below. More generally, as $`m`$ increases, the states become more delocalized, and eventually become independent of $`B.`$ This means that for larger $`m`$ all four terms in Eq. (14) make comparable contributions at an arbitrary typical point $`x,y`$, whereas in the localized case, only one or two terms contribute. This gives interference oscillations in $`\left|\psi _m\right|`$ near $`\left|x\right|\frac{1}{2}`$ as shown in Fig. 9.
### E Momentum localization
We have assumed that $`n>>m,`$ and that $`u_m`$ is slowly varying compared with $`e^{i\pi nx/2}.`$ We may expand $`u_m=\widehat{u}_{m,l}e^{i\pi lx}`$, where $`l`$ is an integer for $`r`$ even and half odd integer for $`r`$ odd. Also $`\widehat{u}_{m,l}=(1)^m\widehat{u}_{m,l}.`$ The unperturbed states \[$`\mathrm{sin}\pi p(x+\frac{1}{2})\mathrm{sin}\pi q(y+\frac{1}{2})`$\] can be labelled by integers $`p,q`$ with unperturbed energies $`(p^2+q^2),`$ dropping a factor $`\pi ^2.`$
Eq. (14) is a superposition of unperturbed states with quantum numbers $`p=\frac{1}{2}n+l,`$ $`q=\frac{1}{2}nl.`$ In particular, the energies $`p^2+q^2=\frac{1}{2}n^2+2l^2`$ are closer to the base energy $`\frac{1}{2}n^2`$ than to the base energy of the next representation, $`ϵ_{n+1}\frac{1}{2}n^2+n.`$ Of course, if the perturbation is symmetric under rotation, the next base energies coupled are $`ϵ_{n\pm 4}ϵ_n\pm 4n.`$
There are, however, other unperturbed states with $`p^2+q^2ϵ_n.`$ For example, $`7^2+49^2=ϵ_{70}.`$ However, the matrix elements of a smoothly perturbed Hamiltonian, $`_{pq,p^{}q^{}}`$, in the unperturbed basis, are small if $`\left|pp^{}\right|`$ or $`\left|qq^{}\right|`$ is large. Because the perturbation due to a uniform field has a singular third derivative, these perturbations drop off as a power law, and this relatively long range effect in momentum space is the source of the diffraction corrections.
In Fig. 11 we show the magnitude of the amplitudes of the unperturbed states combining to make the state $`n=62,`$ $`B=25,`$ $`m=14.`$ The area of a circle is proportional to the square of the amplitude. Above the main diagonal we show the theoretical result. All amplitudes lie on the line $`31l,`$ $`31+l.`$ Below, we show the result for the numerical wavefunction. The circle $`p^2+q^2=231^2`$ is also shown. Our theory is starting to need corrections for $`m`$ this large.
Thus, an interpretation of our method which yields the $`(1,1)`$ resonance states of Eq. (14) is that we effectively diagonalize the Hamiltonian in a basis restricted to the unperturbed states nearly ‘degenerate’ with $`ϵ_n`$ and close to $`\frac{1}{2}n,\frac{1}{2}n.`$ This is the case for the uniform field, and indeed, we achieve agreement between full numerical diagonalization, diagonalization restricted to ‘degenerate’ states, and the procedure using the solution of the differential equation Eq. (8).
## VII Aharonov-Bohm Flux Line
The above approach can be generalized to deal with nonuniform flux configurations. To get the potential $`V`$ associated with the $`(1,1)`$ resonance, all that is needed is to be able to calculate the flux contained within a $`(1,1)`$ periodic orbit. We consider here the case of the Aharonov-Bohm flux line \[ABFL\]. Some further results are published elsewhere.
While much of the above discussion can be carried over to the ABFL, in the ideal case of a zero radius line there are strong diffraction effects which limit the applicability of our theory. We therefore begin by defining the alternative problem of a finite size flux line or tube. This is more realistic if actual experiments are contemplated.
Let $`\rho `$ give the linear scale of the flux tube. We may think of the flux as uniform inside a tube of this radius, or as being distributed in some way, say as a gaussian, with $`\rho `$ giving the scale of the distribution. \[We actually used a square tube of side $`\rho .`$ Even more accurately, we used four quarter strength tubes symmetrically located which allows a symmetry reduction in the numerics. In the approximation of our theory, this gives essentially the same result as a single circular tube. The field inside a single flux tube is $`B_0=\varphi /4\rho ^2`$ and $`\varphi `$ is the total flux.\] The typical angular deflection suffered by a classical particle traversing this field is
$$\delta \theta (\varphi /\varphi _0)/k\rho .$$
(28)
To avoid diffraction we require $`\delta \theta `$ to be small. This can be achieved, of course, if $`k\rho `$ is large, but that is not necessary. In the numerical work shown, we take $`\varphi /\varphi _0=0.1,`$ and $`\rho =0.01,`$ while $`k140.`$ An alternative and equivalent condition is to insist that, on the appropriately defined average, the terms in the Hamiltonian satisfy $`(eA/c)^2/2m`$ $`<<e𝐩𝐀/mc`$.
The results depend on where the flux line is located. We consider first the case that it is located at $`x=0,`$ $`y=\frac{1}{2}+a`$, where $`0a\frac{1}{2}.`$ Again, we consider states related to the (1,1) resonance. This leads as before to Eq. (8) but now we find a potential
$`V_{AB}(x)`$ $`=`$ $`Bk/;x[a,a],`$ (29)
$`V_{AB}(x)`$ $`=`$ $`+Bk/;x[1a,1+a],`$ (30)
$`V_{AB}(x)`$ $`=`$ $`0;x[a,a][1a,1+a]`$ (31)
for points in $`[\frac{3}{2},\frac{1}{2}],`$ and $`V_{AB}`$ is extended periodically by $`V_{AB}(x+2)=V_{AB}(x).`$ For the flux line, $`B=2\pi \varphi /\varphi _0.`$
Actually, Eq. (31) is for the ideal flux line. The ideal case has step function jumps in the potential which are smoothed out at finite $`\rho `$. In other words, the sharp jump at $`x=a`$, is replaced by a smooth rise beginning at $`x=a\rho `$ and ending at $`x=a+\rho .`$ The exact shape of the rise depends on the distribution of flux in the line. As long as the transverse wavelengths in the solution of Eq. (8) are long compared with $`\rho ,`$ the finiteness of $`\rho `$ does not play a significant role in our theory at this level of approximation.
For the ideal, zero radius ABFL, there are significant deviations from this scenario. Indeed, most matrix elements of the ABFL perturbed Hamiltonian in the unperturbed basis are infinite. However, it is a weak, logarithmic infinity, and our theory seems to capture the main shape of the wavefunction, although at the relatively low energies for which numerical results are available there are significant corrections. We consider these to be diffractive corrections, arising from a characteristic length shorter than the wavelength.
It is clear that $`Bk/`$ is the important parameter in the UF case, while for the ABFL, both $`Bk/`$ and $`a`$ are important. We begin with the case $`a=\frac{1}{4},`$ which has a ‘square well potential’ of width $`\frac{1}{2}`$ near $`x=0.`$ For sufficiently large $`Bk/,`$ there will be ‘tight binding’ solutions approximately $`u_m(x)=\mathrm{cos}(m+1)\pi x/2a,\left|x\right|<a,`$ and zero elsewhere. This expression holds for sufficiently small even $`m,`$ and for odd $`m`$ the cosine is replaced by the sine. The energy $`E_mBk/+\pi ^2(m+1)^2/4a^2.`$
Fig. 12 is for the ABFL case, with $`n=86,`$ $`m=0,`$ $`r=0`$ and $`a=\frac{1}{4}.`$ Fig. 12(a) shows $`u_0(x)`$ and $`u_0(x1)`$, its extension into $`x<\frac{1}{2},`$ reflected. For these parameters, $`u_0`$ is not extremely localized, and extends significantly outside $`[\frac{1}{2},\frac{1}{2}].`$ The remaining plots give $`\left|\psi _m\right|=\left|\mathrm{\Psi }/n\right|.`$ Fig. 12(b) plots Eq. (14) and 12(d) is from diagonalization in the limited basis of the ‘degenerate’ states. Clearly, these two approximations are nearly the same. Fig. 12(c) and 12(e) are obtained by numerical diagonalization in the complete unperturbed basis, for the finite size flux tube and the ideal flux line respectively. Clearly, for these parameters, there is significant diffraction from the flux tube. It is somewhat less than for the zero radius line but the two patterns have some resemblance. Diffraction evidently modifies the interference between different parts of the wavefunction in an irregular way. The overall shape of $`\left|\mathrm{\Psi }/n\right|`$ is well predicted by the theoretical $`u_0`$. To give a sense of the irregularity of diffraction effects on the wave function, we show in Fig. 12(f) the function for $`n=70,`$ a somewhat lower energy. In this case, according to theory, the $`u_0`$ is essentially the same, and the interference fringes have a somewhat longer wavelength, on the average. The diffractive effects are quite different in detail, however.
It is clear from this result that the ABFL localizes the particle, and to do so, it must of course exert a force. That a flux line can exert a transverse force is now well established.
In Fig. 13 we show currents from a state with $`a=\frac{1}{2},`$ i.e. the much studied case with the flux line at the center of the square. Again, $`m=0,`$ $`r=0,`$ and $`n=`$ $`82`$. The upper part of the figure shows numerical streamlines for a corner of the square. Although the state is spatially not well localized, it does have a strongly ‘paramagnetic’ current structure. By paramagnetic we mean $`dE_m/d\varphi <0`$.
The lower part shows the theoretical current density $`j_y(x)`$ for $`y=0,x[\frac{1}{2},0],`$ that is, the current density along the upper edge of the upper figure. A simple approximate formula for $`j_y`$ is $`j_y(x)\left(\mathrm{cos}n\pi x/2\right)^2\left[u_m\left(x\frac{1}{2}\right)^2u_m\left(x\frac{1}{2}\right)^2\right]`$ which gives double zeroes of the current at equally spaced $`x=(2l+1)/n.`$ In this case $`j_y`$ is negative and the factor depending on $`u_m`$ has no zeroes except at $`x=0,`$ for $`m=0.`$
The maxima are at $`j_y=0`$; $`j_y0.`$ Theory and numerics closely agree on the period and shape of the oscillations, but the lower envelope of the two differ. We mark with horizontal bars the lower envelope of the numerical calculation, which is considerably more irregular than given by our theory. This structure of zeroes of $`j_x`$ is a consequence of a symmetry, namely $`yy`$, together with complex conjugation. This means that $`\mathrm{\Psi }(x,0)`$ is real and therefore generically will have zeroes as a function of $`x.`$
## VIII Summary
### A Extensions of the results
We have shown how to classify and find eigenstates of a charged particle in a square billiard subjected to a magnetic flux which is classically weak. It is not necessarily weak quantally, however, and large remarkable changes in the wave functions are found which sometimes significantly localize the wave function. We used two basic flux configurations: a uniform field and an Aharonov-Bohm flux line. The placement of the ABFL is important. Of lesser importance is any finite radius to the line, at least in the range of parameters we use. We exhibited some of the wavefunctions in several forms and sequences, concentrating on wavefunctions connected with the $`(1,1)`$ resonance, that is, connected with periodic orbits of the flux free square whose velocities make angles of 45 with the coordinate axes.
These results can be readily generalized to higher order resonances. For a given energy the effects diminish quite rapidly as the order increases, but in principle, at sufficiently high energy any given resonance can show strong magnetic effects. We give a few results in Appendix B.
One can extend these results to integrable systems other than the square billiard, and to other flux configurations. There are other billiard shapes, such as the rectangle, certain triangles, the circle and the ellipse. One may also study ‘soft billiards’, e.g. a confining potential of the form $`U(x,y)=U_1(x)+U_2(y)`$ where $`U_i`$ is some sort of anharmonic potential. Soft billiards are more accurate representations of mesoscopic systems than are hard wall billiards, but the additional theoretical effort they require is not usually made. A practical difficulty, although not one of principle, is that it may be necessary to resort to action-angle variables, and it may be tedious to find the periodic orbits of the flux-free system, if some tiling trick cannot be used.
In any case, the first steps of the program are to choose a convenient surface of section, choose the resonance of interest, find the AB phase $`\mathrm{\Phi }(x,x^{}),`$ and the effective potential $`V(x).`$ Much insight can be gained at this level.
A much larger class of systems solvable by this technique are nearly integrable systems subjected to a flux. For example, one could start from a nearly square trapezoid. Then one would have a double perturbation of the square, one from the flux, the other from the change of shape. Each of these perturbations can have big quantum effects, and the combination of the two can be quite different from each one separately, especially since one breaks time reversal invariance and the other doesn’t. Again, finding the effective potential is key to understanding the qualitative results.
### B Possible experiments
There are possible experiments and even interesting devices which might be made, if time reversal can be broken. Having unusual wavefunctions suggests some of the possibilities. Probes of the system will depend on whether the wavefunction is large or small at the position of the probe.
Mesoscopic systems are of great current interest. In this case, electrons are the particles and there is a magnetic flux. The system of Lévy et al. consisting of isolated metallic ‘squares’ is a case of this type, but one can imagine leads weakly connected to a metallic square whose coupling depends on the shape of the wavefunction at the Fermi level and the position of the lead.
Another type of mesoscopic system is formed of surface electrons in ‘corrals’ on a metallic surface. The wave functions of such electrons can be probed with atomic accuracy with a scanning tunnelling microscope. Achieving an appropriate parameter range will be difficult, but perhaps not impossible. The corrals are leaky and do not really confine the electrons to their interior, but in many cases that idea seems to work, at least qualitatively.
Another kind of system is the shallow square or rectangular container of liquid, which is vibrated to set up standing waves. Instead of a flux, the system can be rotated with a uniform angular velocity. This differs from the uniform magnetic field in that the $`A^2`$ is absent from the Hamiltonian, but since under our assumptions that term can be neglected anyway, the same kind of results are expected. One can also introduce a nonideal ABFL. This is done by making a small hole in the tank and allowing the water to flow out with a certain vorticity. Experiments on such a system have been performed, although scattering rather than eigenstates was studied. Such an experiment would certainly have pedagogical value, and if the experiment is done, it may in fact be the first detailed observation of a nontrivial persistent current state.
Still another system is the thin square microwave cavity. Here the ‘quantum’ waves are microwaves. There is no quantity directly equivalent to the magnetic field. However, one can replace part of the boundary by a ferrite strip, say between $`[a,a]`$ on one side of the cavity. The phase shift of the microwave upon reflection from the ferrite depends on the direction of magnetization of the ferrite and the direction of incidence of the microwave. This gives as an effective potential exactly that of Eq. (31). There is then a localized eigenstate circulating the cavity in one direction, but not the other. A similar situation from the point of view of diffraction and effective potential, but maintaining time reversal invariance, is the ‘step’ billiard. This is a square billiard of side $`L`$ say, with one side moved down by $`\delta L=ϵL`$ for $`\left|x\right|<a.`$ The perturbed $`(1,0)`$ resonance states ‘see’ a square well effective potential to first approximation. If $`k\delta L`$ is small, but $`k\sqrt{L\delta L}1,`$ our theory works, gives nontrivial localization, and is similar to the flux line with small $`\varphi /\varphi _0.`$ The range $`k\delta L=\pi ,`$ i.e. $`\delta L`$ a half wavelength, gives no phase shift for normal incidence and is similar to an ABFL with $`\varphi /\varphi _0=1.`$ Experiments on this system could be carried out in several contexts.
### C Conclusions
We have solved a characteristic example of quantum states in a weakly perturbed integrable system. The states and their energies can be classified and found to good approximation. The wavefunctions are quite nontrivial and interesting, as compared with the states of the integrable system. The states of a hard chaotic system are also not individually interesting, except for some weak ‘scars’, requiring resort to statistical studies and averages over many states in such systems.
Our technique applies to a very large class of systems, which includes experimental systems. The particular case we emphasize breaks time reversal invariance, and the eigenstates have persistent currents. As far as we know, these are the first published nontrivial examples such states.
## IX Acknowledgments
Supported in part by the United States NSF grant DMR-9625549 and United States-Israel Binational Science Foundation, grant 99800319. R.N. was partially supported by the NSF grant DMR98-70681 and the University of Kentucky. We thank Prof. Director Peter Fulde for hospitality at the Max-Planck-Institut für Physik komplexer Systeme in Dresden, where some of this work was done. R.N. and O.Z. thank Dr. R. Seiler for hospitality at the SFB 288 “Differentialgeometrie und Quantenphysik”, TU Berlin. O.Z. thanks Dr. F. Haake for hospitality at the Universität GH Essen.
## A Two dimensional wavefunctions
In this Appendix we find the two dimensional wavefunction. If a solution to $`\psi =T\psi `$ is known, the wavefunction in the square can be found by evaluating the integral over the surface of section.
$$\mathrm{\Psi }(𝐫)=G(𝐫,x^{};E)\psi (x^{})𝑑x^{}.$$
(A1)
In this case
$`G(𝐫,x^{};E)`$ $`=`$ $`{\displaystyle \underset{\text{cl. tr.}}{}}{\displaystyle \frac{1}{2}}\left|{\displaystyle \frac{1}{2\pi k}}{\displaystyle \frac{^2S}{x^{}\mathrm{}_{}}}\right|^{1/2}`$ (A3)
$`\times \mathrm{exp}\left[iS(𝐫,x^{};E)i{\displaystyle \frac{\pi }{2}}\nu \right]`$
in our units. Here the sum is over the (unperturbed) classical trajectories going from a point $`x^{}`$ on the surface of section to a point $`𝐫`$ inside the square, $`S(𝐫,x^{};E)`$ is the reduced action for this trajectory, $`\mathrm{}_{}`$ is the direction perpendicular to the trajectory at point $`𝐫`$, and $`\nu `$ is the Maslov index. In a single square scheme we take the surface of section to be the lower side of the square $`1/2x^{}1/2`$, $`y^{}=1/2`$. Then $`\psi (x^{})`$ is a linear combination of functions $`\psi _I`$ and $`\psi _{II}`$ described in the text such that $`\psi (\pm 1/2)=0`$ (there should be no flow through the ends). It follows then that $`u_m(x^{}+2)=(1)^ru_m(x^{})`$ and
$$\psi (x^{})=e^{i\kappa x^{}}u_m(x^{})+i^re^{i\kappa x^{}}u_m(x^{}1)$$
(A4)
with $`r`$ defined by Eq. (15). Here we have also used a zero field quantization condition $`\kappa =\pi n/2`$ and the property $`u_m(x^{})=(1)^mu_m(x^{})`$. Note that in the surface of section picture $`r`$ does not have a direct interpretation as a representation label.
We evaluate the integral (A1) in the stationary phase approximation. As we shall see, only the 45 orbits survive. For a given point $`𝐫=(x,y)`$ there are four such orbits: two starting at $`x_1^{}`$ and two at $`x_2^{}`$ (Fig. 14). Each orbit gives one term in Eq. (14). Note that the orbits with the positive (negative) $`x`$-projection of the momentum at $`x^{}`$ are generated by the first (second) term in Eq. (A4).
Although the following calculation depends in which sector inside the square the point lies, the final result is independent of sector. So for definitiveness we assume $`y<x<y`$, that is the point lies below the diagonals $`y=x`$ and $`y=x`$. We have numbered the orbits by $`s=0,\mathrm{},3`$ according to the terms in Eq. (14) they represent. For $`s=0`$ orbit, expressing the action in terms of the distance, $`S(𝐫,x^{};E)=kL(𝐫,x^{})k\sqrt{(xx^{})^2+(y+1/2)^2}`$, we find the prefactor in Eq. (A3)
$$\left|\frac{1}{8\pi k}\frac{^2S}{x^{}\mathrm{}_{}}\right|^{1/2}=\left(\frac{1}{8\sqrt{2}\pi L}\right)^{1/2}.$$
(A5)
The stationary point in the integral (A1) is determined by the exponents in Eq. (A3) and in the first term of Eq. (A4). It is $`x_1^{}=xy1/2`$, as in a 45 orbit. The $`s=0`$ wavefunction is then
$$\mathrm{\Psi }_0(x,y)=e^{i\kappa (x+y)}u_m\left(xy\frac{1}{2}\right)$$
(A6)
dropping a constant factor $`\kappa ^{1/2}e^{i(\kappa /2\pi /4)}/2`$. The terms with $`s=1,2,3`$ can be obtained in the same manner with an appropriate addition of the Maslov phase for each bounce.
If $`𝐫`$ is located in a different sector of the square, the results remain the same if proper account is made of the Maslov phases. The labeling of the orbits is invariant if done by the rule: $`s=0`$ trajectory arrives to $`𝐫`$ from the South-West, $`s=1`$ from SE, $`s=2`$ from NE, and $`s=3`$ from NW.
It is worth noticing that Eqs. (A6) or (14) are valid only up to the square root order of the small parameter $`ϵ=B/k`$ of the perturbation theory behind our work. If we wish to obtain the results valid to the first order of $`ϵ`$ we should (a) use a better approximation for $`u_m(x^{})`$, (b) add the vector potential term to the action $`S(𝐫,x^{};E)`$, and (c) find a correction to the stationary point due to $`u_m(x^{})`$.
Alternatively, the above calculation can be carried out in the repeated square scheme defined in Sec. III. The surface of section is then the line $`y=1/2`$ (identified with $`y=3/2`$). We can restrict ourselves to the trajectories with the positive $`x`$\- and $`y`$-projections of the momentum in this extended manifold. Then only $`\psi _I(x^{})`$ is needed to generate a complete two-dimensional wavefunction. In order for $`\psi _I`$ to have a period 2, the condition $`u_m(x^{}+2)=(1)^ru_m(x^{})`$ must be satisfied. A point $`𝐫`$ in the original square will be represented by its images in four (up to a translation by $`\mathrm{\Delta }x=2`$) domains in the manifold. Each of these images is connected by one 45 orbit to the surface of section. These four orbits generate the four terms $`s=0,\mathrm{},3`$ in Eq. (14) as follows: $`s=0`$ is produced by the point $`𝐫`$ in the original domain, $`s=1,3,2`$ by its reflection about the right side of the square, upper side of the square, and the composition of both, respectively. There are no Maslov phases in Eq. (A3) in this case. However the terms $`s=1,3`$, which are obtained by the odd number of reflections, should be summed with an additional minus sign. Indeed, the Jacobian of the coordinate transformation from the original square to the extended manifold is singular on the boundary. So, when the wavefunction on the manifold is folded back to the physical domain a phase difference $`\pi `$ will be accumulated between each pair of domains related by one reflection. Of course, this phase is analogous to the Maslov phase for an impenetrable wall.
## B Higher resonances
For completeness, we give some results for higher resonances, which can be labelled with relatively prime integers, $`p,q.`$ These resonances are closely related to states of a rectangle of side $`1/p,`$ $`1/q.`$ Obviously, for the square, the energies for $`p,q`$ are the same as those for $`q,p.`$ The corresponding states are also the same after a $`90^{}`$ rotation. However, for $`\left|pq\right|/\left|p+q\right|`$ not too small, the $`p,q`$ states and the $`q,p`$ states are different and are nearly uncoupled and there is a nearly doubly degenerate set of energy levels. The splitting of the exact levels, which are combinations of $`pq`$ and $`qp`$ that are eigenstates of the rotation operator $``$, can be estimated by using an analog of ‘chaos assisted tunnelling’.
The classical periodic orbits make $`q`$bounces from the $`x`$ sides and $`p`$ bounces from the $`y`$ sides. These resonant orbits all have the same length $`_{pq}=2\sqrt{p^2+q^2}.`$ The maximal directed area enclosed by an orbit is $`\pm 1/2pq.`$ The energy is $`k_{nm}^2`$ where
$$k_{nm}^2\left(\frac{2\pi n}{_{pq}}\right)^2+\left[1+\frac{p^2}{q^2}\right]E_m^{(q)}$$
(B1)
and $`n`$ is an integer. The first term is much larger than the second. Again, we have the approximate quantization of $`n`$ wavelengths in an orbit length.
The potential $`V_{pq}(x)`$ depends on $`q`$ being associated with the $`x`$-direction, thus the notation $`E_m^{(q)}`$. This ‘transverse’ contribution to the energy is of course symmetric in $`p`$ and $`q,`$ that is, $`q^2E_m^{(p)}=p^2E_m^{(q)}.`$ The potential is
$$V_{pq}(x)=\frac{q}{p}\left[\frac{_{11}}{_{pq}}\right]^3V\left[q\left(x+\frac{1}{2}\right)\frac{1}{2}\right]$$
(B2)
where $`V`$ is given by Eq. (11). $`V_{pq}`$ has period $`2/q`$, rather than $`2.`$ The boundary condition on the eigenstates is $`u(x2)=e^{i\beta }u(x)`$ where $`\beta =2\pi frac\left[pn/\left(p^2+q^2\right)\right]`$ and $`frac`$ indicates the fractional part. The potential can be regarded as weak if $`kB<<pq_{pq}^3.`$ If this condition is satisfied the $`pq`$ resonance can be ignored.
Similar but somewhat more complex results can be obtained for rectangles and certain triangles.
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# Irregular Input Data in Convergence Acceleration and Summation Processes: General Considerations and Some Special Gaussian Hypergeometric Series as Model Problems
## I Introduction
In mathematics and in the mathematical treatment of scientific problems, slowly convergent or divergent sequences and series occur abundantly. Accordingly, many techniques for the acceleration of convergence and the summation of divergent series have been invented, and some of them are even older than calculus (see for instance pp. 90 - 91 of or p. 249 of ).
Sequence transformations are principal tools to overcome convergence problems. Let us assume that $`\{s_n\}_{n=0}^{\mathrm{}}`$ is a slowly convergent or divergent sequence, whose elements $`s_n`$ may for example be the partial sums of an infinite series:
$$s_n=\underset{k=0}{\overset{n}{}}a_k.$$
(1)
The basic assumption of all sequence transformations is that a sequence element $`s_n`$ can for all indices $`n0`$ be partitioned into a (generalized) limit $`s`$ and a remainder or truncation error $`r_n`$ according to
$$s_n=s+r_n.$$
(2)
The conventional approach of evaluating an infinite series consists in adding up so many terms that the remainders $`r_n`$ ultimately become negligible. Unfortunately, this is not always feasible because of obvious practical limitations. Moreover, adding up further terms does not work in the case of a divergent series since their terms usually increase in magnitude with increasing index.
Alternatively, one could try to determine approximations to the remainders $`r_n`$ and to eliminate them from the sequence elements $`s_n`$. At least conceptually, this is what a sequence transformation tries to accomplish. Thus, the original sequence $`\{s_n\}_{n=0}^{\mathrm{}}`$ is transformed into a new sequence $`\{s_n^{}\}_{n=0}^{\mathrm{}}`$ whose elements have the same (generalized) limit $`s`$ but different remainders $`r_n^{}`$:
$$s_n^{}=s+r_n^{}.$$
(3)
The transformation process was successful if the transformed remainders $`r_n^{}`$ have superior numerical properties. For example, in the literature on extrapolation methods it is said that a sequence transformation accelerates convergence if the transformed remainders $`r_n^{}`$ vanish more rapidly than the original remainders $`r_n`$ according to
$$\underset{n\mathrm{}}{lim}\frac{r_n^{}}{r_n}=\underset{n\mathrm{}}{lim}\frac{s_n^{}s}{s_ns}=\mathrm{\hspace{0.33em}0}.$$
(4)
Similarly, a divergent sequence $`\{s_n\}_{n=0}^{\mathrm{}}`$, whose remainders $`r_n`$ do not vanish as $`n\mathrm{}`$, is transformed into convergent sequence $`\{s_n^{}\}_{n=0}^{\mathrm{}}`$ if the transformed remainders $`r_n^{}`$ vanish as $`n\mathrm{}`$.
During the last years, considerable progress has been reached in this field, as documented by the large number of recent monographs and review articles . Moreover, numerous applications of sequence transformations have been reported in the literature. For example, the present author has applied sequence transformations successfully in such diverse fields as the evaluation of special functions , the evaluation of molecular multicenter integrals of exponentially decaying functions , the summation of strongly divergent quantum mechanical perturbation expansions , and the extrapolation of crystal orbital and cluster calculations for oligomers to their infinite chain limits of stereoregular quasi-onedimensional organic polymers .
It should be noted that Padé approximants , which in applied mathematics and in theoretical physics have become the standard tool to overcome convergence problems with power series, can be considered to be a special class of sequence transformations since the partial sums of a power series are transformed into a doubly indexed sequence of rational functions.
As described above, sequence transformations try at least in principle to construct approximations to the actual remainders which are then eliminated from the input data. This is done by detecting and utilizing regularities in the behavior of the elements of the sequence to be transformed. For sufficiently large indices $`n`$, one can expect that certain asymptotic regularities do exist. However, sequence transformations are normally used with the intention of avoiding the asymptotic domain, i.e., the transforms are constructed from the *leading* elements of the input sequence. Unfortunately, sequence elements $`s_n`$ with small indices $`n`$ often behave irregularly. In such a case, a straightforward application of a sequence transformation may be ineffective and even lead to completely nonsensical results. Instead, one should analyze the behavior of the input data as a function of the index and exclude highly irregular sequence elements from the transformation process if necessary. In this way, one has a much better chance of obtaining good and reliable transformation results. It is the intention of this article to describe and classify some of the problems, which can result from irregular input data, and to discuss strategies to overcome them.
Sequence transformations are needed most in cases in which apart from the numerical values of a few elements of a slowly convergent or divergent sequence only very little is known. This is a situation which is not uncommon in scientific applications as for example the summation of strongly divergent perturbation expansions as they occur quantum mechanics or in quantum field theory. Thus, it would in principle be desirable to discuss problems with irregular input data also via examples of numerically determined input data. However, the lack of detailed knowledge about the behavior of the elements of such a sequence makes it hard to fully understand the numerical problems as well as to develop strategies to overcome them. Consequently, complications of that kind are discussed in this article predominantly via suitable mathematical model problems.
In Section II, some formal aspect of sequence transformations are discussed, in particular the concept of a path in the table of a transformation. In this way, different approaches for the determination of the approximations to the (generalized) limit of the input sequence can be classified and formalized. Moreover, the concept of a path helps to understand the impact of irregular input data on the performance of sequence transformations.
At least some of the problems mentioned above can be illuminated by considering the Gaussian hypergeometric function $`{}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ which is defined by a power series that converges in the interior of the unit circle. This function does not only depend on an argument $`z`$ but also on three essentially arbitrary parameters $`a`$, $`b`$, and $`c`$. As discussed in Section III, the convergence of the hypergeometric series can for most values of the parameters $`a`$, $`b`$, and $`c`$ be accelerated quite effectively by a variety of different sequence transformations. Moreover, it is in this way frequently possible to associate a finite value to a hypergeometric series even if its argument does not lie in the interior of the unit circle.
However, as discussed in Section IV, the situation changes dramatically if the third parameter $`c`$ of the hypergeometric series $`{}_{2}{}^{}F_{1}^{}`$ is a negative real number. Then, the terms of this series first increase with increasing summation index even for $`|z|<1`$ and produce partial sums which look like the elements of a mildly divergent sequence. Only for sufficiently large indices, the terms decrease in magnitude and ultimately produce a convergent result. Accordingly, the leading partial sums of such a hypergeometric series display a highly irregular behavior, and they should not be used as input data for a sequence transformation. If the leading irregular coefficients are skipped and only regular coefficients with higher indices are used as input data, then sequence transformations are again able to produce good and reliable results.
Section V contains a summary. In Appendix A, the properties of the sequence transformations, which are used in this article, are discussed. Problems with irregular input data occur also quite frequently in the mathematical treatment of scientific problems. In Appendix B, it is shown that the convergence of extensive summation calculations for the so-called infinite coupling limit $`k_3`$ of the sextic anharmonic oscillator can be improved considerably by excluding the leading irregular coefficients of the divergent perturbation series from the transformation process. Finally, numerous new recurrence formulas for the hypergeometric function $`{}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ are derived in Appendix C.
## II Order-Constant and Index-Constant Paths
In this Section, some aspects of sequence transformations are discussed which admittedly look very formal. Nevertheless, they should not be ignored since they may be very consequential in practical applications.
Obviously, a computational algorithm can only involve a finite number of arithmetic operations. Consequently, a sequence transformation $`𝒯`$ can only use finite subsets of the original sequence $`\{s_n\}_{n=0}^{\mathrm{}}`$ for the computation of new sequence elements $`s_m^{}`$. In addition, these finite subsets normally consist of consecutive elements. Accordingly, only subsets of the type $`\{s_n,s_{n+1},\mathrm{},s_{n+l}\}`$ will be considered in this article.
All the commonly used sequence transformations $`𝒯`$ can be represented by infinite sets of doubly indexed quantities $`T_k^{(n)}`$ with $`k,n0`$ that can be displayed in a two-dimensional array which is called the table of $`𝒯`$.
Here, the convention is used that the superscript $`n`$ always indicates the minimal index occurring in the finite subset of sequence elements used for the computation of a given $`T_k^{(n)}`$. The subscript $`k`$ – usually called the order of the transformation – is a measure for the complexity of the transformation process which yields $`T_k^{(n)}`$.
The elements $`T_k^{(n)}`$ of the table of $`𝒯`$ are gauged in such a way that $`T_0^{(n)}`$ corresponds to an untransformed sequence element,
$$T_0^{(n)}=s_n.$$
(5)
An increasing value of $`k`$ implies that the complexity of the transformation process increases. Moreover, $`l=l(k)`$ also increases. This means that for every $`k,n0`$ the sequence transformation $`𝒯`$ produces a new transform according to
$$T_k^{(n)}=𝒯(s_n,s_{n+1},\mathrm{},s_{n+l(k)}).$$
(6)
The exact relationship, which connects $`k`$ and $`l`$, is specific for a given sequence transformation $`𝒯`$.
Let us assume that a sequence transformation $`𝒯`$ should be used to speed up the convergence of some sequence $`\{s_n\}_{n=0}^{\mathrm{}}`$ to its limit $`s=s_{\mathrm{}}`$. One can try to obtain a better approximation to $`s`$ by proceeding on an in principle unlimited variety of different paths in the table of $`𝒯`$. Two extreme types of paths – and also those which are predominantly used in practical applications – are order-constant paths
$$T_k^{(n)},T_k^{(n+1)},T_k^{(n+2)},\mathrm{}$$
(7)
with fixed transformation order $`k`$ and $`n\mathrm{}`$, and index-constant paths
$$T_k^{(n)},T_{k+1}^{(n)},T_{k+2}^{(n)},\mathrm{}$$
(8)
with fixed minimal index $`n`$ and $`k\mathrm{}`$.
Order-constant and index-constant paths differ significantly. It is not even a priori clear that these two types of paths lead to the same limit in the case of an arbitrary sequence $`\{s_n\}_{n=0}^{\mathrm{}}`$. However, for the sake of simplicity this potential complication will be ignored here, and we shall always tacitly assume that order-constant and index-constant paths lead to the same limit.
In the case of an order-constant path, a fixed number of $`l+1`$ sequence elements $`\{s_n,s_{n+1},\mathrm{}s_{n+l}\}`$ is used for the computation of $`T_k^{(n)}`$, and the starting index $`n`$ of this string of fixed length is increased successively until either convergence is achieved or the number of available elements of the input sequence is exhausted.
In the case of an index-constant path, the starting index $`n`$ is kept fixed at a low value (usually $`n=0`$ or $`n=1`$) and the transformation order $`k`$ is increased and with it the number of elements contained in the subset $`\{s_n,s_{n+1},\mathrm{}s_{n+l(k)}\}`$. Thus, on an index-constant path it is always tried to compute from a given set of input data that element $`T_k^{(n)}`$ which has the highest possible transformation order $`k`$.
In order to clarify the differences between order-constant and index-constant paths, let us consider the probably best known sequence transformation, Wynn’s epsilon algorithm :
$`ϵ_1^{(n)}`$ $`=`$ $`0,ϵ_0^{(n)}=s_n,`$ (10)
$`ϵ_{k+1}^{(n)}`$ $`=`$ $`ϵ_{k1}^{(n+1)}+\mathrm{\hspace{0.17em}1}/[ϵ_k^{(n+1)}ϵ_k^{(n)}].`$ (11)
Wynn showed that if the input data $`s_n`$ for the epsilon algorithm are the partial sums
$$f_n(z)=\underset{\nu =0}{\overset{n}{}}\gamma _\nu z^\nu $$
(12)
of a (formal) power series for some function $`f(z)`$, then the elements $`ϵ_{2k}^{(n)}`$ with even subscripts are Padé approximants to $`f`$ according to
$$ϵ_{2k}^{(n)}=[n+k/k].$$
(13)
Here, the notation of the monograph by Baker and Graves-Morris is used, i.e., a Padé approximant $`[l/m]`$ corresponds to the ratio of two polynomials $`P_l(z)`$ and $`Q_m(z)`$, which are of degrees $`l`$ and $`m`$, respectively, in $`z`$. In contrast, the elements $`ϵ_{2k+1}^{(n)}`$ with odd subscripts are only auxiliary quantities which diverge if the whole process converges.
It follows from (13) that the epsilon algorithm (II) effects the following transformation of the partial sums (12) to Padé approximants:
$$\{f_n(z),f_{n+1}(z),\mathrm{},f_{n+2k}(z)\}[n+k/k].$$
(14)
Thus, if we use a window consisting of $`2k+1`$ partial sums $`f_{n+j}(z)`$ with $`0j2k`$ on an order-constant path and increase the minimal index $`n`$ successively, the epsilon algorithm produces the following sequence of Padé approximants:
$$[n+k/k],[n+k+1/k],\mathrm{},[n+k+m/k],\mathrm{}.$$
(15)
Only $`2k+1`$ partial sums are used for the computation of the Padé approximants, although many more are known. Obviously, the available information is not exploited optimally on such an order-constant path.
Moreover, the degree of the numerator polynomial of a Padé approximant $`[n+k+m/k]`$ increases with increasing $`m0`$, whereas the degree of the denominator polynomial remains fixed. Thus, these Padé look unbalanced. Instead, it seems to be much more natural to use diagonal Padé approximants, i.e., Padé approximants with numerator and denominator polynomials of equal degree, or – if this is not possible – to use Padé approximants with degrees of the numerator and denominator polynomials that differ as little as possible.
This approach has in principle many theoretical as well as practical advantages. For example, Wynn could show that if the partial sums $`f_0(z)`$, $`f_1(z)`$, $`\mathrm{}`$, $`f_{2n}(z)`$ of a Stieltjes series are used for the computation of Padé approximants, then the diagonal approximant $`[n/n]`$ provides the most accurate approximation to the corresponding Stieltjes function $`f(z)`$, and if the partial sums $`f_0(z)`$, $`f_1(z)`$, $`\mathrm{}`$, $`f_{2n+1}(z)`$ are used for the computation of Padé approximants, then either $`[n+1/n]`$ or $`[n/n+1]`$ provides the most accurate approximation (Theorem 5 of ). A detailed discussion of Stieltjes series and their special role in the theory of Padé approximants can for instance be found in Section 5 of the monograph by Baker and Graves-Morris .
Thus, it is apparently an obvious idea to try to use either diagonal Padé approximants or their closest neighbors whenever possible. Let us assume that the partial sums $`f_0(z)`$, $`f_1(z)`$, $`\mathrm{}`$, $`f_m(z)`$ are known. If $`m`$ is even or odd, $`m=2\mu `$ or $`m=2\mu +1`$, respectively, the elements of the epsilon table with the highest possible transformation orders are given by the transformations
$`\{f_0(z),f_1(z),\mathrm{},f_{2\mu }(z)\}`$ $``$ $`ϵ_{2\mu }^{(0)}=[\mu /\mu ],`$ (16)
$`\{f_1(z),f_2(z),\mathrm{},f_{2\mu +1}(z)\}`$ $``$ $`ϵ_{2\mu }^{(1)}=[\mu +1/\mu ].`$ (17)
With the help of the notation $`[[x]]`$ for the integral part of $`x`$, which is the largest integer $`\nu `$ satisfying $`\nu x`$, these two relationships can be expressed by a single equation (Eq. (4.3-6) of ):
$`\{f_{m2[[m/2]]}(z),f_{m2[[m/2]]+1}(z),\mathrm{},f_m(z)\}`$ (18)
$``$ $`ϵ_{2[[m/2]]}^{(m2[[m/2]])}=\left[m[[m/2]]/[[m/2]]\right].`$ (19)
For $`m=0,1,2,\mathrm{}`$, these transformations correspond to the following staircase sequence in the Padé table (Eq. (4.3-7) of ):
$`[0/0],[1/0],[1/1],\mathrm{}`$ (21)
$`\mathrm{},[\nu /\nu ],[\nu +1/\nu ],[\nu +1/\nu +1],\mathrm{}.`$
This staircase sequence exploits the available information optimally if the partial sums $`f_m(z)`$ with $`m0`$ are computed successively and if after the computation of each new partial sum the element of the epsilon table with the highest possible even transformation order is computed. Moreover, the Padé approximants obtained in this way look balanced since the degrees of their numerator and denominator polynomials differ as little as possible.
The example of Wynn’s epsilon algorithm strongly indicates that index-constant paths are at least in principle computationally more efficient than order-constant paths since they exploit the available information optimally. This is in general also true for all other sequence transformations considered in this article.
Another serious disadvantage of order-constant paths is that they cannot be used for the summation of divergent sequences and series since increasing $`n`$ in the set $`\{s_n,s_{n+1},\mathrm{},s_{n+l}\}`$ of input data normally only increases divergence.
In view of the examples given above, it is apparently an obvious idea to use exclusively index-constant paths, and preferably those which start at a very low index $`n`$, for instance at $`n=0`$ or $`n=1`$. This is certainly a good idea if all elements of the input sequence contain roughly the same amount of useful information. If, however, the leading terms of the sequence to be transformed behave irregularly, they cannot contribute useful information, or – to make things worse – they contribute *wrong* information. In such a case it is usually necessary to exclude the leading elements of the input sequence from the transformation process. Thus, it is preferable to use either an order-constant path or an index-constant path with a sufficiently large starting index $`n`$. The use of an order-constant path has the additional advantage that the diminishing influence of irregular input data with small indices $`n`$ should become obvious from the transformation results as $`n`$ increases.
Finally, an important theoretical advantage of order-constant paths should also be mentioned. Normally, it is much easier to perform a theoretical convergence analysis of a sequence transformation on an order-constant path than on an index-constant path. As the starting index $`n`$ of the string $`\{s_n,s_{n+1},\mathrm{},s_{n+l}\}`$ of input data becomes large, asymptotic approximations to the sequence elements $`s_n`$ can be used. Often, this greatly simplifies a theoretical analysis. In the case of index-constant paths, such a simplification is not possible because not all input data have large indices. Accordingly, theoretical convergence properties of sequence transformations are studied almost exclusively on order-constant paths. Notable exceptions are two articles by Sidi where the convergence properties of sequence transformations on both order-constant and index-constant paths are analyzed.
## III The Gaussian Hypergeometric Function
The Gaussian hypergeometric function $`{}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ is one of the most important special functions of mathematical physics, and its properties are discussed in numerous books, for example in those by Abramowitz and Stegun , Erdélyi, Magnus, Oberhettinger, and Tricomi , Magnus, Oberhettinger, and Soni , Seaborn , Slater , Spanier and Oldham , Temme , and Wang and Guo .
The hypergeometric function $`{}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ is defined via the corresponding hypergeometric series (p. 37 of )
$${}_{2}{}^{}F_{1}^{}(a,b;c;z)=\underset{m=0}{\overset{\mathrm{}}{}}\frac{(a)_m(b)_m}{(c)_mm!}z^m,$$
(22)
where $`(a)_m=\mathrm{\Gamma }(a+m)/\mathrm{\Gamma }(a)`$ is a Pochhammer symbol (see for example p. 3 of ). The series (22) terminates after a finite number of terms if either $`a`$ or $`b`$ is a negative integer. Otherwise, it converges in the interior of the unit circle, i.e., for $`|z|<1`$, and it diverges for $`|z|>1`$. On the boundary $`|z|=1`$ of the unit circle, the series (22) diverges if $`\mathrm{Re}(a+bc)1`$, it converges *absolutely* for $`\mathrm{Re}(a+bc)<0`$, and it converges *conditionally* for $`0\mathrm{Re}(a+bc)<1`$, the point $`z=1`$ being excluded.
Thus, a nonterminating hypergeometric series does not suffice for the computation of the corresponding hypergeometric function $`{}_{2}{}^{}F_{1}^{}(a,b;c;z)`$, which is in general a multivalued function defined in the whole complex plane with branch points at $`z=1`$ and $`z=\mathrm{}`$. Instead, techniques which permit an analytic continuation from the interior to the exterior of the unit circle are needed.
Sequence transformations can be used to accelerate the convergence of a hypergeometric series $`{}_{2}{}^{}F_{1}^{}`$ or to sum it in the case of divergence, which corresponds to an analytic continuation. Let us for example consider the following elementary special case of a hypergeometric function $`{}_{2}{}^{}F_{1}^{}`$ (Eq. (15.1.3) of ):
$`\mathrm{ln}(1+z)`$ $`=`$ $`{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^mz^{m+1}}{m+1}}`$ (23)
$`=`$ $`z{}_{2}{}^{}F_{1}^{}(1,1;2;z).`$ (24)
The infinite series converges only for $`|z|<1`$, whereas the logarithm is with the exception of the cut along $`\mathrm{}<z1`$ defined in the whole complex plane.
In Table I, Wynn’s epsilon algorithm, Eq. (II), Brezinski’s theta algorithm, Eq. (A), Levin’s transformation $`d_k^{(n)}(\zeta ,s_n)`$, Eq. (A36), and the closely related sequence transformation $`\delta _k^{(n)}(\zeta ,s_n)`$, Eq. (A40), are applied to the partial sums
$$s_n(z)=\underset{m=0}{\overset{n}{}}\frac{(1)^mz^{m+1}}{m+1}$$
(25)
of the hypergeometric series (24) for $`z=7/2`$. The approximations to the limit in Table I were always chosen in such a way that the transforms with the highest possible transformation order were computed from a given set of input data. Thus, in the case of the epsilon algorithm, the approximations to the limit were chosen according to (19), in the case of the theta algorithm, they were chosen according to (A21), and in the case of the sequence transformations $`d_k^{(n)}(\zeta ,s_n)`$ and $`\delta _k^{(n)}(\zeta ,s_n)`$, they were chosen according to (A41).
The second column of Table I, which displays the partial sums (25), shows that the hypergeometric series (24) for $`\mathrm{ln}(1+z)`$ diverges quite strongly for $`z=7/2`$. Nevertheless, it is apparently possible to sum this divergent series to its correct value.
The results in Table I also show that Wynn’s epsilon algorithm, which in the case of a power series produces Padé approximants according to (13), is contrary to a widespread belief not necessarily the most powerful transformation. The theta algorithm and in particular the two Levin-type transformations $`d_k^{(n)}(\zeta ,s_n)`$ and $`\delta _k^{(n)}(\zeta ,s_n)`$, which use the first term neglected in the partial sum as a remainder estimate according to (A35), produce significantly better summation results.
The other sequence transformations discussed in Appendix A give better results than Wynn’s epsilon algorithm, but are less effective than the transformations $`d_k^{(n)}(\zeta ,s_n)`$ and $`\delta _k^{(n)}(\zeta ,s_n)`$. For example, the two approximations with the highest possible transformation orders, which Aitken’s iterated $`\mathrm{\Delta }^2`$ process (A) produces from the partial sums $`s_0(z)`$, $`s_1(z)`$, …, $`s_{15}(z)`$ according to (A11), are
$`𝒜_7^{(0)}`$ $`=`$ $`1.504077397173,`$ (26)
$`𝒜_7^{(1)}`$ $`=`$ $`1.504077396169.`$ (27)
Similarly, the iteration (A) of Brezinski’s theta algorithm produces according to (A32) the approximants
$`𝒥_4^{(2)}`$ $`=`$ $`1.504077404830,`$ (28)
$`𝒥_5^{(0)}`$ $`=`$ $`1.504077394094.`$ (29)
These results show that sequence transformations can be very useful. Nevertheless, in the case of a *real* real argument $`z`$ it is actually not necessary to to use sequence transformations for doing the analytic continuations or for speeding up convergence. Instead, it is often simpler to exploit some known mathematical properties: Unless certain linear combinations of the parameters $`a`$, $`b`$, and $`c`$ are positive or negative integers, a hypergeometric function $`{}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ can be expressed as the sum of two other $`{}_{2}{}^{}F_{1}^{}`$’s with a transformed argument $`w=1z`$, $`w=1/z`$, $`w=1/(1z)`$, or $`w=11/z`$, respectively. Thus, the argument $`w`$ of the two resulting hypergeometric series can normally be chosen in such a way that the two new series in $`w`$ either converge, if the original series in $`z`$ diverges, or that they converge more rapidly if the original series converges too slowly to be numerically useful.
For example, if $`|1z|<1`$ and if $`cab`$ is not a positive or negative integer, then we can use the analytic continuation formula (Eq. (15.3.6) of )
$`{}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ (30)
$`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }(c)\mathrm{\Gamma }(cab)}{\mathrm{\Gamma }(ca)\mathrm{\Gamma }(cb)}}{}_{2}{}^{}F_{1}^{}(a,b;a+bc+1;1z)`$ (33)
$`+{\displaystyle \frac{\mathrm{\Gamma }(c)\mathrm{\Gamma }(a+bc)}{\mathrm{\Gamma }(a)\mathrm{\Gamma }(b)}}(1z)^{cab}`$
$`\times {}_{2}{}^{}F_{1}^{}(ca,cb;cab+1;1z).`$
Let us now assume that $`z`$ is only slightly smaller than 1. Then, the convergence of the original hypergeometric series $`{}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ will be very bad. However, the two hypergeometric series on the right-hand side with argument $`1z`$ will converge rapidly in the vicinity of $`z=1`$.
With the help of this or similar analytic continuation formulas it is normally possible to compute a hypergeometric function $`{}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ with *real* argument $`z`$ effectively since it is possible to find for every argument $`z(\mathrm{},+\mathrm{})`$ two hypergeometric series with an argument $`|w|1/2`$ (see p. 127 of or Table I of ).
Unfortunately, this approach does not necessarily work in the case of complex arguments $`z`$. Consider the points
$$z_{1,2}=\frac{1\pm \sqrt{3}}{2},$$
(34)
which both lie on the boundary of the circle of convergence because of $`|z_{1,2}|=1`$. In practice, it is either impossible or not feasible to evaluate a nonterminating hypergeometric series $`{}_{2}{}^{}F_{1}^{}`$ by adding up its terms if its argument $`z`$ lies on the boundary of the unit circle. Consequently, something has to be done to speed up convergence or to accomplish a summation in the case of divergence. Unfortunately, the analytic continuation formulas of the type of (33) do not improve the situation if $`z=z_{1,2}`$ since
$`1z_{1,2}`$ $`=`$ $`z_{2,1},`$ (36)
$`1/z_{1,2}`$ $`=`$ $`z_{2,1},`$ (37)
$`1/(1z_{1,2})`$ $`=`$ $`z_{1,2},`$ (38)
$`11/z_{1,2}`$ $`=`$ $`z_{1,2}.`$ (39)
Hence, the analytic continuation formulas cannot help if the $`z`$ is close to $`z_{1,2}`$.
However, sequence transformations work for $`z=z_{1,2}`$. Let us consider the following elementary special case of a hypergeometric function $`{}_{2}{}^{}F_{1}^{}`$ (p. 38 of ):
$$(1+z)(1z)^{2\alpha 1}={}_{2}{}^{}F_{1}^{}(2\alpha ,\alpha +1;\alpha ;z).$$
(40)
For $`\alpha >0`$, the hypergeometric series converges in the interior of the unit circle, but it diverges on its boundary.
In Table II, Wynn’s epsilon algorithm, Eq. (II), and Levin’s transformation $`d_k^{(n)}(\zeta ,s_n)`$, Eq. (A36), are used to sum the hypergeometric series (40) with $`\alpha =1/3`$ on the boundary of the unit circle, i.e., they are applied to the partial sums
$$s_n(z)=\underset{m=0}{\overset{n}{}}\frac{(2/3)_m(4/3)_m}{(1/3)_mm!}z^m$$
(41)
with $`z=z_1=(1+\mathrm{i}\sqrt{3})/2`$.
The partial sums (41) in the second column of Table II display a very unusual sign pattern and grow slowly in magnitude. Nevertheless, both the epsilon algorithm as well as the Levin transformation are apparently able to sum the hypergeometric series (40) for $`z=(1+\mathrm{i}\sqrt{3})/2`$.
The other sequence transformations discussed in Appendix A are apparently also able to sum the hypergeometric series (40) for $`z=(1+\mathrm{i}\sqrt{3})/2`$. They give better results than Wynn’s epsilon algorithm, but are less effective than Levin’s transformations $`d_k^{(n)}(\zeta ,s_n)`$. The two approximations with the highest possible transformation orders, that can be produced from the partial sums $`s_0(z)`$, $`s_1(z)`$, …, $`s_{15}(z)`$ by Aitken’s iterated $`\mathrm{\Delta }^2`$ process (A), by Brezinski’s theta algorithm (A) and its iteration (A), and by the Levin-type transformation (A40), are
$`𝒜_7^{(0)}=`$ (43)
$`\mathrm{\hspace{0.17em}1.113}340798057+\mathrm{i}\mathrm{\hspace{0.17em}1.326}827896288,`$
$`𝒜_7^{(1)}=`$ (45)
$`\mathrm{\hspace{0.17em}1.113}340798408+\mathrm{i}\mathrm{\hspace{0.17em}1.326}827896424,`$
$`\theta _8^{(2)}=`$ (47)
$`\mathrm{\hspace{0.17em}1.113}340797528+\mathrm{i}\mathrm{\hspace{0.17em}1.326}827893689,`$
$`\theta _{10}^{(0)}=`$ (49)
$`\mathrm{\hspace{0.17em}1.113}340799160+\mathrm{i}\mathrm{\hspace{0.17em}1.326}827895539,`$
$`𝒥_4^{(2)}=`$ (51)
$`\mathrm{\hspace{0.17em}1.113}340798249+\mathrm{i}\mathrm{\hspace{0.17em}1.326}827894967,`$
$`𝒥_5^{(0)}=`$ (53)
$`\mathrm{\hspace{0.17em}1.113}340798314+\mathrm{i}\mathrm{\hspace{0.17em}1.326}827895955,`$
$`\delta _{14}^{(0)}(1,s_0(z))=`$ (55)
$`\mathrm{\hspace{0.17em}1.113}340798314+\mathrm{i}\mathrm{\hspace{0.17em}1.326}827895955,`$
$`\delta _{15}^{(0)}(1,s_0(z))=`$ (57)
$`\mathrm{\hspace{0.17em}1.113}340798414+\mathrm{i}\mathrm{\hspace{0.17em}1.326}827896325.`$
The numerical results presented here should suffice to support the claim of the author that sequence transformations can be extremely useful numerical tools for the evaluation of special functions in general and for the evaluation of Gaussian hypergeometric series $`{}_{2}{}^{}F_{1}^{}`$ in special. Moreover, the numerical results shown above indicate that a computational algorithm, which would be capable of evaluating of a hypergeometric function $`{}_{2}{}^{}F_{1}^{}`$ with essentially arbitrary *complex* argument $`z`$ and parameters $`a`$, $`b`$, and $`c`$, should be a suitable combination of analytic continuation formulas of the type of (33) with sequence transformations.
Here it should be taken into account that a hypergeometric series $`{}_{2}{}^{}F_{1}^{}`$ depends on four essentially arbitrary complex quantities $`a`$, $`b`$, $`c`$, and $`z`$. This makes the development of a general algorithm for its computation difficult since a very large variety of special cases and computationally different situations have to be taken into account. Consequently, the development of such an algorithm would first require extensive numerical studies which – although undeniably interesting – would clearly be beyond the scope of this article.
## IV Hypergeometric Series with a Negative Third Parameter
It is a direct consequence of its definition (22) that a hypergeometric series $`{}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ terminates after a finite number of terms if either $`a`$ or $`b`$ is a negative integer. Moreover, $`{}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ makes sense if both $`a`$ and $`c`$ are negative integers such that $`a=m`$ and $`c=mk`$ with $`k,m=1,2,\mathrm{}`$:
$${}_{2}{}^{}F_{1}^{}(m,b;mk;z)=\underset{\mu =0}{\overset{m}{}}\frac{(m)_\mu (b)_\mu }{(mk)_\mu \mu !}z^\mu .$$
(58)
If $`c`$ is a negative integer and if neither $`a`$ nor $`b`$ is a negative integer, then it follows from (22) that the hypergeometric series $`{}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ is in general undefined. However, the following limit exists for $`m=0,1,2,\mathrm{}`$ (p. 38 of ):
$`\underset{cm}{lim}{\displaystyle \frac{1}{\mathrm{\Gamma }(c)}}{}_{2}{}^{}F_{1}^{}(a,b;c;z)={\displaystyle \frac{(a)_{m+1}(b)_{m+1}z^{m+1}}{(m+1)!}}`$ (60)
$`\times {}_{2}{}^{}F_{1}^{}(a+m+1,b+m+1;m+2;z).`$
If $`c`$ is not a negative integer, there are no problems with Pochhammer symbols in the denominators of the terms of the series that could become zero. Nevertheless, unpleasant numerical problems occur even if $`c`$ is just a nonintegral negative real number. These problems can be demonstrated convincingly by trying to accelerate the convergence of two hypergeometric series $`{}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ which differ only by the sign of the third parameter $`c`$.
For that purpose, we consider in Table III and in Tables IV and V, respectively, the hypergeometric series with $`a=3/7`$, $`b=5/2`$, $`z=77/100`$, and $`c=\pm 7/2`$.
In Table III, Wynn’s epsilon algorithm, Eq. (II), Brezinski’s theta algorithm, Eq. (A), and the Levin-type transformations $`d_k^{(n)}(\zeta ,s_n)`$, Eq. (A36), and $`\delta _k^{(n)}(\zeta ,s_n)`$, Eq. (A40), are applied to the partial sums
$$s_n(z)=\underset{m=0}{\overset{n}{}}\frac{(3/7)_m(5/2)_m}{(7/2)_mm!}z^m$$
(61)
of the hypergeometric series $`{}_{2}{}^{}F_{1}^{}(3/7,5/2;7/2;z)`$ with $`z=77/100`$. The results in Table III show that the convergence of this series can indeed be accelerated quite effectively by sequence transformations. This is also true for the other sequence transformations discussed in Appendix A. The approximations with the highest transformation orders, that can be obtained from the partial sums $`s_0(z)`$, $`s_1(z)`$, …, $`s_{16}(z)`$ by Aitken’s iterated $`\mathrm{\Delta }^2`$ process (A) and by the iteration (A) of Brezinski’s theta algorithm, are
$`𝒜_7^{(1)}`$ $`=`$ $`1.463807099629,`$ (62)
$`𝒜_8^{(0)}`$ $`=`$ $`1.463807099563.`$ (63)
$`𝒥_5^{(0)}`$ $`=`$ $`1.463807143254,`$ (64)
$`𝒥_5^{(1)}`$ $`=`$ $`1.463807103421.`$ (65)
In Tables IV and V, we now consider a hypergeometric series $`{}_{2}{}^{}F_{1}^{}`$ which is identical with the one in Table III except that its third parameter is negative, i.e., we now have $`c=7/2`$. Thus, in Table IV we use Aitken’s iterated $`\mathrm{\Delta }^2`$ process, Eq. (A), Wynn’s epsilon algorithm, Eq. (II), and Brezinski’s theta algorithm, Eq. (A), for the acceleration of the convergence of the partial sums
$$s_n(z)=\underset{m=0}{\overset{n}{}}\frac{(3/7)_m(5/2)_m}{(7/2)_mm!}z^m.$$
(66)
of the hypergeometric series $`{}_{2}{}^{}F_{1}^{}(3/7,5/2;7/2;z)`$ with $`z=77/100`$, and in Table V we use the iteration of Brezinski’s theta algorithm, Eq. (A), and the Levin-type transformations $`d_k^{(n)}(\zeta ,s_n)`$, Eq. (A36), and $`\delta _k^{(n)}(\zeta ,s_n)`$, Eq. (A40).
So far, the transformation results had always been very good as well as very reliable. In contrast, the results in Tables IV and V are very bad and not reliable at all. For small transformation orders $`n`$, all transformations produce results which are by 5 orders of magnitude too small, and in the case of Brezinski’s theta algorithm, the wrong results even seem to have converged with an accuracy of 4 decimal digits. For increasing transformation orders $`n`$, there occur sudden and unmotivated sign changes, and only if $`n`$ approaches 30, at least the epsilon algorithm in Table IV and the Levin transformation $`d_n^{(0)}`$ in Table V converge to the correct result. For $`n=30`$, the other transformations show no indication of convergence and produce results which are still by some orders of magnitude too small.
How can the disturbingly bad performance of sequence transformations in Tables IV and V be explained. The terms of a hypergeometric series $`{}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ satisfy the 2-term recursion
$`{\displaystyle \frac{(a)_{n+1}(b)_{n+1}z^{n+1}}{(c)_{n+1}(n+1)!}}`$ (68)
$`={\displaystyle \frac{(a+n)(b+n)z}{(c+n)(n+1)}}{\displaystyle \frac{(a)_n(b)_nz^n}{(c)_nn!}}.`$
Obviously, the factor $`(a+n)(b+n)z/[(c+n)(n+1)]`$ on the right-hand side determines whether the terms increase or decrease in magnitude with increasing $`n`$. As long as this factor is greater than one in magnitude, the terms increase with increasing $`n`$, and as soon as this factor is smaller than one, the terms decrease.
Thus, we only have to determine those values of $`n`$ which satisfy
$$\left|\frac{(a+n)(b+n)z}{(c+n)(n+1)}\right|=\mathrm{\hspace{0.33em}1}$$
(69)
for given $`a`$, $`b`$, $`c`$, and $`z`$ in order to find out for which values of the index $`n`$ the terms of a hypergeometric series $`{}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ change their growth pattern.
In the case of the hypergeometric series in Table III with $`a=3/7`$, $`b=5/2`$, $`c=7/2`$, and $`z=77/100`$, there is no $`n>0`$ which satisfies condition (69). This implies that the terms of this series decrease monotonously in magnitude with increasing $`n0`$. Moreover, the results in Table III show that the convergence of this series can be accelerated quite effectively.
In the case of the hypergeometric series in Tables IV and V with $`c=7/2`$, the situation is more complicated since condition (69) is satisfied by $`n22`$. Thus, the terms of this series initially increase up to $`n=22`$, and only for $`n>22`$ they decrease and ultimately produce a convergent result. Accordingly, the partial sums (66) of this hypergeometric series initially look like the partial sums of a mildly divergent series, and only for $`n>22`$, they behave like the partial sums of a convergent series. Therefore, it should not be too surprising that sequence transformations perform poorly if they use as input data only the partial sums (66) with $`n22`$. However, even for $`22n30`$, only the epsilon algorithm and the Levin transformation $`d_n^{(0)}`$ ultimately converge to the correct result. This provides strong evidence that the irregular input data with small indices $`n`$ have a detrimental effect on the performance of sequence transformations with large transformation orders, which also use input data with a correct behavior.
We can test the hypothesis, that the irregular behavior of the initial partial sums (66) with $`n22`$ leads to the poor performance of sequence transformations in Tables IV and V, by skipping the terms up to $`n=22`$ in the transformation processes. Accordingly, Wynn’s epsilon algorithm, Eq. (II), and the Levin-type transformations $`d_k^{(n)}(\zeta ,s_n)`$, Eq. (A36), and $`\delta _k^{(n)}(\zeta ,s_n)`$, Eq. (A40), use in in Table VI the modified partial sums
$$s_n^{(22)}(z)=\underset{m=0}{\overset{n+22}{}}\frac{(3/7)_m(5/2)_m}{(7/2)_mm!}z^m$$
(70)
of the hypergeometric series $`{}_{2}{}^{}F_{1}^{}(3/7,5/2;7/2;z)`$ with $`z=77/100`$ as input data. This approach is possible since all sequence transformations considered in this article are quasi-linear, i.e., they satisfy (A45).
The results in Table VI indeed confirm our hypothesis since they are nearly as good as the results in Table III, at least with respect to the transformation orders that are needed to achieve a given relative accuracy.
The other sequence transformations discussed in Appendix A produce results that are less good than those shown in Table VI. The two approximations with the highest possible transformation orders, that can be produced from the partial sums $`s_0^{(22)}(z)`$, $`s_1^{(22)}(z)`$, …, $`s_{20}^{(22)}(z)`$ by Aitken’s iterated $`\mathrm{\Delta }^2`$ process (A), by Brezinski’s theta algorithm (A) and its iteration (A), are
$`𝒜_9^{(1)}`$ $`=`$ $`1.01014772243910^{+5},`$ (71)
$`𝒜_{10}^{(0)}`$ $`=`$ $`1.01014753770110^{+5}.`$ (72)
$`\theta _{12}^{(1)}`$ $`=`$ $`1.01146205162810^{+5},`$ (73)
$`\theta _{12}^{(2)}`$ $`=`$ $`1.01146201150110^{+5},`$ (74)
$`𝒥_7^{(0)}`$ $`=`$ $`1.01023390882510^{+5},`$ (75)
$`𝒥_7^{(1)}`$ $`=`$ $`1.01017605478610^{+5}.`$ (76)
The results in Table VI show that a hypergeometric series $`{}_{2}{}^{}F_{1}^{}`$ with a negative third parameter can be evaluated reliably with the help of sequence transformations if the nonregular leading terms are excluded from the transformation processes. Unfortunately, this may lead to new problems since the number of terms, that initially grow in magnitude, may become quite large, in particular if $`z`$ is close to one. For example, if we increase the argument of the hypergeometric series in Table VI from $`z=77/100`$ to $`z=87/100`$ or to $`z=97/100`$, then the number of terms, which initially grow in magnitude and have to be skipped, grow from $`n=22`$ to $`n=40`$ or even to $`n=179`$. Moreover, more negative values of $`c`$ also increase the number of terms that have to be skipped. For instance, if we consider the hypergeometric series $`{}_{2}{}^{}F_{1}^{}(3/7,5/2;13/2;z)`$ with $`z=77/100`$, $`z=87/100`$, or $`z=97/100`$, then we obtain $`n=35`$, $`n=63`$, or $`n=279`$, respectively. These examples show that it is in principle possible to construct hypergeometric series $`{}_{2}{}^{}F_{1}^{}`$ which can only be evaluated reliably with the help of sequence transformations if a very large number of terms is skipped in the transformation process.
In the case of a Gaussian hypergeometric series, this poses no unsurmountable problems. Firstly, it is a triviality to compute the terms, even for very large indices. Consequently, a brute force evaluation of a Gaussian hypergeometric series by adding up the terms is possible as long as the argument is not too close to the boundary of the circle of convergence. Secondly, the highly developed mathematical theory of these functions makes it possible to simplify the numerical task with the help of known transformation formulas. For example, if the number of terms, that have to be skipped, is large because the argument $`z`$ of the hypergeometric series with a negative third parameter is close to one, then it may be a good idea to use the analytic continuation formula (33). This would not necessarily solve the principal problems due to a negative third parameter, but the argument $`1z`$ of the two new hypergeometric series would then be small.
However, the probably simplest approach would be the use of recurrence formulas. In this approach, one would have to evaluate two hypergeometric series $`{}_{2}{}^{}F_{1}^{}`$ with suitable positive values of the third parameter, for example with the help of sequence transformations, and to compute recursively the numerical value of desired hypergeometric series with a negative third parameter. Of course, many recurrence formulas are known. Nevertheless, numerous new three-term recurrence formulas satisfied by the Gaussian hypergeometric function $`{}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ are derived in Appendix C.
Unfortunately, these alternative approaches are in general not available if we have to evaluate an infinite series whose terms are determined numerically and only behave like the terms of a hypergeometric series $`{}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ with a negative third parameter. In such a case, it cannot be excluded that the number of terms, that have to be skipped, would be so large that their computation would no longer be feasible. Then, neither the conventional process of successively adding up the terms of the series nor sequence transformations would be able to provide reliable approximations to the value of such an infinite series.
## V Summary and Conclusions
A sequence transformation is a rule $`𝒯`$ which transforms a sequence $`\{s_n\}_{n=0}^{\mathrm{}}`$ with the (generalized) limit $`s`$ to another sequence $`\{s_n^{}\}_{n=0}^{\mathrm{}}`$ having the same (generalized) limit $`s`$ but different remainders $`r_n^{}=s_n^{}s`$. The transformation process was successful if the new sequence has better numerical properties than the original sequence. For example, a sequence transformation $`𝒯`$ accelerates convergence, if the transformed remainders $`r_n^{}=s_n^{}s`$ converge more rapidly than the original remainders $`r_n=s_ns`$ according to (4), and $`𝒯`$ sums a divergent sequence, whose remainders $`r_n`$ do not vanish as $`n\mathrm{}`$, to its generalized limit $`s`$ if the transformed remainders $`r_n^{}`$ approach zero as $`n\mathrm{}`$.
As discussed in Section II, a sequence transformation $`𝒯`$ can be represented by an infinite set of doubly indexed quantities $`T_k^{(n)}`$ with $`k,n0`$ that can be displayed in a two-dimensional array called the table of $`𝒯`$. The superscript $`n`$ denotes the minimal index occurring in the finite string $`\{s_n,s_{n+1},\mathrm{}s_{n+l}\}`$ of sequence elements used for the computation of a given $`T_k^{(n)}`$. The subscript $`k`$ – usually called the order of the transformation – is a measure for the complexity of the transformation process which yields $`T_k^{(n)}`$.
A convergence acceleration or summation process tries to obtain a better approximation to the (generalized) limit of the input sequence by proceeding on a certain path in the table of $`𝒯`$. There is an in principle unlimited variety of different paths, but in practical applications either order-constant paths defined in (7) or index-constant paths defined in (8) are normally used.
On an order-constant path, a set $`\{s_n,s_{n+1},\mathrm{}s_{n+l}\}`$ of input data of fixed length is used and the starting index $`n`$ of this set is increased successively. In contrast, an index-constant path uses sets of input data of increasing length, and it is always tried to compute from a given set of input data that element $`T_k^{(n)}`$ which has the highest possible transformation order $`k`$.
Order-constant and index-constant paths differ substantially. For example, on an index-constant path the available information is exploited more efficiently than on an order-constant path. Accordingly, index-constant paths normally produce better transformation results. Moreover, order-constant paths cannot be used for the summation of a divergent sequence since increasing $`n`$ in the set $`\{s_n,s_{n+1},\mathrm{}s_{n+l}\}`$ of input data normally only increases divergence.
If, however, the leading elements of the input sequence $`\{s_n\}_{n=0}^{\mathrm{}}`$ behave irregularly, the principal advantages of index-constant paths can easily turn into disadvantages: If the sets $`\{s_n,s_{n+1},\mathrm{}s_{n+l}\}`$ of input data with increasing $`l`$ have a sufficiently small starting index $`n`$, then all transforms $`T_k^{(n)}`$ will be affected by irregular input data, albeit to a different degree. As shown by Tables IV and V, this can lead to unreliable or even completely nonsensical transformation results. In such a case, it is necessary to exclude the irregular input data from the transformation process. This can be accomplished by using either an order-constant path or an index-constant path with a sufficiently large starting index.
In Section III, the Gaussian hypergeometric function $`{}_{2}{}^{}F1(a,b;c;z)`$ is discussed, which is in general a multivalued function defined in the whole complex plane with branch points at $`z=1`$ and $`\mathrm{}`$. However, it is defined by the power series (22) which only converges for $`|z|<1`$. Accordingly, sequence transformations can either be used for speeding up convergence or for accomplishing an analytic continuation in the case of divergence. In Table I, an alternating hypergeometric series for $`\mathrm{ln}(1+z)`$ is summed effectively by sequence transformations for an argument $`z=7/2`$ which is is far away from the unit circle, and in Table II, another special hypergeometric series, which converges only in the interior of the unit circle, is evaluated for an argument $`z=(1+\mathrm{i}\sqrt{3})/2`$ that is located on the boundary of the unit chicle.
For most parameters $`a`$, $`b`$, and $`c`$, sequence transformations greatly facilitate the evaluation of a hypergeometric series $`{}_{2}{}^{}F_{1}^{}`$, and the good transformation results presented in Tables I and II are fairly typical. However, as discussed in Section IV, there is an important and instructive exception: If the third parameter $`c`$ of a hypergeometric series is a negative real number, then the terms of this series initially increase in magnitude and look even for $`|z|<1`$ like the terms of a mildly divergent series. Only for sufficiently large values of the index, the terms decrease and ultimately produce a convergent result.
The use of these irregular terms as input data seriously affects the performance of sequence transformations and leads to unreliable and sometimes even completely nonsensical results. This is demonstrated by applying in Table III and in Tables IV and V, respectively, sequence transformations to the partial sums of the hypergeometric series $`{}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ with $`a=3/7`$, $`b=5/2`$, $`z=77/100`$, and $`c=\pm 7/2`$.
In Table III, sequence transformations are applied to the hypergeometric series with the positive value of the third parameter. As expected, the transformation results are very good. However, in Tables IV and V, where the hypergeometric series with the negative value of the third parameter is considered, the transformation results are both unreliable and bad.
Nevertheless, it is possible to compute the hypergeometric series $`{}_{2}{}^{}F_{1}^{}(3/7,5/2;7/2;77/100)`$ efficiently and reliably with the help of sequence transformations. However, one cannot use an index-constant path with a small minimal index, as it was done in Tables IV and V. Instead, one should either use an order-constant path or an index-constant path with a sufficiently large minimal index, as it was done in Table VI, where all irregular terms were excluded from the transformation processes.
The sequence transformations, which are used in this article, are all described in Appendix A. The use of several transformations was quite intentional. The author wanted to make clear that problems due to irregular input data are not restricted to some special sequence transformations only. Of course, the results in Tables IV and V show that different sequence transformations respond differently to irregular input data. However, this is quite helpful and can protect us against misinterpretations. For example, in Table IV Brezinski’s theta algorithm produced transformation results which seemed to have converged with an accuracy of 4 decimal digits, but were actually by 5 orders of magnitude too small. Fortunately, the other transformations in Tables IV and V produced different results, which provided strong evidence that the transformation results were unreliable. Consequently, it is recommendable to use in convergence acceleration and summation processes more than a single transformation whenever possible. This is particularly important if numerically determined data are to be transformed, about which very little is known. If several different sequence transformations produce consistent results, then it is very likely that these results are indeed correct although it is of course clear that purely numerical results cannot be a substitute for a rigorous mathematical proof.
It looks like a contradiction that in Appendix B only the Levin-type transformation (A40) was used for the computation of the infinite coupling limit $`k_3`$ of the sextic anharmonic oscillator. However, the divergent perturbation series (B9), whose leading coefficients also show an irregular behavior, constitutes a very demanding summation problem, for which the other transformations discussed in Appendix A are not powerful enough.
Here, it should be emphasized once more that the two examples considered in this article – the Gaussian hypergeometric series $`{}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ with a negative third parameter and the divergent perturbation series (B9) for the infinite coupling limit $`k_3`$ – are comparatively simple model problems, and the leading irregular terms of their series expansions pose no unsurmountable computational problems. This is largely due to the fact that it is relatively easy to find out which terms behave irregularly. In the case of the perturbation series (B9), we are in the fortunate situation that the leading large-$`n`$ asymptotics (B11) of the coefficients $`c_n^{(3)}`$ is known, and in the case of the hypergeometric series, one only has to solve (69) in order to find out for which indices $`n`$ the terms change their growth pattern. Moreover, in the case of a $`{}_{2}{}^{}F_{1}^{}`$ there are numerous alternative computational approaches. For example, with the help of recurrence formulas the evaluation of Gaussian hypergeometric series with a negative third parameter can be avoided completely. Consequently, in Appendix C numerous new recurrence formulas satisfied by a $`{}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ are derived.
Finally, the author wishes to express his hope that this article will inspire additional research on the evaluation of special functions with the help of sequence transformations. There can be no doubt that sequence transformations are normally extremely useful tools for the evaluation of special functions, and since the terms of the series expansions for special functions are comparatively simple and explicitly known, we can even hope to gain additional insight from those cases in which sequence transformations fail to produce good transformation results.
###### Acknowledgements.
The author thanks the Fonds der Chemischen Industrie for financial support.
## A Sequence Transformations
This appendix gives a short description of all the sequence transformations that are used in this article. Further details plus additional references can be found in . Here, the same notation as in is used.
One of the oldest sequence transformations (see for instance pp. 90 - 91 of ) is Aitken’s $`\mathrm{\Delta }^2`$ formula :
$$𝒜_1^{(n)}=s_n\frac{[\mathrm{\Delta }s_n]^2}{\mathrm{\Delta }^2s_n}.$$
(A1)
The (forward) difference operator $`\mathrm{\Delta }`$ acts for all integers $`n0`$ on a function $`f(n)`$ according to
$$\mathrm{\Delta }f(n)=f(n+1)f(n).$$
(A2)
The Aitken formula (A1) is by construction exact for model sequences of the type $`s_n=s+c\lambda ^n`$ with $`c0`$ and $`\lambda 1`$. If the numerical values of three consecutive elements $`s_n`$, $`s_{n+1}`$, and $`s_{n+2}`$ of this model sequence are known, then the (generalized) limit $`s`$ of this sequence can be computed according to
$$𝒜_1^{(n)}=s,$$
(A3)
no matter whether the sequence converges ($`|\lambda |<1`$) or diverges ($`|\lambda |>1`$).
The power and practical usefulness of Aitken’s $`\mathrm{\Delta }^2`$ formula is of course limited since it is designed to eliminate only a single exponential term from the elements of the model sequence mentioned above. However, the quantities $`𝒜_1^{(n)}`$ can again be used as input data in (A1). Hence, the $`\mathrm{\Delta }^2`$ process can be iterated, yielding the following nonlinear recursive scheme \[13, Eq. (5.1-15)\]:
(A5)
$`𝒜_0^{(n)}`$ $`=`$ $`s_n,`$ (A6)
$`𝒜_{k+1}^{(n)}`$ $`=`$ $`𝒜_k^{(n)}{\displaystyle \frac{\left[\mathrm{\Delta }𝒜_k^{(n)}\right]^2}{\mathrm{\Delta }^2𝒜_k^{(n)}}}.`$ (A7)
In this article, the difference operator $`\mathrm{\Delta }`$ acts only on the superscript $`n`$ and not on the subscript $`k`$ of a doubly indexed quantity like $`𝒜_k^{(n)}`$. A more detailed discussion of Aitken’s iterated $`\mathrm{\Delta }^2`$ process as well as additional references can for instance be found in Section 5 of or in . The iteration of other sequence transformations is discussed in .
In the case of Aitken’s iterated $`\mathrm{\Delta }^2`$ process, the approximation to the limit of the input sequence with the highest possible transformation order depends upon the index $`m`$ of the last sequence element $`s_m`$ which was used in the recursion. If $`m`$ is either even or odd, $`m=2\mu `$ or $`m=2\mu +1`$, respectively, the approximations to the limit are chosen according to
$`\{s_0,s_1,\mathrm{},s_{2\mu }\}`$ $``$ $`𝒜_\mu ^{(0)},`$ (A8)
$`\{s_1,s_2,\mathrm{},s_{2\mu +1}\}`$ $``$ $`𝒜_\mu ^{(1)}.`$ (A9)
As in the case of Wynn’s $`ϵ`$ algorithm, these two relationships can with the help of the notation $`[[x]]`$ for the integral part of $`x`$ be expressed by a single equation (Eq. (5.2-6) of ):
$`\{s_{m2[[m/2]]},s_{m2[[m/2]]+1},\mathrm{},s_m\}`$ (A10)
$``$ $`𝒜_{[[m/2]]}^{(m2[[m/2]])}.`$ (A11)
The behavior of many practically relevant convergent sequences $`\{s_n\}_{n=0}^{\mathrm{}}`$ can be characterized by the asymptotic condition
$$\underset{n\mathrm{}}{lim}\frac{s_{n+1}s}{s_ns}=\rho ,$$
(A12)
where $`s=s_{\mathrm{}}`$ is the limit of the sequence $`\{s_n\}_{n=0}^{\mathrm{}}`$. This condition closely resembles the well known ratio test for infinite series. A convergent sequence satisfying (A12) with $`|\rho |<1`$ is called linearly convergent, and it is called logarithmically convergent if $`\rho =1`$.
It is one of the major weaknesses of the otherwise very powerful and very useful epsilon algorithm (II) that it does not work in the case of logarithmic convergence. Brezinski showed that this principal weakness can be overcome by a suitable modification of the recursive scheme (II), which leads to the so-called theta algorithm :
$`\theta _1^{(n)}`$ $`=`$ $`0,\theta _0^{(n)}=s_n,`$ (A14)
$`\theta _{2k+1}^{(n)}`$ $`=`$ $`\theta _{2k1}^{(n+1)}+\mathrm{\hspace{0.17em}1}/[\mathrm{\Delta }\theta _{2k}^{(n)}],`$ (A15)
$`\theta _{2k+2}^{(n)}`$ $`=`$ $`\theta _{2k}^{(n+1)}+{\displaystyle \frac{[\mathrm{\Delta }\theta _{2k}^{(n+1)}][\mathrm{\Delta }\theta _{2k+1}^{(n+1)}]}{\mathrm{\Delta }^2\theta _{2k+1}^{(n)}}}.`$ (A16)
As in the case of Aitken’s iterated $`\mathrm{\Delta }^2`$ process (A), it is assumed that the difference operator $`\mathrm{\Delta }`$ acts only upon the superscript $`n`$ and not on the subscript $`k`$.
Again, the approximation to the limit of the input sequence depends upon the index $`m`$ of the last sequence element $`s_m`$ which was used in the recursion. If we have $`m=3\mu `$, $`m=3\mu +1`$, or $`m=3\mu +2`$, respectively, the approximations to the limit with the highest transformation orders are chosen according to
$`\{s_0,s_1,\mathrm{},s_{3\mu }\}`$ $``$ $`\theta _{2\mu }^{(0)},`$ (A17)
$`\{s_1,s_2,\mathrm{},s_{3\mu +1}\}`$ $``$ $`\theta _{2\mu }^{(1)},`$ (A18)
$`\{s_2,s_3,\mathrm{},s_{3\mu +2}\}`$ $``$ $`\theta _{2\mu }^{(2)}.`$ (A19)
These three relationships can be expressed by a single equation (Eq. (10.2-8) of ):
$`\{s_{m3[[m/3]]},s_{m3[[m/3]]+1},\mathrm{},s_m\}`$ (A20)
$``$ $`\theta _{2[[m/3]]}^{(m3[[m/3]])}.`$ (A21)
Further details on the theta algorithm as well as additional references can be found in Section 2.9 of or in Sections 10 and 11 of .
As for example discussed in , new sequence transformations can be constructed by iterating explicit expressions for sequence transformations with low transformation orders. The best known example of such an iterated sequence transformation is probably Aitken’s iterated $`\mathrm{\Delta }^2`$ process (A) which is obtained by iterating (A1).
The same approach is also possible in the case of the theta algorithm (A). A suitable closed-form expression, which may be iterated, is (Eq. (10.3-1) of )
$`\vartheta _2^{(n)}=s_{n+1}`$ (A23)
$`{\displaystyle \frac{\left[\mathrm{\Delta }s_n\right]\left[\mathrm{\Delta }s_{n+1}\right]\left[\mathrm{\Delta }^2s_{n+1}\right]}{\left[\mathrm{\Delta }s_{n+2}\right]\left[\mathrm{\Delta }^2s_n\right]\left[\mathrm{\Delta }s_n\right]\left[\mathrm{\Delta }^2s_{n+1}\right]}}.`$
The iteration of this expression yields the following nonlinear recursive scheme (Eq. (10.3-6) of ):
$`𝒥_0^{(n)}=s_n,`$ (A25)
$`𝒥_{k+1}^{(n)}=𝒥_k^{(n+1)}`$ (A26)
$`{\displaystyle \frac{\left[\mathrm{\Delta }𝒥_k^{(n)}\right]\left[\mathrm{\Delta }𝒥_k^{(n+1)}\right]\left[\mathrm{\Delta }^2𝒥_k^{(n+1)}\right]}{\left[\mathrm{\Delta }𝒥_k^{(n+2)}\right]\left[\mathrm{\Delta }^2𝒥_k^{(n)}\right]\left[\mathrm{\Delta }𝒥_k^{(n)}\right]\left[\mathrm{\Delta }^2𝒥_k^{(n+1)}\right]}}.`$ (A27)
In convergence acceleration and summation processes, the iterated transformation $`𝒥_k^{(n)}`$ has similar properties as Brezinski’s theta algorithm from which it was derived: Both transformations are very powerful as well as very versatile. $`𝒥_k^{(n)}`$ is not only an effective accelerator for linear convergence as well as able to sum divergent alternating series, but it is also able to accelerate the convergence of many logarithmically convergent sequences and series .
Again, the approximation to the limit of the input sequence depends upon the index $`m`$ of the last sequence element $`s_m`$ which was used in the recursion. If we have $`m=3\mu `$, $`m=3\mu +1`$, or $`m=3\mu +2`$, respectively, the approximations to the limit with the highest transformation orders are chosen according to
$`\{s_0,s_1,\mathrm{},s_{3\mu }\}`$ $``$ $`𝒥_\mu ^{(0)},`$ (A28)
$`\{s_1,s_2,\mathrm{},s_{3\mu +1}\}`$ $``$ $`𝒥_\mu ^{(1)},`$ (A29)
$`\{s_2,s_3,\mathrm{},s_{3\mu +2}\}`$ $``$ $`𝒥_\mu ^{(2)}.`$ (A30)
These three relationships can be expressed by a single equation (Eq. (10.4-7) of ):
$`\{s_{m3[[m/3]]},s_{m3[[m/3]]+1},\mathrm{},s_m\}`$ (A31)
$``$ $`𝒥_{[[m/3]]}^{(m3[[m/3]])}.`$ (A32)
So far, only sequence transformations were considered which use as input data the elements of the sequence to be transformed. However, in some cases structural information on the dependence of the remainders $`r_n`$ on the index $`n`$ is available. For example, it is well known that the truncation error of a convergent series with strictly alternating and monotonously decreasing terms is bounded in magnitude by the first term not included in the partial sum and has the same sign as this term (see for instance p. 132 of ). The first term neglected is also the best simple estimate for the truncation error of a strictly alternating nonterminating hypergeometric series $`{}_{2}{}^{}F_{0}^{}(\alpha ,\beta ;x)`$ with $`\alpha ,\beta ,x>0`$ (Theorem 5.12-5 of ). Such an information on the behavior of the truncation errors can be extremely helpful in a convergence acceleration or summation process. Unfortunately, the sequence transformations considered so far are not able to benefit from it.
A convenient way of incorporating such an information into the transformation process consists in the use of remainder estimates $`\{\omega _n\}_{n=0}^{\mathrm{}}`$. Because of the explicit incorporation of the information contained in the remainder estimates, sequence transformations of that kind are potentially very powerful and as well as very versatile.
The best-known example of such a sequence transformation is Levin’s transformation which is both very versatile and very powerful :
$`_k^{(n)}(\zeta ,s_n,\omega _n)`$ (A34)
$`={\displaystyle \frac{{\displaystyle \underset{j=0}{\overset{k}{}}}(1)^j\left({\displaystyle \genfrac{}{}{0pt}{}{k}{j}}\right){\displaystyle \frac{(\zeta +n+j)^{k1}}{(\zeta +n+k)^{k1}}}{\displaystyle \frac{s_{n+j}}{\omega _{n+j}}}}{{\displaystyle \underset{j=0}{\overset{k}{}}}(1)^j\left({\displaystyle \genfrac{}{}{0pt}{}{k}{j}}\right){\displaystyle \frac{(\zeta +n+j)^{k1}}{(\zeta +n+k)^{k1}}}{\displaystyle \frac{1}{\omega _{n+j}}}}}.`$
The shift parameter $`\zeta `$ has to be positive in order to admit $`n=0`$ in (A34). The most obvious choice, which is always used in this article, is $`\zeta =1`$. Recurrence formulas for the numerator and denominator sums of $`_k^{(n)}(\zeta ,s_n,\omega _n)`$ can be found in Section 7.2 of .
Levin’s transformation is based on the assumption that the remainders $`r_n`$ of the input sequence can for all $`n0`$ be approximated by a remainder estimate $`\omega _n`$, which should be chosen such that $`s_ns=\omega _n\left[c+O(n^1)\right]`$ as $`n\mathrm{}`$, multiplied by a polynomial in $`1/(n+\zeta )`$ with $`\zeta >0`$. Levin introduced several simple remainder estimates for infinite series which give rise to several variants of Levin’s sequence transformation. Further details on Levin’s transformation can for instance be found in Section 7 of .
In this article, we only consider the remainder estimate
$$\omega _n=\mathrm{\Delta }s_n=a_{n+1},$$
(A35)
which was first proposed by Smith and Ford . It yields the following variant of Levin’s transformation (Eq. (7.3-9) of ):
$$d_k^{(n)}(\zeta ,s_n)=_k^{(n)}(\zeta ,s_n,\mathrm{\Delta }s_n).$$
(A36)
Levin’s transformation is based on the implicit assumption that the ratio $`[s_ns]/\omega _n`$ can be expressed as a power series in $`1/(n+\zeta )`$. A different class of sequence transformations can be derived by assuming that the ratio $`[s_ns]/\omega _n`$ can be expressed as a so-called factorial series according to (Section 8 of )
$$s_n=s+\omega _n\underset{j=0}{\overset{\mathrm{}}{}}c_j/(n+\zeta )_j,$$
(A37)
where $`(n+\zeta )_j=\mathrm{\Gamma }(n+j+\zeta )/\mathrm{\Gamma }(n+\zeta )`$ is a Pochhammer symbol (p. 3 of ). In this way, the sequence transformation (Eq. (8.2-7) of )
$`𝒮_k^{(n)}(\zeta ,s_n,\omega _n)`$ (A39)
$`={\displaystyle \frac{{\displaystyle \underset{j=0}{\overset{k}{}}}(1)^j\left({\displaystyle \genfrac{}{}{0pt}{}{k}{j}}\right){\displaystyle \frac{(\zeta +n+j)_{k1}}{(\zeta +n+k)_{k1}}}{\displaystyle \frac{s_{n+j}}{\omega _{n+j}}}}{{\displaystyle \underset{j=0}{\overset{k}{}}}(1)^j\left({\displaystyle \genfrac{}{}{0pt}{}{k}{j}}\right){\displaystyle \frac{(\zeta +n+j)_{k1}}{(\zeta +n+k)_{k1}}}{\displaystyle \frac{1}{\omega _{n+j}}}}}`$
can be derived which is formally very similar to Levin’s sequence transformation. The only difference between the transformations $`_k^{(n)}(\zeta ,s_n,\omega _n)`$ and $`𝒮_k^{(n)}(\zeta ,s_n,\omega _n)`$ is that the powers $`(\zeta +n+j)^{k1}`$ in (A34) are replaced by the Pochhammer symbols $`(\zeta +n+j)_{k1}`$ in (A39). Again, the shift parameter $`\zeta `$ has to be positive in order to admit $`n=0`$ in (A39), and the most obvious choice is also $`\zeta =1`$ which is exclusively used in this article. Recurrence formulas for the numerator and denominator sums of $`𝒮_k^{(n)}(\zeta ,s_n,\omega _n)`$ can be found in Section 8.3 of .
If we use the remainder estimate (A35) also in (A39), we obtain the following sequence transformation (Eq. (8.4-4) of ):
$$\delta _k^{(n)}(\zeta ,s_n)=𝒮_k^{(n)}(\zeta ,s_n,\mathrm{\Delta }s_n).$$
(A40)
It was shown in several articles that the transformation (A39) as well as its variant (A40) can be very effective , in particular if strongly divergent alternating series are to be summed.
In the case of the transformations (A36) and (A40), the approximation to the limit with the highest transformation order is given by
$$\{s_0,s_1,\mathrm{},s_{m+1}\}\mathrm{\Xi }_m^{(0)}(\zeta ,s_0),$$
(A41)
where $`\mathrm{\Xi }_k^{(n)}(\zeta ,s_n)`$ stands for either $`d_k^{(n)}(\zeta ,s_n)`$ or $`\delta _k^{(n)}(\zeta ,s_n)`$.
If the input data $`s_n`$ are the partial sums of a (formal) power series for some function $`f(z)`$ according to (12), $`s_n=f_n(z)`$, then the transformations (A36) and (A40) produce rational functions $`d_k^{(n)}(\zeta ,f_n(z))`$ and $`\delta _k^{(n)}(\zeta ,f_n(z))`$, whose numerator and denominator polynomials are of degrees $`k+n`$ and $`k`$ in $`z`$, respectively (Eqs. (4.25) and (4.26) of ). Moreover, the rational approximants $`d_k^{(n)}(\zeta ,f_n(z))`$ and $`\delta _k^{(n)}(\zeta ,f_n(z))`$ satisfy the following asymptotic error estimates as $`z0`$ (Eqs. (4.28) and (4.29) of ),
$`f(z)d_k^{(n)}(\zeta ,f_n(z))`$ $`=`$ $`O(z^{k+n+2}),`$ (A42)
$`f(z)\delta _k^{(n)}(\zeta ,f_n(z))`$ $`=`$ $`O(z^{k+n+2}),`$ (A43)
which are very similar to the well known accuracy-through-order relationships satisfied by Padé approximants .
It is a typical feature of all sequence transformations discussed in this Appendix that they are both homogeneous and translative: If the elements of two sequences $`\{s_n\}_{n=0}^{\mathrm{}}`$ and $`\{\sigma _n\}_{n=0}^{\mathrm{}}`$ satisfy
$$\sigma _n=as_n+b,$$
(A44)
where $`a`$ and $`b`$ are suitable constants, then
$$𝒯(\sigma _n,\sigma _{n+1},\mathrm{})=a𝒯(s_n,s_{n+1},\mathrm{})+b.$$
(A45)
Sequence transformations $`𝒯`$ satisfying this condition are called quasi-linear in the book by Brezinski and Redivo Zaglia . In Section 1.4 of this book, a detailed discussion of the properties of quasi-linear sequence transformations as well as further references can be found.
## B The Infinite Coupling Limit of the Sextic Anharmonic Oscillator
The detrimental effect of the irregular behavior of the leading elements of a sequence in convergence acceleration and summation processes is not restricted to mathematical model problems but occurs also in the mathematical treatment of scientific problems.
This will be shown by performing extensive summation calculations for the so-called strong coupling limit $`k_3`$ of the sextic anharmonic oscillator. The quartic ($`m=2`$), sextic ($`m=3`$), and octic ($`m=4`$) anharmonic oscillators are defined by the Hamiltonians
$$\widehat{H}(\beta )=\widehat{p}^2+\widehat{x}^2+\beta \widehat{x}^{2m},m=2,3,4,$$
(B1)
and the strong coupling limit $`k_m`$ of the ground state energy eigenvalue $`E^{(m}(\beta )`$ of this Hamiltonian is defined by
$$k_m=\underset{\beta \mathrm{}}{lim}E^{(m)}(\beta )/\beta ^{1/(m+1)}.$$
(B2)
Ever since the seminal work of Bender and Wu , the divergent weak coupling perturbation expansion
$$E^{(m)}(\beta )=\underset{n=0}{\overset{\mathrm{}}{}}b_n^{(m)}\beta ^n$$
(B3)
for the ground state energy of an anharmonic oscillator has been considered to be the model example of a strongly divergent quantum mechanical perturbation expansion which has to be summed in order to produce numerically useful results. Accordingly, there is an extensive literature on the summation of the divergent perturbation expansions of the anharmonic oscillators (see for example and references therein).
In addition to the divergent weak coupling expansion (B3), there is also a strong coupling expansion
$$E^{(m)}(\beta )=\beta ^{1/(m+1)}\underset{n=0}{\overset{\mathrm{}}{}}K_n^{(m)}\beta ^{2n/(m+1)}.$$
(B4)
It can be shown that this expansion converges for sufficiently large values of $`\beta `$ . However, the computation of the perturbative coefficients $`K_n^{(m)}`$ is very difficult (see for example and references therein).
It follows from (B2), (B3), and (B4) that the infinite coupling limit $`k_m`$ corresponds to the leading coefficient of the strong coupling expansion (B4) according to
$$k_m=K_0^{(m)}.$$
(B5)
The weak coupling perturbation expansion (B3) cannot be used in a straightforward way for a calculation of the strong coupling limit $`k_m`$. However, this can be accomplished comparatively easily with the help of the following renormalized weak coupling perturbation expansion (Eqs. (3.30) - (3.31) of ):
$$E^{(m)}(\beta )=(1\kappa )^{1/2}\underset{n=0}{\overset{\mathrm{}}{}}c_n^{(m)}\kappa ^n.$$
(B6)
This expansion is based on a renormalization scheme introduced by Vinette and Čížek . In this approach, the original coupling constant $`\beta [0,\mathrm{})`$ is transformed into a renormalized and explicitly $`m`$-dependent coupling constant $`\kappa [0,1)`$ according to (Eq. (3.19) of )
$$\beta =\frac{1}{B_m}\frac{\kappa }{(1\kappa )^{(m+1)/2}},$$
(B7)
where (Eq. (3.17) of )
$$B_m=\frac{m(2m1)!!}{2^{m1}}.$$
(B8)
For the sextic ($`m=3`$) case, these expressions correspond to $`B_3=45/4`$ and $`\beta =4\kappa /[45(1\kappa )^2]`$.
Thus, the infinite coupling limits $`k_3`$ of the sextic anharmonic oscillator can be expressed by the renormalized weak coupling expansion (B6) according to (Eqs. (3.43) and (3.44) of )
$$k_3=[45/4]^{1/4}\underset{n=0}{\overset{\mathrm{}}{}}c_n^{(3)}.$$
(B9)
The summation of either this or the perturbation series (B6) with $`m=3`$, from which (B9) was derived, is a formidable computational problem. This follows at once from the large-$`n`$ asymptotics of the renormalized perturbative coefficients for the sextic anharmonic oscillator (Eq. (3.34) of ):
$`c_n^{(3)}(1)^{n+1}{\displaystyle \frac{(128)^{1/2}}{\pi ^2}}`$ (B11)
$`\times \mathrm{\Gamma }(2n+1/2)(64/[45\pi ^2])^n,n\mathrm{}.`$
It should be noted that the summation of the perturbation series (B9) for $`k_3`$ is much more demanding than the summation of the divergent asymptotic expansions for special functions since their coefficients $`c_n`$ grow essentially like $`n!`$ . Although Padé approximants – or equivalently Wynn’s epsilon algorithm – are in principle capable of summing alternating divergent power series whose coefficients $`c_n`$ grow essentially like $`(2n)!`$ in magnitude, the convergence of Padé approximants is too slow to be practically useful. Moreover, it was shown in that the Levin transformation (A36) produces in the case of the perturbation expansions for the anharmonic oscillators sequences of approximants which initially seem to converge but which ultimately diverge. In contrast, the Levin-type transformation (A40) produces comparatively good results.
Thus, in analogy to we sum the perturbation series (B9) with the help of the Levin-type transformation $`\delta _k^{(n)}(\zeta ,s_n)`$ defined in (A40). It was shown in several articles that the sequence transformation (A40) as well as the transformation (A39), from which it was derived, is apparently very effective, in particular if strongly divergent alternating series are to be summed .
In our summation calculations for $`k_3`$ we use the renormalized coefficients $`c_\nu ^{(3)}`$ with $`0\nu 300`$. The coefficients were calculated using the exact rational arithmetics of Maple by solving a system of nonlinear difference equations as described in the Appendix of . Unfortunately, Eq. (A22) in the Appendix of , which specifies the system of nonlinear equations, contains a typographical error. Correct is
$`4jG_j^{(n)}`$ $`=`$ $`2(j+1)(2j+1)G_{j+1}^{(n)}+{\displaystyle \frac{1}{B_m}}G_{jm}^{(n1)}`$ (B13)
$`G_{j1}^{(n1)}\mathrm{\hspace{0.17em}2}{\displaystyle \underset{k=1}{\overset{n1}{}}}G_1^{(k)}G_j^{(nk)}.`$
The topic of this article is the study of the impact of irregular input data on the performance of sequence transformations. Accordingly, we have to investigate whether the renormalized coefficients $`c_n^{(3)}`$ behave irregularly for small indices $`n`$. For that purpose, we list in Table VII selected renormalized coefficients $`c_n^{(3)}`$ as well as the corresponding ratios
$$𝒞_n^{(3)}=\frac{(1)^{n+1}\pi ^2c_n^{(3)}}{\sqrt{128}\mathrm{\Gamma }(2n+1/2)}\left(\frac{45\pi ^2}{64}\right)^n,$$
(B14)
that are obtained by dividing the renormalized coefficients $`c_n^{(3)}`$ by the leading order of their large-$`n`$ asymptotics according to (B11).
The last column in Table VII shows quite clearly that the renormalized coefficients $`c_n^{(3)}`$ deviate for small values of $`n`$ considerably from their larger-order behavior. Firstly, the coefficients $`c_0^{(3)}`$ and $`c_1^{(3)}`$ apparently possess “wrong” sign. Secondly, the coefficients $`c_n^{(3)}`$ initially decrease in magnitude, and only for $`n4`$ they grow as they should according to (B11). Nevertheless, it is only a relatively mild irregularity, which affects only a few of the available coefficients $`c_\nu ^{(3)}`$ with $`0\nu 300`$.
The impact of the irregular coefficients $`c_n^{(3)}`$ with small indices $`n`$ can be checked by computing for $`l=0,1,2,\mathrm{}`$ and for $`n299l`$ the approximants
$$k_3^{(n,l)}=\delta _n^{(0)}(1,s_0^{(l)}),$$
(B15)
to the infinite coupling limit $`k_3`$ of the sextic anharmonic oscillator. Here, $`\delta _n^{(0)}`$ is the Levin-type transformation (A40), and
$$s_n^{(l)}=[45/4]^{1/4}\underset{\nu =0}{\overset{n+l}{}}c_\nu ^{(3)}$$
(B16)
is a partial sum of the perturbation series (B9) which skips the first $`l`$ terms in the transformation process.
In Table VIII, the approximants $`k_3^{(l,n)}`$ with the three highest possible values of $`n299l`$ are listed for $`l12`$.
If we compare the results with the extremely accurate result of Vinette and Čížek (Eq. (69) of )
$$k_3=\mathrm{\hspace{0.33em}1.144}802453797052763765457534149549,$$
(B17)
which was obtained nonperturbatively, we see that we gain 5 decimal digits by skipping the first 7 terms of the perturbation series (B9) for $`k_3`$. For $`l8`$, the accuracy of the summation results deteriorates again.
This is a very remarkable gain of accuracy, if we take into account that the summation of the perturbation series (B9) is a formidable task and that the leading terms $`c_n^{(3)}`$ display an only relatively mild irregularity, as shown in Table VII. Nevertheless, the results show quite clearly that the transformation order is not the only criterion which affects the performance of a sequence transformation. The results in Table VIII show that may be more effective to use smaller sets of more regular input data.
## C Recurrence Formulas for the Gaussian Hypergeometric Function
Many three-term recurrence formulas satisfied by the Gaussian hypergeometric function $`{}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ are known. For example, on pp. 557 - 558 of or on pp. 46 - 47 of , the following formulas can be found:
$`(ca){}_{2}{}^{}F_{1}^{}(a1,b;c;z)`$ (C1)
$`+`$ $`[2ac(ab)z]{}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ (C3)
$`+a(z1){}_{2}{}^{}F_{1}^{}(a+1,b;c;z)=\mathrm{\hspace{0.33em}0},`$
$`(cb){}_{2}{}^{}F_{1}^{}(a,b1;c;z)`$ (C4)
$`+`$ $`[2bc(ba)z]{}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ (C6)
$`+b(z1){}_{2}{}^{}F_{1}^{}(a,b+1;c;z)=\mathrm{\hspace{0.33em}0},`$
$`c(1c)(1z){}_{2}{}^{}F_{1}^{}(a,b;c1;z)`$ (C7)
$`+`$ $`c[c1(2cab1)z]{}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ (C9)
$`+(ca)(cb)z{}_{2}{}^{}F_{1}^{}(a,b;c+1;z)=\mathrm{\hspace{0.33em}0},`$
$`(ba){}_{2}{}^{}F_{1}^{}(a,b;c;z)+a{}_{2}{}^{}F_{1}^{}(a+1,b;c;z)`$ (C11)
$`b{}_{2}{}^{}F_{1}^{}(a,b+1;c;z)=\mathrm{\hspace{0.33em}0},`$
$`(bc){}_{2}{}^{}F_{1}^{}(a,b1;c;z)`$ (C12)
$`+`$ $`(cab){}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ (C14)
$`+a(1z){}_{2}{}^{}F_{1}^{}(a+1,b;c;z)=\mathrm{\hspace{0.33em}0},`$
$`c[a(cb)z]{}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ (C15)
$``$ $`ac(1z){}_{2}{}^{}F_{1}^{}(a+1,b;c;z)`$ (C17)
$`+(ca)(cb)z{}_{2}{}^{}F_{1}^{}(a,b;c+1;z)=\mathrm{\hspace{0.33em}0},`$
$`(1c){}_{2}{}^{}F_{1}^{}(a,b;c1;z)`$ (C18)
$`+`$ $`(ca1){}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ (C20)
$`+a{}_{2}{}^{}F_{1}^{}(a+1,b;c;z)=\mathrm{\hspace{0.33em}0},`$
$`(ac){}_{2}{}^{}F_{1}^{}(a1,b;c;z)`$ (C21)
$`+`$ $`(cab){}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ (C23)
$`+b(1z){}_{2}{}^{}F_{1}^{}(a,b+1;c;z)=\mathrm{\hspace{0.33em}0},`$
$`(ac){}_{2}{}^{}F_{1}^{}(a1,b;c;z)`$ (C24)
$`+`$ $`(cb){}_{2}{}^{}F_{1}^{}(a,b1;c;z)`$ (C26)
$`+(ba)(1z){}_{2}{}^{}F_{1}^{}(a,b;c;z)=\mathrm{\hspace{0.33em}0},`$
$`(c){}_{2}{}^{}F_{1}^{}(a1,b;c;z)`$ (C27)
$`+`$ $`c(1z){}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ (C29)
$`+(cb)z{}_{2}{}^{}F_{1}^{}(a,b;c+1;z)=\mathrm{\hspace{0.33em}0},`$
$`(ca){}_{2}{}^{}F_{1}^{}(a1,b;c;z)`$ (C30)
$``$ $`(c1)(1z){}_{2}{}^{}F_{1}^{}(a,b;c1;z)`$ (C32)
$`+[a1(cb1)z]{}_{2}{}^{}F_{1}^{}(a,b;c;z)=\mathrm{\hspace{0.33em}0},`$
$`c[b(ca)z]{}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ (C33)
$``$ $`bc(1z){}_{2}{}^{}F_{1}^{}(a,b+1;c;z)`$ (C35)
$`+(ca)(cb)z{}_{2}{}^{}F_{1}^{}(a,b;c+1;z)=\mathrm{\hspace{0.33em}0},`$
$`(1c){}_{2}{}^{}F_{1}^{}(a,b;c1;z)`$ (C36)
$`+`$ $`(cb1){}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ (C38)
$`+b{}_{2}{}^{}F_{1}^{}(a,b+1;c;z)=\mathrm{\hspace{0.33em}0},`$
$`(c){}_{2}{}^{}F_{1}^{}(a,b1;c;z)`$ (C39)
$`+`$ $`c(1z){}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ (C41)
$`+(ca)z{}_{2}{}^{}F_{1}^{}(a,b;c+1;z)=\mathrm{\hspace{0.33em}0},`$
$`(cb){}_{2}{}^{}F_{1}^{}(a,b1;c;z)`$ (C42)
$``$ $`(c1)(1z){}_{2}{}^{}F_{1}^{}(a,b;c1;z)`$ (C44)
$`+[b1(ca1)z]{}_{2}{}^{}F_{1}^{}(a,b;c;z)=\mathrm{\hspace{0.33em}0}.`$
It is a typical feature of these recurrence formulas that there is a hypergeometric function $`{}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ plus two other $`{}_{2}{}^{}F_{1}^{}`$’s which differ with respect to only one of the three parameters by $`\pm 1`$. However, recurrence formulas, which contain a hypergeometric function $`{}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ plus two other $`{}_{2}{}^{}F_{1}^{}`$’s which differ with respect to two or even three parameters by $`\pm 1`$, can be constructed comparatively easily. For that purpose, we combine those of the recurrence formulas given above, whose the third parameter $`c`$ assumes at least two different values, with the linear transformation formulas (see for example p. 559 of or p. 47 of )
$`{}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ $`=`$ $`(1z)^{cab}{}_{2}{}^{}F_{1}^{}(ca,cb;c;z)`$ (C45)
$`=`$ $`(1z)^a{}_{2}{}^{}F_{1}^{}(a,cb;c;z/(z1))`$ (C46)
$`=`$ $`(1z)^b{}_{2}{}^{}F_{1}^{}(ca,b;c;z/(z1)).`$ (C47)
For the derivation of new recurrence formulas, we replace in (C9) $`a`$ by $`ca`$ and $`b`$ by $`cb`$. This yields:
$`c(1c)(1z){}_{2}{}^{}F_{1}^{}(ca,cb;c1;z)`$ (C48)
$`+`$ $`c[c1(a+b1)z]{}_{2}{}^{}F_{1}^{}(ca,cb;c;z)`$ (C50)
$`+abz{}_{2}{}^{}F_{1}^{}(ca,cb;c+1;z)=\mathrm{\hspace{0.33em}0}.`$
If we now combine this relationship with the linear transformation (C45), we obtain the following recurrence formula, where all three parameters of the hypergeometric functions change simultaneously:
$`c(1c){}_{2}{}^{}F_{1}^{}(a1,b1;c1;z)`$ (C51)
$`+`$ $`c[c1(a+b1)z]{}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ (C53)
$`+abz(1z){}_{2}{}^{}F_{1}^{}(a+1,b+1;c+1;z)=\mathrm{\hspace{0.33em}0}.`$
If we now proceed in (C17), (C20), and in (C29) - (C44) in exactly the same way, we obtain the following recurrence formulas:
$`c(cabz){}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ (C54)
$``$ $`c(ca){}_{2}{}^{}F_{1}^{}(a1,b;c;z)`$ (C56)
$`+abz(1z){}_{2}{}^{}F_{1}^{}(a+1,b+1;c+1;z)=\mathrm{\hspace{0.33em}0},`$
$`(1c){}_{2}{}^{}F_{1}^{}(a1,b1;c1;z)`$ (C57)
$`+`$ $`(a1)(1z){}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ (C59)
$`+(ca){}_{2}{}^{}F_{1}^{}(a1,b;c;z)=\mathrm{\hspace{0.33em}0},`$
$`c{}_{2}{}^{}F_{1}^{}(a+1,b;c;z)`$ (C60)
$``$ $`c{}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ (C62)
$`bz{}_{2}{}^{}F_{1}^{}(a+1,b+1;c+1;z)=\mathrm{\hspace{0.33em}0},`$
$`a(1z){}_{2}{}^{}F_{1}^{}(a+1,b;c;z)`$ (C63)
$`+`$ $`(1c){}_{2}{}^{}F_{1}^{}(a1,b1;c1;z)`$ (C65)
$`+[ca1(b1)z]{}_{2}{}^{}F_{1}^{}(a,b;c;z)=\mathrm{\hspace{0.33em}0},`$
$`c(cbaz){}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ (C66)
$``$ $`c(cb){}_{2}{}^{}F_{1}^{}(a,b1;c;z)`$ (C68)
$`+abz(1z){}_{2}{}^{}F_{1}^{}(a+1,b+1;c+1;z)=\mathrm{\hspace{0.33em}0},`$
$`(1c){}_{2}{}^{}F_{1}^{}(a1,b1;c1;z)`$ (C69)
$`+`$ $`(b1)(1z){}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ (C71)
$`+(cb){}_{2}{}^{}F_{1}^{}(a,b1;c;z)=\mathrm{\hspace{0.33em}0},`$
$`c{}_{2}{}^{}F_{1}^{}(a,b+1;c;z)`$ (C72)
$``$ $`c{}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ (C74)
$`az{}_{2}{}^{}F_{1}^{}(a+1,b+1;c+1;z)=\mathrm{\hspace{0.33em}0},`$
$`b(1z){}_{2}{}^{}F_{1}^{}(a,b+1;c;z)`$ (C75)
$`+`$ $`(1c){}_{2}{}^{}F_{1}^{}(a1,b1;c1;z)`$ (C77)
$`+[cb1(a1)z]{}_{2}{}^{}F_{1}^{}(a,b;c;z)=\mathrm{\hspace{0.33em}0}.`$
Not all of these recurrence formulas are independent. For example, (C56) and (C68) can be transformed into each other by interchanging $`a`$ and $`b`$. This is also true for the pairs (C59) and (C71), (C62) and (C74), and (C65) and (C77), which can be transformed into each other by interchanging $`a`$ and $`b`$.
For the derivation of recurrence formulas which differ with respect to two parameters by $`\pm 1`$, we replace in (C9) $`a`$ by $`ca`$ and $`z`$ by $`z/(z1)`$. This yields:
$`{\displaystyle \frac{c(1c)}{1z}}{}_{2}{}^{}F_{1}^{}(ca,b;c1;{\displaystyle \frac{z}{z1}})`$ (C78)
$`+`$ $`{\displaystyle \frac{c[c1+(ab)z]}{z1}}{}_{2}{}^{}F_{1}^{}(ca,b;c;{\displaystyle \frac{z}{z1}})`$ (C80)
$`+{\displaystyle \frac{a(cb)z}{z1}}{}_{2}{}^{}F_{1}^{}(ca,b;c+1;{\displaystyle \frac{z}{z1}})=\mathrm{\hspace{0.33em}0}.`$
If we now combine this relationship with the linear transformation (C47), we obtain the following recurrence formula, where the first and the third parameter of the hypergeometric series change simultaneously:
$`c(1c){}_{2}{}^{}F_{1}^{}(a1,b;c1;z)`$ (C81)
$``$ $`c[c1+(ab)z]{}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ (C83)
$`+a(cb)z{}_{2}{}^{}F_{1}^{}(a+1,b;c+1;z)=\mathrm{\hspace{0.33em}0}.`$
If we now proceed in (C17), (C20), and in (C29) - (C44) in exactly the same way, we obtain the following recurrence formulas:
$`c[ca+(ab)z]{}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ (C84)
$``$ $`c(ca){}_{2}{}^{}F_{1}^{}(a1,b;c;z)`$ (C86)
$`a(cb)z{}_{2}{}^{}F_{1}^{}(a+1,b;c+1;z)=\mathrm{\hspace{0.33em}0},`$
$`(1c){}_{2}{}^{}F_{1}^{}(a1,b;c1;z)`$ (C87)
$`+`$ $`(a1){}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ (C89)
$`+(ca){}_{2}{}^{}F_{1}^{}(a1,b;c;z)=\mathrm{\hspace{0.33em}0},`$
$`c(1z){}_{2}{}^{}F_{1}^{}(a+1,b;c;z)`$ (C90)
$``$ $`c{}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ (C92)
$`+(cb)z{}_{2}{}^{}F_{1}^{}(a+1,b;c+1;z)=\mathrm{\hspace{0.33em}0},`$
$`a(1z){}_{2}{}^{}F_{1}^{}(a+1,b;c;z)`$ (C93)
$`+`$ $`(1c){}_{2}{}^{}F_{1}^{}(a1,b;c1;z)`$ (C95)
$`+[ca1+(ab)z]{}_{2}{}^{}F_{1}^{}(a,b;c;z)=\mathrm{\hspace{0.33em}0},`$
$`c[b+(ab)z]{}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ (C96)
$``$ $`bc(1z){}_{2}{}^{}F_{1}^{}(a,b+1;c;z)`$ (C98)
$`a(cb)z{}_{2}{}^{}F_{1}^{}(a+1,b;c+1;z)=\mathrm{\hspace{0.33em}0},`$
$`(1c){}_{2}{}^{}F_{1}^{}(a1,b;c1;z)`$ (C99)
$`+`$ $`(cb1){}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ (C101)
$`+b(1z){}_{2}{}^{}F_{1}^{}(a,b+1;c;z)=\mathrm{\hspace{0.33em}0},`$
$`c{}_{2}{}^{}F_{1}^{}(a,b1;c;z)`$ (C102)
$``$ $`c{}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ (C104)
$`+az{}_{2}{}^{}F_{1}^{}(a+1,b;c+1;z)=\mathrm{\hspace{0.33em}0},`$
$`(cb){}_{2}{}^{}F_{1}^{}(a,b1;c;z)`$ (C105)
$`+`$ $`(1c){}_{2}{}^{}F_{1}^{}(a1,b;c1;z)`$ (C107)
$`+[b1+(ab)z]{}_{2}{}^{}F_{1}^{}(a,b;c;z)=\mathrm{\hspace{0.33em}0}.`$
Next, we replace in (C9) $`b`$ by $`cb`$ and $`z`$ by $`z/(z1)`$. This yields:
$`{\displaystyle \frac{c(1c)}{1z}}{}_{2}{}^{}F_{1}^{}(a,cb;c1;{\displaystyle \frac{z}{z1}})`$ (C108)
$`+`$ $`{\displaystyle \frac{c[c1(ab)z]}{z1}}{}_{2}{}^{}F_{1}^{}(a,cb;c;{\displaystyle \frac{z}{z1}})`$ (C110)
$`+{\displaystyle \frac{(ca)bz}{z1}}{}_{2}{}^{}F_{1}^{}(a,cb;c+1;{\displaystyle \frac{z}{z1}})=\mathrm{\hspace{0.33em}0}.`$
If we now combine this relationship with the linear transformation (C46), we obtain the following recurrence formula, where the second and the third parameter of the hypergeometric series change simultaneously:
$`c(1c){}_{2}{}^{}F_{1}^{}(a,b1;c1;z)`$ (C111)
$`+`$ $`c[c1(ab)z]{}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ (C113)
$`(ca)bz{}_{2}{}^{}F_{1}^{}(a,b+1;c+1;z)=\mathrm{\hspace{0.33em}0}.`$
If we now proceed in (C17), (C20), and in (C29) - (C44) in exactly the same way, we obtain the following recurrence formulas:
$`c[(ab)za]{}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ (C114)
$`+`$ $`ac(1z){}_{2}{}^{}F_{1}^{}(a+1,b;c;z)`$ (C116)
$`+(ca)bz{}_{2}{}^{}F_{1}^{}(a,b+1;c+1;z)=\mathrm{\hspace{0.33em}0},`$
$`(1c){}_{2}{}^{}F_{1}^{}(a,b1;c1;z)`$ (C117)
$`+`$ $`(ca1){}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ (C119)
$`+a(1z){}_{2}{}^{}F_{1}^{}(a+1,b;c;z)=\mathrm{\hspace{0.33em}0},`$
$`c{}_{2}{}^{}F_{1}^{}(a1,b;c;z)`$ (C120)
$``$ $`c{}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ (C122)
$`+bz{}_{2}{}^{}F_{1}^{}(a,b+1;c+1;z)=\mathrm{\hspace{0.33em}0},`$
$`(ca){}_{2}{}^{}F_{1}^{}(a1,b;c;z)`$ (C123)
$`+`$ $`(1c){}_{2}{}^{}F_{1}^{}(a,b1;c1;z)`$ (C125)
$`+[a1(ab)z]{}_{2}{}^{}F_{1}^{}(a,b;c;z)=\mathrm{\hspace{0.33em}0},`$
$`c[cb(ab)z]{}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ (C126)
$``$ $`c(cb){}_{2}{}^{}F_{1}^{}(a,b1;c;z)`$ (C128)
$`(ca)bz{}_{2}{}^{}F_{1}^{}(a,b+1;c+1;z)=\mathrm{\hspace{0.33em}0},`$
$`(1c){}_{2}{}^{}F_{1}^{}(a,b1;c1;z)`$ (C129)
$`+`$ $`(b1){}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ (C131)
$`+(cb){}_{2}{}^{}F_{1}^{}(a,b1;c;z)=\mathrm{\hspace{0.33em}0},`$
$`c(1z){}_{2}{}^{}F_{1}^{}(a,b+1;c;z)`$ (C132)
$``$ $`c{}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ (C134)
$`+(ca)z{}_{2}{}^{}F_{1}^{}(a,b+1;c+1;z)=\mathrm{\hspace{0.33em}0},`$
$`b(1z){}_{2}{}^{}F_{1}^{}(a,b+1;c;z)`$ (C135)
$`+`$ $`(1c){}_{2}{}^{}F_{1}^{}(a,b1;c1;z)`$ (C137)
$`+[cb1(ab)z]{}_{2}{}^{}F_{1}^{}(a,b;c;z)=\mathrm{\hspace{0.33em}0}.`$
The two groups of recursions (C83) - (C107) and (C113) - (C137), respectively, are not independent. They can be transformed into each other by interchanging $`a`$ and $`b`$.
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# Communication Capacity of Quantum Computation
## Abstract
By considering quantum computation as a communication process, we relate its efficiency to a communication capacity. This formalism allows us to rederive lower bounds on the complexity of search algorithms. It also enables us to link the mixedness of a quantum computer to its efficiency. We discuss the implications of our results for quantum measurement.
Any computation, both classical and quantum, is formally identical to a communication in time. At time $`t=0`$, the programmer sets the computer to accomplish any one of several possible tasks. Each of these tasks can be regarded as embodying a different message. Another programmer can obtain this message by looking at the output of the computer when the computation is finished at time $`t=t_1`$. Recent years have witnessed a surge of interest in both quantum computation and quantum communication . Computation based on quantum principles allows for more efficient algorithms for solving certain problems than algorithms based on purely classical principles. Quantum communication, on the other hand, can be used for unconditionally secure secret key distribution . However, till date, these two areas (i.e quantum computation and quantum communication) have developed independently. In this letter we connect the classical capacity of a quantum communication channel with the efficiency of quantum computation. This offers an unifying framework for quantum information processing.
Let us first introduce a few definitions and a communication model of quantum computation. We have two programmers, the sender and the reciever and two registers, the memory ($`M`$) register and the computational ($`C`$) register. The sender prepares the memory register in a certain quantum state $`|i_M`$ which encodes the problem to be solved. For example, in the case of factorization , this register will store the number to be factored. In case of a database search , this register will store the state of the database to be searched. The number $`N`$ of possible states $`|i_M`$ will, of course, be limited by the greatest number that the given computer could factor or the largest database it could search. The reciever prepares the computational register in some initial state $`\rho _C^0`$. Both the sender and the reciever feed the registers (prepared by them) to the quantum computer. The quantum computer implements the following general transformation on the registers
$$(|ii|)_M\rho _C^0(|ii|)_MU_i\rho _CU_i^{}.$$
(1)
The resulting state $`\rho _C(i)=U_i\rho _C^0U_i^{}`$ of the computational register contains the answer to the computation and is measured by the reciever. As the quantum computation should work for any $`|i_M`$, it should also work for any mixture $`_i^Np_i(|ii|)_M`$, where $`p_i`$ are probabilities. For the sender to use the above computation as a communication protocol, he has to prepare any one of the states $`|i_M`$ with an apriori probability $`p_i`$. The entire input ensemble is thus $`_i^Np_i(|ii|)_M\rho _C^0`$. Due to the quantum computation, this becomes
$$\underset{i}{\overset{N}{}}p_i(|ii|)_M\rho _C^0\underset{i}{\overset{N}{}}p_i(|ii|)_M\rho _C(i).$$
(2)
Whereas before the quantum computation, the two registers where completely uncorrelated (mutual information is zero), at the end, the mutual information becomes
$`I_{MC}:`$ $`=`$ $`S(\rho _M)+S(\rho _C)S(\rho _{MC})`$ (3)
$`=`$ $`S(\rho _C){\displaystyle \underset{i}{\overset{N}{}}}p_iS(\rho _C(i)),`$ (4)
where $`\rho _M`$ and $`\rho _C`$ are the reduced density operators for the two registers, $`\rho _{MC}`$ is the density operator of entire $`M+C`$ system and $`S(\rho )=\text{Tr}\rho \mathrm{log}\rho `$ is the von Neumann entropy (for conventional reasons we will use $`\mathrm{log}_2`$ in all calculations). Notice that the value of the mutual information (i.e correlations) is equal to the Holevo bound $`H=S(\rho _C)_i^Np_iS(\rho _C(i))`$ for the classical capacity of a quantum communication channel (Note that $`\rho _C=_i^Np_i\rho _C(i)`$). This tells us how much information the reciever can obtain about the choice $`|i_M`$ made by the sender by measuring the computational register. The maximum value of $`H`$ is obtained when the states $`\rho _C(i)`$ are pure and orthogonal. Moreover, the sender conveys the maximum information when all the message states have equal apriori probability (which also maximizes the channel capacity). In that case the mutual information (channel capacity) at the end of the computation is $`\mathrm{log}N`$. Thus the communication capacity $`I_{MC}`$ (given by Eq.(3)) gives an index of the efficiency of a quantum computation. The target of a quantum computation is to achieve the maximum possible communication capacity consistent with given initial states of the quantum computer. If one breaks down the general unitary transformation $`U_i`$ of a quantum algorithm into several succesive unitary transformations, then the maximum capacity may be achieved only after several steps. In each of the smaller unitary transformations, the mutual information between the $`M`$ and the $`C`$ registers (i.e the communication capacity) increases by a certain amount. When its total value reaches the maximum possible value consistent with a given initial state of the quantum computer, the computation is regarded as being complete.
We now proceed to illustrate one immediate application of the above formalism. Any general quantum algorithm has to have a certain number of queries into the memory register (this is neccessiated by the fact that the transformation on the computational register has to depend on the problem at hand, encoded in $`|i_M`$). These queries can be considered to be implemented by a black box into which the states of both the memory and the computational registers are fed. The number of such queries needed in a certain quantum algorithm gives the black box complexity of that algorithm and is a lower bound on the complexity of the whole algorithm. Recently, Ambainis showed in a very elegant paper that if the memory register was prepared initially in the superposition $`_i^N|i_M`$, then, in a search algorithm, $`O(\sqrt{N})`$ queries would be needed to completely entangle it with the computational register. This gives a lower bound on the number of queries in a search algorithm. In a manner analogous to his, we will calculate the change in mutual information between the memory and the computational registers (from Eq.(3)) in one query step. The number of queries needed to increase the mutual information to $`\mathrm{log}N`$ (for perfect communication between the sender and the reciever), is then a lower bound on the complexity of the algorithm.
Any search algorithm (whether quantum or classical, irrespective of its explicit form), will have to find a match for the state $`|i_M`$ of the $`M`$ register among the states $`|j_C`$ of the $`C`$ register and associate a marker to the state that matches (Here, $`|j_C`$ is a complete orthonormal basis for the $`C`$ register). The most general way of doing such a query in the quantum case is the black box unitary transformation
$$U_B|i_M|j_C=(1)^{\delta _{ij}}|i_M|j_C.$$
(5)
Any other unitary transformation performing a query matching the states of the $`M`$ and the $`C`$ registers, could be constructed from the above type of query. We would like to put a bound on the change of the mutual information in one such black box step. Let the memory states $`|i_M`$ be available to the sender with equal apriori probability so that the communication capacity is a maximum. His initial ensemble is then $`\frac{1}{N}_i^N(|ii|)_M`$. Let the reciever prepare the $`C`$ register in an initial pure state $`\psi ^0`$ (in fact, the power of quantum computation stems from the ability of the reciever to prepare pure state superpositions of form $`\frac{1}{N}_j^N|j_C`$). In general, there will be many black box steps on the initial ensemble before perfect correlations between the $`M`$ and the $`C`$ registers is set up. Let, after the $`k`$th black box step, the state of the system be
$$\rho ^k=\frac{1}{N}\underset{i}{\overset{N}{}}(|ii|)_M(|\psi ^k(i)\psi ^k(i)|)_C$$
(6)
where
$$|\psi ^k(i)_C=\underset{j}{}\alpha _{ij}^k|j_C.$$
(7)
The $`(k+1)`$th black box step changes this state to $`\rho ^{k+1}=\frac{1}{N}_i^N(|ii|)_M(|\psi ^{k+1}(i)\psi ^{k+1}(i)|)_C`$ with
$$|\psi ^{(k+1)}(i)=\underset{i,j}{\overset{N}{}}\alpha _{ij}^k(1)^{\delta _{ij}}|i_M|j_C.$$
(8)
Thus we only have to evaluate the difference of mutual information between the $`M`$ and the $`C`$ register for the states. This difference of mutual information (when computed from Eq.(3)) can be shown to be the difference $`|S(\rho _C^{k+1})S(\rho _C^k)|`$ . This quantity is bounded from the above by
$`|S(\rho _C^{k+1})`$ $``$ $`S(\rho _C^k)|d_B(\rho _C^k,\rho _C^{k+1})\mathrm{log}N`$ (9)
$``$ $`d_B(\rho _C^k,\rho _C^{k+1})\mathrm{log}d_B(\rho _C^k,\rho _C^{k+1})`$ (10)
where, $`d_B(\sigma ,\rho )=\sqrt{1F^2(\sigma ,\rho )}`$ is the Bures metric and $`F(\sigma ,\rho )=\text{Tr}\sqrt{\sqrt{\rho }\sigma \sqrt{\rho }}`$ is the fidelity. Using methods similar to Ambainis , it can be shown that $`F(\rho _C^k,\rho _C^{k+1})\frac{N2}{N}`$ from which it follows that
$$|S(\rho _C^{k+1})S(\rho _C^k)|\frac{3}{\sqrt{N}}\mathrm{log}N.$$
(11)
This means that at least $`O(\sqrt{N})`$ steps are needed to produce full correlations (maximum mutual information of value $`\mathrm{log}N`$) between the two registers. This gives the black box lower bound on the complexity of any quantum search algorithm. Of course, we know that there also exists an algorithm achieving this bound due to Grover and this has been proven to be optimal .
We now use Grover’s algorithm to show how the mutual information varies with time in a quantum search. The general sequence described by Cleve et. al for Grover’s algorithm will be used in this letter. The algorithm consists of repeated blocks, each consisting of a Hadamard transform on each qubit of the $`C`$ register, followed by a $`U_B`$ (our black box transformation), followed by another Hadamard transform on each qubit of the $`C`$ register and finally a phase flip $`f_0`$ of the the $`|00\mathrm{}0_C`$ state of the $`C`$ register (See fig.1). This block can then be repeated as many times as is necessary to bring the mutual information to its maximum value of $`\mathrm{log}N`$, which, as we have shown in Eq.(11) to be $`O(\sqrt{N})`$. Note that the only transformation correlating the $`M`$ and $`C`$ registers is the black box transformation $`U_B`$ and all the other transformations are done only on the $`C`$ register and therefore do not change the mutual information between the two registers. In fig.2 we have plotted the variation of mutual information between the $`M`$ and the $`C`$ registers (i.e the communication capacity of the quantum computation) with the number of iterations of the block in Grover’s algorithm. It is seen that the mutual information oscillates with the number of iterations. Fig.2 is plotted for a four qubit computational register which can search a database of $`16`$ entries. It is seen that the period is roughly $`6`$, which means that the number of steps needed to achieve maximum mutual information is roughly $`3`$. This is well above our bound for the minimum number of steps, which is $`4/3`$ in this case.
The three graphs (a), (b) and (c) in Fig.2 are for different values of initial mixedness of the $`C`$ register. We find that the mutual information fails to rise to the maximum value of $`\mathrm{log}N`$ when the state of the computational register is mixed. Our formalism thus allows us to calculate the performance of a quantum computation as a function of the mixedness (quantified by the von Neumann entropy) of the computational register. We can put a bound on the entropy of the second register after which the quantum search becomes as inefficient as the classical search. If the initial entropy $`S(\rho _C^0)`$ of the $`C`$ register exceeds $`\frac{1}{2}\mathrm{log}N`$, then the change in mutual information between the $`M`$ and the $`C`$ registers in the course of the entire quantum computation would be at most $`\mathrm{log}\sqrt{N}`$. This can be achieved by a classical database search in $`\sqrt{N}`$ steps. So there is no advantage in using quantum evolution when the initial state is too mixed. Note that our condition
$$S(\rho _C^0)\frac{1}{2}\mathrm{log}N$$
(12)
for no quantum speedup in the search algorithm is only a sufficient condition and not a neccessary condition. This is similar to the entropic conditions sufficient to ensure no quantum benefit from teleportation and dense coding . Analogous analysis can be applied to any other algorithm.
Finally, we point out that the states of the $`M`$ register need not be a mixture, but could be an arbitrary superposition of states $`|i_M`$ (such a state was used by Ambainis in his argument ). All the above arguments still hold in that case, and the $`M`$ and the $`C`$ registers become quantum mechanically entangled and not just classically correlated. Thus our analysis implies that any quantum computation is mathematically identical to a measurement process . The system being measured is the $`M`$ register and the apparatus is the $`C`$ register of the quantum computer. As the time progresses the apparatus (register $`C`$) becomes more and more correlated (or entangled) to the system (register $`M`$). This means that the states of register $`C`$ become more and more distinguishable which allows us to extract more information about the $`M`$ register by measuring the $`C`$ register. The analysis in the last paragraph, where we showed the limitations on the efficiency of quantum computation imposed by the mixedness of the $`C`$ register, applies also to the efficiency of a quantum measurement when the apparatus is in a mixed state. Mixedness of an apparatus, to the best of our knowledge, has never been considered in the analysis of quantum measurement. In general practice, any apparatus, however macroscopic, is considered to be in a pure quantum state before the measurement. Our approach highlighting the formal analogy between measurement and computation offers a way to analyse measurement in a much more general context.
L.R. would like to thank Invensys Plc for financial support.
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# 1 Introduction
## 1 Introduction
The aim of this paper is to proceed a bit further in search of a unified algorithm for achieving unitary and irreducible representations (unirreps for short) of Lie groups in the context of quantization. Our starting point here is a rather developed Group Approach to Quantization (GAQ) (see and references there in), which generalizes and improves Geometric Quantization (GQ) and/or the Coadjoint-Orbit Method (COM) in many respects, and particularly in the treatment of the non-Kähler orbits of the Virasoro group , denominated “non-quantizable orbits” in Ref. .
GAQ inherited, however, the technical problem of finding an appropriate and natural integration measure on the polarized submanifold of the original symplectic coadjoint orbits (or classical phase space). In fact, even though a symplectic manifold $`(M^{2n},\omega )`$ is canonicaly endowed with a volume, that is, $`\omega ^n`$, a maximally isotropic submanifold associated with a Polarization (half a symplectic manifold, so to speak) does not necessarily possess a canonical measure invariant under the action of the group generators (or quantum operators, in physical language).
Nevertheless, the virtue of GAQ working directly on a group manifold, rather than on a coadjoint orbit, taking advantage of the tools available on any Lie group (left- and right-invariant vector fields, Haar measure, etc.) brings out again the solution to the present problem of finding invariant measures. The precise technique of pseudo-extensions employed here was introduced in on an equal footing with non-trivial central extensions, and was further elaborated in , emphasizing its relation with COM. Now the main trick consists in considering pseudo-extensions by the multiplicative real line $`^+`$ along with (pseudo)-extensions by $`U(1)`$. Central (even trivial) extensions by $`^+`$ can modify well the common factor accompanying the wave functions (the weight) with an extra non-unimodular real function, thus providing half of the correction needed to make a quasi-invariant measure strictly invariant. The resulting construction shed new light on the cryptic language of “half-forms” , which came to faint the beauty of the original scheme of GQ.
This paper is organized as follows. In Sec. 2 we provide a general background on pseudo-extensions and the explicit connection with the coadjoint orbits of a general simply connected Lie group. In Sec. 3 the existence and uniqueness of a quasi-invariant measure $`\mu `$ with Radon-Nikodym derivative $`\lambda `$ on a homogeneous space is translated to the group $`G`$ itself, providing a constructive proof of the existence of such a $`\lambda `$. Then, with the aid of this function, we find a specific $`^+`$ pseudo-extension of $`G`$ making $`\mu `$ strictly invariant. The results above are applied, as an example, to the explicit construction of the representations of $`SL(2,)`$, including the Mock representation.
## 2 Pseudo-extensions
A pseudo-extension of a simply connected Lie group $`G`$ is a central extension $`\stackrel{~}{G}`$of $`G`$ by $`U(1)`$ by means of a 2-cocycle<sup>1</sup><sup>1</sup>1We shall consider, following Bargmann , local exponents $`\xi :G\times GR`$ such that $`\omega =e^{i\xi }`$ defines a 2-cocycle (or factor), $`w:G\times GU(1)`$. $`\xi _\lambda :G\times GR`$, which is a coboundary and therefore defines a trivial central extension; i.e. there exists a function $`\lambda :GR`$, the generating function of the coboundary, such that $`\xi _\lambda (g^{},g)=\lambda (g^{}g)\lambda (g^{})\lambda (g)`$, but with the property that the Lie derivative of $`\lambda `$ at the identity is different from zero for some left-invariant vector fields. In other words, the gradient of $`\lambda `$ at the identity, $`\lambda _i^0\frac{\lambda (g)}{g^i}`$, with respect to a basis of local canonical coordinates $`\{g^i\}`$ at a neighbourhood of the identity of $`G`$, is not zero.
It should be emphasized that $`\stackrel{}{\lambda }^0(\lambda _1^0,\mathrm{},\lambda _n^0)`$ defines an element in the dual $`𝒢^{}`$ of the Lie algebra $`𝒢`$ of $`G`$. Before going further into the properties of pseudo-extensions and their classification into equivalence classes (in the same way as true extensions), we must introduce some definitions.
Let $`\{X_i^L\}`$ be a basis of left-invariant vector fields associated with the canonical coordinates $`\{g^i\},i=1,\mathrm{},n=\mathrm{dim}G`$ at the identity. Let $`\{\theta ^{Li}\}`$ be the dual basis of left invariant 1-forms on $`G`$. They verify the relations:
$`i_{X_i^L}\theta ^{Lj}`$ $`=`$ $`\delta _i^j`$
$`L_{X_i^L}\theta ^{Lj}`$ $`=`$ $`C_{ik}^j\theta ^{Lk},`$ (1)
where $`C_{ik}^j`$ are the structure constants of the Lie algebra $`𝒢`$ generated by $`\{X_i^L\}`$.
Right-invariant vector fields $`\{X_i^R\}`$ can also be introduced together with the dual basis of right-invariant 1-forms $`\{\theta ^{R(i)}\}`$, satisfying properties similar to (1), but changing $`C_{ik}^j`$ by $`C_{ik}^j`$, since right-invariant vector fields generate an algebra isomorphic to that of left-invariant ones but with the structure constant with opposite sign<sup>2</sup><sup>2</sup>2This is due to our choice for the left and right action of the group on functions: $`R_g^{}f(g)=f(gg^{})`$ and $`L_g^{}f(g)=f(g^{}g)`$ instead of $`L_g^{}f(g)=f(g^{}{}_{}{}^{1}g)`$, as is used in other contexts.. Left-invariant 1-forms have zero Lie derivative with respect to right-invariant vector fields and vice versa, as it should be. An important formula which will be extensively used in this paper is the Maurer-Cartan equations:
$$d\theta ^{Li}=\frac{1}{2}C_{jk}^i\theta ^{Lj}\theta ^{Lk},$$
(2)
with analogous expression for the right-invariant counterpart, but changing the sign to the structure constants, as before. These equations state, for instance, that, for an Abelian group, all left- and right-invariant 1-forms are closed, and that left- and right-invariant 1-forms dual to vector fields that are not in the commutant of $`𝒢`$ are also closed. These properties will be relevant below.
Let us consider a central extension $`\stackrel{~}{G}`$of $`G`$ by $`U(1)`$ characterized by a 2-cocycle $`\xi :G\times GR`$, which has to satisfy the equations:
$`\xi (g_1,g_2)+\xi (g_1g_2,g_3)`$ $`=`$ $`\xi (g_1,g_2g_3)+\xi (g_2,g_3)`$
$`\xi (e,e)`$ $`=`$ $`0,`$ (3)
for all $`g_1,g_2,g_3G`$, in order to define a (associative) group law. This group law is given by:
$`g^{\prime \prime }`$ $`=`$ $`g^{}g`$
$`\zeta ^{\prime \prime }`$ $`=`$ $`\zeta ^{}\zeta e^{i\xi (g^{},g)},`$ (4)
where $`\zeta ,\zeta ^{},\zeta ^{\prime \prime }U(1)`$. Left- and right-invariant vector fields for the extended group $`\stackrel{~}{G}`$, denoted with a tilde, can be derived from the ones of $`G`$ and from the 2-cocycle as follows:
$`\stackrel{~}{X}_i^L`$ $`=`$ $`X_i^L+{\displaystyle \frac{\xi (g^{},g)}{g^i}}|_{g=e,g^{}=g}{\displaystyle \frac{}{\varphi }}`$
$`\stackrel{~}{X}_i^R`$ $`=`$ $`X_i^R+{\displaystyle \frac{\xi (g^{},g)}{g^{}^i}}|_{g^{}=e}{\displaystyle \frac{}{\varphi }},`$ (5)
where we have introduced $`\zeta =e^{i\varphi }`$. Left- and right invariant 1-forms do not change, and, of course, there are new left- and right invariant vectors fields and 1-forms associated with the new variable $`\zeta U(1)`$. These are:
$`\stackrel{~}{X}_\zeta ^L`$ $`=`$ $`{\displaystyle \frac{}{\varphi }}=2\mathrm{R}\mathrm{e}(i\zeta {\displaystyle \frac{}{\zeta }})\mathrm{\Xi }`$
$`\stackrel{~}{X}_\zeta ^R`$ $`=`$ $`\mathrm{\Xi }`$
$`\theta ^{L(\zeta )}`$ $`=`$ $`{\displaystyle \frac{d\zeta }{i\zeta }}+{\displaystyle \frac{\xi (g^{},g)}{g^i}}|_{g^{}=g^1}dg^i`$ (6)
$`\theta ^{R(\zeta )}`$ $`=`$ $`{\displaystyle \frac{d\zeta }{i\zeta }}+{\displaystyle \frac{\xi (g^{},g)}{g^i}}|_{g=g^1}dg^i,`$
where $`\frac{d\zeta }{i\zeta }=d\varphi `$. We shall call $`\mathrm{\Theta }\theta ^{L(\zeta )}`$ the Quantization 1-form. This 1-form defines a connection on the fibre bundle $`U(1)\stackrel{~}{G}G`$, and will play an important role in our formalism, since it contains all the information about the dynamics of the system under study. In fact, $`\mathrm{\Theta }_{PC}=\mathrm{\Theta }\frac{d\zeta }{i\zeta }`$ is the Poincaré-Cartan 1-form, and $`d\mathrm{\Theta }=d\mathrm{\Theta }_{PC}`$ is a presymplectic 2-form on $`G`$ which defines a symplectic 2-form once the distribution generated by its kernel is removed.
Now let us assume that we add to $`\xi `$ the coboundary $`\xi _\lambda `$, generated by the function $`\lambda `$, $`\xi _\lambda (g^{},g)=\lambda (g^{}g)\lambda (g^{})\lambda (g)`$, with $`\lambda `$ satisfying $`\lambda (e)=0`$ for $`\xi _\lambda `$ to verify (3). Then $`\xi ^{}=\xi +\xi _\lambda `$ determines a new extended group $`\stackrel{~}{G}^{}`$, and a new Quantization 1-form $`\mathrm{\Theta }^{}=\mathrm{\Theta }+\mathrm{\Theta }_\lambda `$, with
$$\mathrm{\Theta }_\lambda =\lambda _i^0\theta ^{Li}d\lambda .$$
(7)
The new presymplectic 2-form is $`d\mathrm{\Theta }^{}=d\mathrm{\Theta }+d\mathrm{\Theta }_\lambda `$, with $`d\mathrm{\Theta }_\lambda =\frac{1}{2}\lambda _i^0C_{jk}^i\theta ^{Lj}\theta ^{Lk}`$ (making use of the Maurer-Cartan equations). We shall use this decomposition of $`\mathrm{\Theta }^{}`$ and $`d\mathrm{\Theta }^{}`$ to split an arbitrary 2-cocycle $`\xi ^{}`$ in the form
$$\xi ^{}=\xi +\xi _\lambda ,$$
(8)
for some $`\lambda (g)`$. The term $`\xi `$ is such that, when considered on its own, it determines a pure central extension, i.e. a central extension for which the Lie algebra satisfies: If $`C_{ij}^\zeta 0`$, then $`C_{ij}^k=0k\zeta `$.
The term $`\xi _\lambda `$ is such that, when considered on its own, it determines a pure pseudo-extension, i.e a central extension for which the Lie algebra satisfies: $`C_{ij}^\zeta =\lambda _k^0C_{ij}^k,i,j`$, with $`\stackrel{}{\lambda }^0`$ the gradient at the identity of $`\lambda (g)`$.
An arbitrary central extension determined by $`\xi `$ will belong to a given cohomology class $`[[\xi ]]`$ constituted by all 2-cocycles $`\xi ^{}`$ differing from $`\xi `$ by coboundaries with arbitrary generating functions $`\lambda :GR`$. This is the usual definition of the $`2^{\mathrm{nd}}`$ cohomology group $`H^2(G,U(1))`$ (see, for instance ). Now we are going to introduce subclasses $`[\xi ]`$ inside $`[[\xi ]]`$, called pseudo-cohomology classes. For the sake of simplicity, we shall restrict to the trivial cohomology class $`[[\xi ]]_0`$ of 2-cocycles which admit a generating function and are therefore coboundaries. The partition of $`[[\xi ]]_0`$ into pseudo-cohomology subclasses can be translated to any other cohomology class using the decomposition (8). The equivalence relation defining the subclasses $`[\xi ]`$ is given by:
Two coboundaries $`\xi _\lambda `$ and $`\xi _\lambda ^{}`$ with generating functions $`\lambda `$ and $`\lambda ^{}`$, respectively, are in the same subclass $`[\xi ]_{\stackrel{}{\lambda }^0}`$ if and only if their gradients at the identity verify $`\stackrel{}{\lambda }^0{}_{}{}^{}=Ad^{}(g)\stackrel{}{\lambda }^0`$, for some $`gG`$.
In particular, if $`\stackrel{}{\lambda }^0{}_{}{}^{}=\stackrel{}{\lambda }^0`$, $`\xi _\lambda `$ and $`\xi _\lambda ^{}`$ are in the same pseudo-cohomology class. This allows us always to choose representatives that are linear in the canonical coordinates, $`\xi _{\stackrel{}{\lambda }^0}=\lambda _i^0g^i`$.
The condition $`\stackrel{}{\lambda }^0{}_{}{}^{}=Ad^{}(g)\stackrel{}{\lambda }^0`$ simply says that $`\stackrel{}{\lambda }^0^{}`$ and $`\stackrel{}{\lambda }^0`$ lie in the same coadjoint orbit in $`𝒢^{}`$, and it is justified because $`d\mathrm{\Theta }_{\stackrel{}{\lambda }^0^{}}`$ and $`d\mathrm{\Theta }_{\stackrel{}{\lambda }^0}`$ are symplectomorphic, the symplectomorphism being the pull-back of the coadjoint action (see ).
The equivalence relation we have just introduced constitutes a partition of the trivial cohomology class $`[[\xi ]]_0`$ of coboundaries (or of any cohomology class once translated by the relation (8)), but there is not a one to one correspondence between pseudo-cohomology classes and coadjoint orbits, since the coadjoint orbits must satisfy the integrality condition (see , and for the proof) for $`\xi _\lambda `$ to define a central extension. This restriction can be expressed in a different manner:
The gradient at the identity $`\stackrel{}{\lambda }^0𝒢^{}`$ defines a linear functional of $`𝒢`$ on $`R`$. But it also defines a one-dimensional representation of the isotropy lie subalgebra $`𝒢_{\stackrel{}{\lambda }^0}`$ of the point $`\stackrel{}{\lambda }^0`$ under the coadjoint action of $`G`$ on $`𝒢^{}`$. In particular, if $`\stackrel{}{\lambda }^0`$ is invariant under the coadjoint action (i.e. it constitutes a zero dimensional coadjoint orbit), it defines a one-dimensional representation of the whole Lie algebra $`𝒢`$. The condition of integrability of the coadjoint orbit passing through $`\stackrel{}{\lambda }^0`$ is nothing more than the condition for $`\stackrel{}{\lambda }^0`$ to be exponentiable (integrable) to a character of the isotropy subgroup $`G_{\stackrel{}{\lambda }^0}`$ (whose Lie algebra is $`𝒢_{\stackrel{}{\lambda }^0}`$).
The introduction of a pseudo-extension generated by $`\lambda (g)`$ in $`G`$, defining a central extension $`\stackrel{~}{G}`$, has the effect of modifying left- and right-invariant vector fields in the following way:
$$\stackrel{~}{X}_i^L=X_i^L+(X_i^L.\lambda \lambda _i^0)\mathrm{\Xi },\stackrel{~}{X}_i^R=X_i^R+(X_i^R.\lambda \lambda _i^0)\mathrm{\Xi }.$$
(9)
It also modifies the commutation relations in the Lie algebra $`𝒢`$ of $`G`$ (defining the commutation relations of $`\stackrel{~}{𝒢}`$):
$$[\stackrel{~}{X}_i^L,\stackrel{~}{X}_j^L]=C_{ij}^k(\stackrel{~}{X}_k^L+\lambda _k^0\mathrm{\Xi }),$$
(10)
where $`C_{ij}^k`$ are the structure constants of the original algebra $`𝒢`$. For right-invariant vector fields, we get the same commutation relations up to a sign. Once the representations of $`\stackrel{~}{G}`$have been obtained (using a technique like GAQ, for instance), we recover the representations of $`G`$ by simply redefining the operators (right-invariant vector fields) in the following manner:
$$\stackrel{~}{X}_i^R\stackrel{~}{X}_i^R{}_{}{}^{}=\stackrel{~}{X}_i^R+\lambda _i^0\mathrm{\Xi }=X_i^R+(X_i^R.\lambda )\mathrm{\Xi }.$$
(11)
It is trivial to check that the new generators $`\stackrel{~}{X}_i^R^{}`$ satisfy the (original) commutation relations of $`𝒢`$.
Once that the pseudo-extensions have been introduced and classified according to equivalence classes, they can be treated as if they were true extensions and the ordinary quantization techniques, in particular GAQ, can be applied. We refer the reader to for a detailed description of GAQ, and here we shall simply use it to arrive at the irreducible representations of $`SL(2,)`$ in Sec. 4.
## 3 Quasi-invariant measures
For any Lie group $`G`$, there exists a measure, the Haar measure, which is invariant under the left or right action of the group on itself. However, if $`M`$ is a manifold on which there is a transitive action of $`G`$ (that is, $`M`$ is a homogenous space under $`G`$), the existence of an invariant measure on $`M`$ is not guaranteed, despite that $`M`$ is locally diffeomorphic to the quotient $`G/H`$ of $`G`$ by a certain closed subgroup $`H`$, which is the isotropy group of an arbitrary point $`x_0M`$. More precisely, each point in $`M`$ has a different isotropy group, although all of them are conjugate to each other; in particular all are isomorphic.
It can be proven (see and ), however, that $`M`$ admits quasi-invariant measures. A measure $`d\mu (x)`$ on $`M`$ is called quasi-invariant if $`d\mu (gx)`$ is equivalent to $`d\mu (x)`$ for all $`gG`$, where $`gx`$ denotes the action of $`G`$ on $`M`$, and the equivalence relation is defined among measures that have the same sets of measure zero. Then the Radon-Nikodym theorem asserts that there exists a positive function $`\lambda `$ (the Radon-Nikodym derivative) on $`M`$ such that $`d\mu (gx)/d\mu (x)=\lambda (g,x)`$.
Furthermore, it turns out that any two quasi-invariant measures are equivalent (). Therefore, up to equivalence, there exists a unique quasi-invariant measure $`d\mu (x)`$ with Radon-Nikodym derivative $`\lambda (,x)`$ on $`M`$. The function $`\lambda `$ can be derived from a strictly positive, locally integrable, Borel function $`\rho (g)`$ satisfying<sup>3</sup><sup>3</sup>3Since we are considering the quotient space $`G/H`$ instead of $`H\backslash G`$, i.e. we are changing left by right with respect to , modular functions get inverted.
$$\rho (gh)=\frac{\mathrm{\Delta }_G(h)}{\mathrm{\Delta }_H(h)}\rho (g),$$
(12)
where $`\mathrm{\Delta }_G,\mathrm{\Delta }_H`$ are the modular function of $`G`$ and $`H`$, respectively (a modular function of $`G`$ is a non-negative functions on $`G`$ such that, if $`\mu _G()`$ is the left-invariant Haar measure on $`G`$, then $`\mu _G(R_gf)=\mathrm{\Delta }(g)\mu _G(f)`$, where $`R_g`$ means right translation by the element $`g`$). A modular function is a homomorphism of $`G`$ into the positive reals with the product as composition law). The Radon-Nikodym derivative is given by:
$$\lambda (g,x)=\frac{\rho (gg^{})}{\rho (g^{})},$$
(13)
where $`g^{}`$ is any element whose image under the natural projection $`GM`$ is $`x`$. This definition makes sense since $`\frac{\rho (gg^{})}{\rho (g^{})}`$ depends only on $`x`$ and not on the particular choice of $`g^{}`$.
Note that if $`\mathrm{\Delta }_H(h)=\mathrm{\Delta }_G(h),hH`$, then $`\rho (gh)=\rho (g)`$, so that we can choose $`\rho (g)=1`$ and $`\lambda (g,x)=1`$ as the Radon-Nikodym derivative. Thus, in this case, $`M`$ admits an invariant measure under $`G`$.
Let us rewrite the above considerations in infinitesimal terms. Defining the modular constants $`k_i^G\frac{\mathrm{\Delta }_G(g)}{g^i}|_{g=e},i=1,\mathrm{},n=\mathrm{dim}G`$, and similarly for $`k_i^H,i=1,\mathrm{},p=\mathrm{dim}H`$, we can rephrase the $`\rho `$-function condition (12) as:
$$X_i^L\rho (g)=k_i^{G/H}\rho (g),$$
(14)
where $`k_i^{G/H}k_i^Gk_i^H,i=1,\mathrm{},p`$. Modular constants possess properties derived from those of modular functions. Firstly, it can be proven that $`k_i^G=_{j=1}^nC_{ij}^j`$, and accordingly, $`k_i^{G/H}=_{j=p+1}^nC_{ij}^j`$, where we have assumed that the first $`p=\mathrm{dim}H`$ elements of $`𝒢`$ belong to $``$, the Lie algebra of $`H`$. In addition, $`k_i^G,i=1,\mathrm{},n`$ define a character $`k^G`$ of the Lie algebra $`𝒢`$ of $`G`$, coming from the fact that $`\mathrm{\Delta }_G(g)`$ defines a character of $`G`$, in such a way that $`k^G(X_i^L)=k_i^G`$. This property implies linearity, and also $`C_{ij}^lk_l^G=0`$, since $`k^G([X_i^L,X_j^L])=0`$. As a result, $`k_i^G=0`$ for $`G`$ semisimple.
However, $`k_i^H`$ can be non-trivial, even if $`H`$ is a subgroup of a semisimple group $`G`$, allowing for non-trivial $`k_i^{G/H}`$, and, according to (14), for the possibility of homogeneous spaces with non-invariant, although quasi-invariant, measures.
Let us develop a constructive technique for obtaining quasi-invariant measures on homogeneous spaces. That is, a procedure for constructing $`\rho `$-functions satisfying (12) (or (14)). According to Mackey , such a function always exists, although the proof of his theorem is not constructive.
Consider the left-invariant Haar measure $`\mathrm{\Omega }^L`$ on $`G`$. This is an n-form, with $`n=\mathrm{dim}G`$, and can be written, up to a constant, as $`\mathrm{\Omega }^L=\theta ^{L\mathrm{\hspace{0.17em}1}}\theta ^{L\mathrm{\hspace{0.17em}2}}\mathrm{}\theta ^{Ln}`$, where $`\theta ^{Li},i=1,\mathrm{},n`$, is the set of left invariant 1-forms on $`G`$ dual to a given basis $`\{X_i^L\}`$ of left-invariant vector fields. Let us suppose that the first $`p=\mathrm{dim}H`$ elements in these bases correspond to left-invariant 1-forms and vector fields of $`H`$, respectively. Then we tentatively define a measure on $`G/H`$ as:
$$\mathrm{\Omega }_H^L=i_{X_p^L}i_{X_{p1}^L}\mathrm{}i_{X_1^L}\mathrm{\Omega }^L=\theta ^{Lp+1}\mathrm{}\theta ^{Ln}.$$
(15)
In general, $`\mathrm{\Omega }_H^L`$ is not an invariant measure on $`G/H`$; in fact, it is not even a measure on $`G/H`$, in the sense that it does not fall down to the quotient. This can be checked by computing its invariance properties under $`X_i^L,i=1,\mathrm{},p`$. After a few computations we get $`L_{X_i^L}\mathrm{\Omega }_H^L=k_i^{G/H}\mathrm{\Omega }_H^L`$. Therefore, if $`k_i^{G/H}0`$ for some $`i`$, $`\mathrm{\Omega }_H^L`$ does not fall down to the quotient, and this is the same condition for $`G/H`$ not to have a strictly invariant measure. Therefore, these two facts seem to be related. Indeed, if we look for a function $`\rho `$ on $`G`$ such that $`L_{X_i^L}(\rho \mathrm{\Omega }_H^L)=0,i=1,\mathrm{},p`$, we find that $`\rho `$ has to be a $`\rho `$-function, satisfying $`X_i^L\rho =k_i^{G/H}\rho `$, as in (14).
Now we have to prove that equation (14) always has non-trivial solutions. We know from Mackey , that equation (12) always has a solution, but we would like to provide a proof in infinitesimal terms and, moreover, we would like to construct the solutions explicitly.
Let us consider the Radical of $``$, $`\mathrm{Rad}`$ – the maximal solvable ideal of $``$. We know that $`/\mathrm{Rad}`$ is semisimple. According to the previous considerations, the $`k_i`$’s vanish on this quotient. Thus, the non-trivial $`k_i`$’s lie only on $`\mathrm{Rad}`$, which is solvable. According to one of Lie’s theorems , a solvable algebra of operators always possesses a common eigenvector. We proceed to construct it as follows:
Let us consider the equation $`X_i^L\rho =k_i^{G/H}\rho ,i=1,\mathrm{},p`$. Let $`\chi `$ be the general solution of $`X_i^L\chi =0`$, which always exists and which we know how to construct, according to the Frobenious theorem. Then we can write $`\rho =\chi h`$, where $`h`$ is a particular solution of $`X_i^Lh=k_i^{G/H}h`$, with $`X_i^L\mathrm{Rad}`$ (the rest of the equations give zero, and since $`h`$ is a particular solution, we can choose so as not to depend on the corresponding variables). Then Lie’s Theorem guarantees the existence of such a function $`h`$, since $`\mathrm{Rad}`$ is solvable.
Once we have constructed the measure $`\rho \mathrm{\Omega }_H^L`$ on $`G/H`$, we must check its invariance properties under the action of $`G`$. For this, we compute $`L_{X_i^R}(\rho \mathrm{\Omega }_H^L)=\frac{1}{\rho }(X_i^R.\rho )(\rho \mathrm{\Omega }_H^L),i=1,\mathrm{},n`$. The result is that $`\rho \mathrm{\Omega }_H^L`$ is quasi-invariant under $`G`$ and the divergence of the vertor field $`X_i^R`$ is $`\frac{1}{\rho }(X_i^R.\rho )`$. Once the divergence of all vector fields have been computed, it is very easy to modify the (infinitesimal) action of the group $`G`$ in order to restore the invariance of $`\rho \mathrm{\Omega }_H^L`$, by defining the new vector fields:
$$\stackrel{~}{X}_i^R=X_i^R+\frac{1}{2\rho }(X_i^R.\rho ),$$
(16)
i.e., right-invariant vector fields are modified with the addition of a multiplicative term, half the divergence of the corresponding vector field. In the context of Sec. 2, we could think of this redefinition as coming from a pseudo-extension of $`G`$ by means of some pseudo-cocycle generated by a certain function $`\lambda `$ on $`G`$. In fact, this is the case, since the extra term can be written as $`X_i^R(\frac{1}{2}\mathrm{log}\rho )`$, i.e., the function $`\lambda `$, according to equation (11), would be $`\lambda =i\frac{1}{2}\mathrm{log}\rho `$. Note the presence of the imaginary constant $`i`$ in $`\lambda `$ (so that $`\lambda `$ is a pure imaginary function) revealing that $`G`$ has been centrally pseudo-extended by $`R^+`$ instead of $`U(1)`$. Therefore, the invariance of a measure on a quotient space $`G/H`$ can be restored by means of a central extension of $`G`$ by $`R^+`$ with generating function $`i\frac{1}{2}\mathrm{log}\rho `$, where $`\rho `$ is a $`\rho `$-function.
If we compute the commutation relations of the redefined vector fields, we get:
$$[\stackrel{~}{X}_i^R,\stackrel{~}{X}_j^R]=C_{ij}^k\stackrel{~}{X}_k^R,$$
(17)
showing that this pseudo-extension does not modify the commutation relations. As in Sec. 2, we can compute the gradient of the generating function $`\lambda `$ at the identity, proving to be $`\lambda _i^0=\frac{i}{2}k_i^{G/H},i=1,\mathrm{},n`$. It is pure imaginary, as would be expected of a pseudo-extension by $`R^+`$.
## 4 Example: Representations of $`SL(2,)`$
Let us consider, as an example of application of the formalism developed above, the study of the unitary and irreducible representations of $`G=SL(2,)`$. Since this group is non-simply connected, in order to apply our previous considerations, we shall consider its universal covering group $`\overline{G}`$, with $`p:\overline{G}SL(2,)`$ the covering map, which is a group homomorphism. The kernel of $`p`$ is $`Z`$, the first homotopy group of $`SL(2,)`$. It is easy to check that a unirrep $`U`$ of $`\overline{G}`$ is also a unirrep of $`SL(2,)`$ if and only if Ker$`p`$ is represented as phases, i.e $`U(g)=e^{i\alpha _g},g\mathrm{ker}p`$. Therefore, we shall compute the representations $`U`$ of $`\overline{G}`$ and then retain only the ones that verify $`U(g)=e^{i\alpha _g},\alpha _gR,g\mathrm{Ker}p`$. For simplicity, we shall denote $`\overline{G}`$ just by $`G`$, bearing in mind that at the end we wish to get the representations of $`SL(2,)`$.
Since $`SL(2,)`$ is semisimple, it has no non-trivial central extensions by $`U(1)`$; i.e. its second cohomology group $`H^2(G,U(1))=\{e\}`$. However, as shown in Sec. 2, this group admits non-trivial pseudo-extensions by $`U(1)`$, which can be classified into pseudo-cohomology classes. These pseudo-cohomology classes are in one-to-one correspondence with the coadjoint orbits of $`SL(2,)`$ with integral symplectic 2-form (see ).
Thus, we must first study the coadjoint orbits of $`SL(2,)`$. These can be classified into three types: the 1-sheet hyperboloids, the 2-sheets hyperboloids, and the cones. The cones are really three different orbits, the upper and lower cones and the origin. The origin is the only zero-dimensional orbit, and is associated with the only one-dimensional representation (character) of $`SL(2,)`$, the trivial one.
As we shall see below, the 1-sheet hyperboloids are associated with the Principal continuous series of unirreps of $`SL(2,)`$, the 2-sheet hyperboloids are associated with the Principal discrete series of unirreps and the two cones are associated with the Mock representations.
### 4.1 The group law
The $`SL(2,)`$ group can be parameterized by:
$$SL(2,)=\left\{\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)M_2()/adbc=1\right\}.$$
(18)
If $`a0`$ (the case $`a=0`$ is treated in an analogous manner, changing $`a`$ by $`c`$), we can eliminate $`d`$, $`d=\frac{1+bc}{a}`$, and we arrive at the following group law from matrix multiplication:
$`a^{\prime \prime }`$ $`=`$ $`a^{}a+b^{}c`$
$`b^{\prime \prime }`$ $`=`$ $`a^{}b+b^{}{\displaystyle \frac{1+bc}{a}}`$ (19)
$`c^{\prime \prime }`$ $`=`$ $`c^{}a+{\displaystyle \frac{1+b^{}c^{}}{a^{}}}c.`$
Left- and right-invariant vector fields are easily derived from the group law:
$$\begin{array}{ccc}X_a^L\hfill & =& a\frac{}{a}+c\frac{}{c}b\frac{}{b}\hfill \\ X_b^L\hfill & =& a\frac{}{b}\hfill \\ X_c^L\hfill & =& \frac{1+bc}{a}\frac{}{c}+b\frac{}{a}\hfill \end{array}\begin{array}{ccc}X_a^R\hfill & =& a\frac{}{a}+b\frac{}{b}c\frac{}{c}\hfill \\ X_b^R\hfill & =& \frac{1+bc}{a}\frac{}{b}+c\frac{}{a}\hfill \\ X_c^R\hfill & =& a\frac{}{c}.\hfill \end{array}$$
(20)
The Lie algebra satisfied by the (say, left-invariant) vector fields is:
$`[X_a^L,X_b^L]`$ $`=`$ $`2X_b^L`$
$`[X_a^L,X_c^L]`$ $`=`$ $`2X_c^L`$ (21)
$`[X_b^L,X_c^L]`$ $`=`$ $`X_a^L,`$
and the Casimir for this Lie algebra is given by $`\widehat{C}=\frac{1}{2}(X_a^L)^2+X_b^LX_c^L+X_c^LX_b^L`$. The left-invariant 1-forms (dual to the set of left-invariant vector fields) are given by:
$`\theta ^{L(a)}`$ $`=`$ $`{\displaystyle \frac{1+bc}{a}}dabdc`$
$`\theta ^{L(b)}`$ $`=`$ $`{\displaystyle \frac{1}{a}}db{\displaystyle \frac{b^2}{a}}dc+{\displaystyle \frac{b}{a}}{\displaystyle \frac{1+bc}{a}}da`$ (22)
$`\theta ^{L(c)}`$ $`=`$ $`adccda.`$
The exterior product of all left-invariant 1-forms constitutes a (left-invariant) volume form on the whole group (Haar measure):
$$\mathrm{\Omega }^L=\theta ^{L(a)}\theta ^{L(b)}\theta ^{L(c)}=\frac{1}{a}dadbdc.$$
(23)
### 4.2 Pseudo-extensions
The different (classes of) pseudo-extensions of $`SL(2,)`$ by $`U(1)`$ are classified, according to the discussion in Sec. 2, by the coadjoints orbits of the group $`SL(2,)`$. Let us parameterize $`𝒢^{}`$ by $`\{\alpha ,\beta ,\gamma \}`$, a coordinate system associated with the base $`\{X_a^L,X_b^L,X_c^L\}`$ of $`𝒢`$. Instead of looking for the different coadjoint orbits by direct computation, we can classify them by means of the Casimir functions. The Casimirs $`C_i`$ are invariant functions under the coadjoint action of the group on $`𝒢^{}`$, so that the equations $`C_i=c_i`$ define hypersurfaces on $`𝒢^{}`$ invariant under the coadjoint action. Of course, these hypersurfaces could be the union of two or more coadjoint orbits, and we shall need extra conditions to characterize them (these are called invariant relations, see ).
The only (independent) Casimir function for $`SL(2,)`$ is $`C=\frac{1}{2}\alpha ^2+\beta \gamma `$. This is a quadratic function, and therefore its level sets are conic sections.
It is more appropriate for our purposes to perform the change of variables $`\alpha =\alpha ,\beta =\mu +\nu ,\gamma =\mu \nu `$. In terms of the new variables, the Casimir function is written $`C=\frac{1}{2}\alpha ^2+2\mu ^22\nu ^2`$. In this form, it is easy to identify the conics, of which there are essentially three types, depending on whether $`C>0,C=0`$ or $`C<0`$. The case $`C>0`$ corresponds to 1-sheet hyperboloids; the case $`C=0`$ corresponds to the two cones and the origin, i.e. the union of three coadjoint orbits; and finally, the case $`C<0`$ corresponds to 2-sheets hyperboloids (i.e. the union of two coadjoint orbits).
Now we select a particular point $`\stackrel{}{\lambda }^0`$ in each coadjoint orbit, which will be used to define a pseudo-extension in $`SL(2,)`$ (different choices of $`\stackrel{}{\lambda }^0`$ in the same coadjoint orbit will lead to equivalent pseudo-extensions). For the case $`C>0`$, the easiest choice is $`\stackrel{}{\lambda }^0=(\alpha ,0,0)`$. For $`C=0`$, we have $`\stackrel{}{\lambda }^0=(0,0,0)`$ for the origin, and we can choose $`\stackrel{}{\lambda }^0=(0,0,\gamma <0)`$ for the upper cone and $`\stackrel{}{\lambda }^0=(0,0,\gamma >0)`$ for the lower cone. Finally, for the case $`C<0`$, we select $`\stackrel{}{\lambda }^0=(0,\beta >0,\gamma =\beta )`$ for the upper sheet and $`\stackrel{}{\lambda }^0=(0,\beta <0,\gamma =\beta )`$ for the lower sheet of the 2-sheets hyperboloid.
### 4.3 Representations associated with the 1-sheet hyperboloid: Principal Continuous Series
According to the above discussion, let us choose $`\stackrel{}{\lambda }^0=(\alpha ,0,0)`$ as the representative point in the 1-sheet hyperboloids. We need to look for a function $`\lambda `$ on $`SL(2,)`$ satisfying $`\frac{}{g^i}\lambda (g)|_{g=e}=\lambda _i^0`$. The easiest one would be a function linear on the coordinate $`a`$, but we should take into account that $`a`$ is not a canonical coordinate, since its composition law is multiplicative. That is, the uniparametric subgroup associated with it is $`R^+`$ instead of $`R`$ (the value of $`a`$ at the identity of the group is $`1`$ instead of $`0`$). Thus, we can select for $`\lambda (g)=\alpha \mathrm{log}a`$ or rather $`\lambda (g)=\alpha (a1)`$, since the generating function $`\lambda `$ must satisfy $`\lambda (e)=0`$ for $`\xi _\lambda `$ to satisfy (3).
Let us fix $`\lambda (g)=\alpha (a1)`$, to be precise (the other choice would lead to en equivalent result). The representation achieved when applying GAQ to the resulting group will be associated with the coadjoint orbit for which the Casimir is $`C=\frac{1}{2}\alpha ^2>0`$. The resulting group law for $`SL(2,)`$ pseudo-extended by $`U(1)`$ by means of the two-cocycle $`\xi _\lambda `$ is:
$`a^{\prime \prime }`$ $`=`$ $`a^{}a+b^{}c`$
$`b^{\prime \prime }`$ $`=`$ $`a^{}b+b^{}{\displaystyle \frac{1+bc}{a}}`$
$`c^{\prime \prime }`$ $`=`$ $`c^{}a+{\displaystyle \frac{1+b^{}c^{}}{a^{}}}c`$ (24)
$`\zeta ^{\prime \prime }`$ $`=`$ $`\zeta ^{}\zeta e^{i\alpha (a^{}a+b^{}ca^{}a+1)}.`$
Left- and right-invariant vector field, obtained as usual from the group law, are:
$$\begin{array}{ccc}\stackrel{~}{X}_a^L\hfill & =& a\frac{}{a}+c\frac{}{c}b\frac{}{b}+\alpha (a1)\mathrm{\Xi }\hfill \\ \stackrel{~}{X}_b^L\hfill & =& a\frac{}{b}\hfill \\ \stackrel{~}{X}_c^L\hfill & =& \frac{1+bc}{a}\frac{}{c}+b\frac{}{a}+\alpha b\mathrm{\Xi }\hfill \\ \stackrel{~}{X}_\zeta ^L\hfill & =& \frac{}{\varphi }=2\mathrm{R}\mathrm{e}(i\zeta \frac{}{\zeta })\mathrm{\Xi }\hfill \end{array}\begin{array}{ccc}\stackrel{~}{X}_a^R\hfill & =& a\frac{}{a}+b\frac{}{b}c\frac{}{c}+\alpha (a1)\mathrm{\Xi }\hfill \\ \stackrel{~}{X}_b^R\hfill & =& \frac{1+bc}{a}\frac{}{b}+c\frac{}{a}+\alpha c\mathrm{\Xi }\hfill \\ \stackrel{~}{X}_c^R\hfill & =& a\frac{}{c}\hfill \\ \stackrel{~}{X}_\zeta ^R\hfill & =& \mathrm{\Xi }.\hfill \end{array}$$
(25)
Left- and right-invariant 1-forms associated with the variables of $`SL(2,)`$ remain the same, and there are extra left- and right-invariant 1-forms associated with the variable $`\zeta `$. We are interested in the left-invariant one, which is:
$$\mathrm{\Theta }\theta ^{L(\zeta )}=\frac{d\zeta }{i\zeta }+\alpha (\theta ^{L(a)}da)=\frac{d\zeta }{i\zeta }+\alpha (\frac{1+bca}{a}dabdc).$$
(26)
The resulting Lie algebra is that of $`SL(2,)`$ with one of the commutators modified:
$`[\stackrel{~}{X}_a^L,\stackrel{~}{X}_b^L]`$ $`=`$ $`2\stackrel{~}{X}_b^L`$
$`[\stackrel{~}{X}_a^L,\stackrel{~}{X}_c^L]`$ $`=`$ $`2\stackrel{~}{X}_c^L`$ (27)
$`[\stackrel{~}{X}_b^L,\stackrel{~}{X}_c^L]`$ $`=`$ $`\stackrel{~}{X}_a^L+\alpha \mathrm{\Xi }.`$
The 2-form
$$d\mathrm{\Theta }=\alpha (dcdb+\frac{c}{a}dbda+\frac{b}{a}dcda)$$
(28)
defines a presymplectic structure on $`\stackrel{~}{G}`$. The characteristic module, or more precisely, ker$`d\mathrm{\Theta }`$ker$`\mathrm{\Theta }`$, is generated by the characteristic subalgebra, $`𝒢_C=<\stackrel{~}{X}_a^L>`$. We should remember that the characteristic subalgebra is nothing more than the isotropy subalgebra $`𝒢_{\stackrel{}{\lambda }^0}`$ of the point $`\stackrel{}{\lambda }^0𝒢`$.
Now we have to look for polarization subalgebras. These should contain the characteristic subalgebra $`𝒢_C`$ and must be horizontal (i.e., in the kernel of $`\mathrm{\Theta }`$). There are essentially two, and these lead to unitarily equivalent representations (since they are related by the adjoint action of the Lie algebra on itself, and this turns out to be a unitary transformation). We shall choose as polarization
$$𝒫=<\stackrel{~}{X}_a^L,\stackrel{~}{X}_b^L>,$$
(29)
and this, by solving the equation $`\stackrel{~}{X}_a^L\mathrm{\Psi }=\stackrel{~}{X}_b^L\mathrm{\Psi }=0`$, provides the wave functions $`\mathrm{\Psi }=\zeta e^{i\alpha (\kappa 1)}\kappa ^{i\alpha }\mathrm{\Phi }(\tau )`$, where $`\kappa a`$ and $`\tau \frac{c}{a}`$. The action of the right-invariant vector fields on polarized wave functions is:
$`\stackrel{~}{X}_a^R\mathrm{\Psi }`$ $`=`$ $`\zeta e^{i\alpha (\kappa 1)}\kappa ^{i\alpha }[2\tau {\displaystyle \frac{d}{d\tau }}]\mathrm{\Phi }(\tau )`$
$`\stackrel{~}{X}_b^R\mathrm{\Psi }`$ $`=`$ $`\zeta e^{i\alpha (\kappa 1)}\kappa ^{i\alpha }[i\alpha \tau \tau ^2{\displaystyle \frac{d}{d\tau }}]\mathrm{\Phi }(\tau )`$ (30)
$`\stackrel{~}{X}_c^R\mathrm{\Psi }`$ $`=`$ $`\zeta e^{i\alpha (\kappa 1)}\kappa ^{i\alpha }[{\displaystyle \frac{d}{d\tau }}]\mathrm{\Phi }(\tau ).`$
According to Sec. 2, the right-invariant generators should be redefined as $`\stackrel{~}{X}_{g^i}^R\stackrel{~}{X}_{g^i}^R{}_{}{}^{}=\stackrel{~}{X}_{g^i}^R+\lambda _i^0\mathrm{\Xi }`$ in order to obtain the representations of $`G`$, and this affects only the generators $`\stackrel{~}{X}_a^R`$, which changes to $`\stackrel{~}{X}_a^R{}_{}{}^{}=\stackrel{~}{X}_a^R+\alpha \mathrm{\Xi }`$. Its action on polarized wave functions turns out to be:
$$\stackrel{~}{X}_a^R{}_{}{}^{}\mathrm{\Psi }=\zeta e^{i\alpha (\kappa 1)}\kappa ^{i\alpha }[i\alpha 2\tau \frac{d}{d\tau }]\mathrm{\Phi }(\tau ).$$
(31)
The representation of $`SL(2,)`$ here constructed is irreducible but not unitary. One way of viewing it (before discussing integration measures) is to consider the Casimir operator, which is the quadratic operator $`\widehat{C}=\frac{1}{2}(\stackrel{~}{X}_a^R)^2+\stackrel{~}{X}_b^R\stackrel{~}{X}_a^R+\stackrel{~}{X}_c^R\stackrel{~}{X}_b^R`$. After the pseudoextension and redefinition of operators ($`\stackrel{~}{X}_a^R`$ should be changed by $`\stackrel{~}{X}_a^R^{}`$), the resulting Casimir operator, $`\widehat{C}^{}`$, acts on polarized wave functions as $`\widehat{C}{}_{}{}^{}\mathrm{\Psi }=(\alpha ^2/2+i\alpha )\mathrm{\Psi }`$. The fact that it is a number reveals that the representation is irreducible, but since it is not real, the representation cannot be unitary (the Casimir is a quadratic function of (anti-)Hermitian operators, and should therefore be a self-adjoint operator in any unitary representation).
The reason for this lack of unitarity is that the support manifold for the representation does not admit an invariant measure. Since the process of polarizing wave functions really amounts to reducing the space of functions to those defined in the quotient $`G/G_𝒫`$, where $`G_𝒫`$ is the group associated with the polarization subalgebra $`𝒫`$), the support manifold is given by $`G/G_𝒫`$, which is naturally an homogeneous space under $`G`$. According to Sec. 3, it may well happen that $`G/G_𝒫`$ does not admit an invariant measure, and in fact this is the case. However, the existence of quasi-invariant measures is granted, and this fact will allow us to restore the unitarity of the representation.
If we compute the measure on $`G/G_𝒫`$, derived from the left Haar measure $`\mathrm{\Omega }^L`$ on $`G`$, we obtain $`\mathrm{\Omega }_𝒫^L=i_{\stackrel{~}{X}_b^L}i_{\stackrel{~}{X}_a^L}\mathrm{\Omega }^L=adccda`$. When expressed in terms of the new variables $`\kappa `$ and $`\tau `$, it takes the form $`\mathrm{\Omega }_𝒫^L=\kappa ^2d\tau `$. Taking into account that $`G/G_𝒫`$ is parameterized by $`\tau `$, now becomes clear why the representation is not unitary: the measure does not even fall down to the quotient.
A solution to this problem consists in choosing any quasi-invariant measure on $`G/G_𝒫`$ and introducing the appropriate Radon-Nikodym derivative . Here, we propose another, yet equivalent, solution to this lack of unitarity, giving a new insight into the problem according to Sec. 3. We shall consider a pseudo-extension of $`G`$ by $`R^+`$, rather than $`U(1)`$. The reason is that we wish to restore the unitarity of a non-unitary representation, and for this we need a “piece” of non-unitary representation, in such a way that the resulting representation is unitary. To enable a direct comparison with the treatment of Mackey, we shall employ the equivalent technique of non-horizontal polarizations instead of that of pseudo-extensions. A non-horizontal polarization $`𝒫^{\mathrm{n}.\mathrm{h}.}`$ is a polarization in which the horizontality condition has been relaxed. The polarization equations acquire the form: $`\stackrel{~}{X}_j^L\mathrm{\Psi }=i\alpha _j\mathrm{\Psi },\stackrel{~}{X}_j^L𝒫^{\mathrm{n}.\mathrm{h}.}`$ (see for a discussion on the equivalence between pseudo-extensions and non-horizontal polarizations).
The key point is to keep $`\mathrm{\Omega }_𝒫^L`$ as the measure on $`G/G_𝒫`$, and to impose the polarization conditions $`\stackrel{~}{X}_i^L\stackrel{~}{\mathrm{\Psi }}=\frac{1}{2}k_i^{G/G_𝒫}\stackrel{~}{\mathrm{\Psi }}`$, instead of $`\stackrel{~}{X}_i^L\mathrm{\Psi }=0,\stackrel{~}{X}_i^L𝒫`$. In finite terms, this condition is written as:
$$\stackrel{~}{\mathrm{\Psi }}(gh)=\sqrt{\frac{\mathrm{\Delta }_G(h)}{\mathrm{\Delta }_H(h)}}\stackrel{~}{\mathrm{\Psi }}(g).$$
(32)
We can rephrase this by saying that $`\stackrel{~}{\mathrm{\Psi }}`$ is a $`\frac{1}{2}`$-$`\rho `$-function<sup>4</sup><sup>4</sup>4Note that, according to Sec. 3, the generating function for the pseudo-extension by $`R^+`$ would be $`\lambda =\frac{i}{2}\mathrm{log}\rho =i\mathrm{log}\rho ^{\frac{1}{2}}`$, with $`\lambda _i^0=\frac{i}{2}k_i^{G/G_𝒫}=\alpha _i`$.. The purpose of this definition is to make $`\stackrel{~}{\mathrm{\Psi }}^{}\stackrel{~}{\mathrm{\Psi }}^{}`$ a $`\rho `$-function, with two $`\frac{1}{2}`$-$`\rho `$-functions $`\stackrel{~}{\mathrm{\Psi }}`$ and $`\stackrel{~}{\mathrm{\Psi }^{}}`$, in such a way that $`\stackrel{~}{\mathrm{\Psi }}^{}\stackrel{~}{\mathrm{\Psi }}^{}\mathrm{\Omega }_𝒫^L`$ is a well-defined quantity on $`G/G_𝒫`$ and can be integrated with respect to $`\tau `$. In other words, $`\stackrel{~}{\mathrm{\Psi }}^{}\stackrel{~}{\mathrm{\Psi }}^{}`$ is a $`\rho `$-function necessary to make $`\mathrm{\Omega }_𝒫^L`$ a quasi-invariant measure on $`G/G_𝒫`$.
To begin, we must compute the modular constants $`k_i^{G/G_𝒫}=k_i^Gk_i^{G_𝒫},i=1,\mathrm{},p`$. Firstly, since $`G=SL(2,)`$ is semi-simple, $`k_i^G=0,i=1,\mathrm{},n`$. Secondly, we have $`k_a^{G_𝒫}=2`$ and $`k_b^{G_𝒫}=0`$. Therefore, $`k_a^{G/G_𝒫}=2`$ and $`k_b^{G/G_𝒫}=0`$.
Accordingly, the new polarization equations we have to solve are:
$$\stackrel{~}{X}_a^L\stackrel{~}{\mathrm{\Psi }}=\stackrel{~}{\mathrm{\Psi }},\stackrel{~}{X}_b^L\stackrel{~}{\mathrm{\Psi }}=0.$$
(33)
It is easy to verify that the solutions of these new polarization equations are of the form:
$$\stackrel{~}{\mathrm{\Psi }}(g)=a^1\mathrm{\Psi }(g),$$
(34)
where $`\mathrm{\Psi }(g)`$ is a solution of the previous (horizontal) polarization equations. Thus, the form of the solutions is:
$$\stackrel{~}{\mathrm{\Psi }}=\zeta \kappa ^1e^{i\alpha (\kappa 1)}\kappa ^{i\alpha }\mathrm{\Phi }(\tau ).$$
(35)
Now it it clear why $`\stackrel{~}{\mathrm{\Psi }}^{}\stackrel{~}{\mathrm{\Psi }}^{}\mathrm{\Omega }_𝒫^L=\mathrm{\Phi }(\tau )^{}\mathrm{\Phi }^{}(\tau )d\tau `$ can be integrated in $`G/G_𝒫`$; the $`\kappa `$ dependence has been removed.
The right-invariant vector fields, when acting on $`\frac{1}{2}`$-$`\rho `$-functions, acquire extra terms that restore the unitarity of the representation<sup>5</sup><sup>5</sup>5The difference between pseudo-extensions and non-horizontal polarizations lie in the fact that pseudo-extensions modify the left- and right-invariant vector fields and non-horizontal polarizations modify the wave functions. The extra term in the reduced operators is a consequence of their acting on modified wave functions.:
$$\stackrel{~}{X}_i^R\stackrel{~}{\mathrm{\Psi }}=\kappa ^1\stackrel{~}{X}_i^R\mathrm{\Psi }+\kappa ^1(\kappa \stackrel{~}{X}_i^R.\kappa ^1)\mathrm{\Psi }.$$
(36)
In this way, the final representation has the form, restricted to its action on $`\mathrm{\Phi }(\tau )`$:
$`\stackrel{~}{X}_a^R{}_{}{}^{}\mathrm{\Phi }(\tau )`$ $`=`$ $`[1+i\alpha 2\tau {\displaystyle \frac{d}{d\tau }}]\mathrm{\Phi }(\tau )`$
$`\stackrel{~}{X}_b^R\mathrm{\Phi }(\tau )`$ $`=`$ $`[\tau +i\alpha \tau \tau ^2{\displaystyle \frac{d}{d\tau }}]\mathrm{\Phi }(\tau )`$ (37)
$`\stackrel{~}{X}_c^R\mathrm{\Phi }(\tau )`$ $`=`$ $`[{\displaystyle \frac{d}{d\tau }}]\mathrm{\Phi }(\tau ).`$
W can readly verify that these operators are self-adjoint with respect to the quasi-invariant measure $`d\tau `$ (what remains of $`\mathrm{\Omega }_𝒫^L`$ after multiplication by the factor $`\kappa ^2`$ contained in the wave functions). Even more, the Casimir operator, acting on the new wave functions, turns out to be real, revealing that the representation is now unitary:
$$\widehat{C}^{}\mathrm{\Phi }(\tau )=\frac{1}{2}(1+\alpha ^2)\mathrm{\Phi }(\tau ).$$
(38)
### 4.4 Representations associated with the cones: Mock representation
In accordance with Sec. 4.2, let us choose $`\stackrel{}{\lambda }^0=(0,0,\gamma )`$ as the representative point in the cone. If $`\gamma <0`$ we are in the upper cone and if $`\gamma >0`$ we are in the lower cone. We have to look for a function $`\lambda `$ on $`SL(2,)`$ satisfying $`\frac{}{g^i}\lambda (g)|_{g=e}=\lambda _i^0`$. The easiest one is the function linear on the coordinate $`c`$, since here $`c`$ is a true canonical coordinate, and therefore, we fix $`\lambda (g)=\gamma c`$.
The representation obtained when applying GAQ to the resulting group will be associated with one of the coadjoint orbit for which the Casimir is $`C=0`$. The resulting group law for $`SL(2,)`$ pseudo-extended by $`U(1)`$ by means of the two-cocycle $`\xi _\lambda `$ is:
$`a^{\prime \prime }`$ $`=`$ $`a^{}a+b^{}c`$
$`b^{\prime \prime }`$ $`=`$ $`a^{}b+b^{}{\displaystyle \frac{1+bc}{a}}`$
$`c^{\prime \prime }`$ $`=`$ $`c^{}a+{\displaystyle \frac{1+b^{}c^{}}{a^{}}}c`$ (39)
$`\zeta ^{\prime \prime }`$ $`=`$ $`\zeta ^{}\zeta e^{i\gamma (c^{}a+\frac{1+b^{}c^{}}{a^{}}cc^{}c)}.`$
Left- and right-invariant vector field, derived as usual from the group law, are:
$$\begin{array}{ccc}\stackrel{~}{X}_a^L\hfill & =& a\frac{}{a}+c\frac{}{c}b\frac{}{b}+\gamma c\mathrm{\Xi }\hfill \\ \stackrel{~}{X}_b^L\hfill & =& a\frac{}{b}\hfill \\ \stackrel{~}{X}_c^L\hfill & =& \frac{1+bc}{a}\frac{}{c}+b\frac{}{a}+\gamma (\frac{1+bc}{a}1)\mathrm{\Xi }\hfill \\ \stackrel{~}{X}_\zeta ^L\hfill & =& \frac{}{\varphi }=2\mathrm{R}\mathrm{e}(i\zeta \frac{}{\zeta })\mathrm{\Xi }\hfill \end{array}\begin{array}{ccc}\stackrel{~}{X}_a^R\hfill & =& a\frac{}{a}+b\frac{}{b}c\frac{}{c}\gamma c\mathrm{\Xi }\hfill \\ \stackrel{~}{X}_b^R\hfill & =& \frac{1+bc}{a}\frac{}{b}+c\frac{}{a}\hfill \\ \stackrel{~}{X}_c^R\hfill & =& a\frac{}{c}+\gamma (a1)\mathrm{\Xi }\hfill \\ \stackrel{~}{X}_\zeta ^R\hfill & =& \mathrm{\Xi }.\hfill \end{array}$$
(40)
The left-invariant 1-form associated with the variable $`\zeta `$ is:
$$\mathrm{\Theta }\theta ^{L(\zeta )}=\frac{d\zeta }{i\zeta }+\gamma (\theta ^{L(c)}dc)=\frac{d\zeta }{i\zeta }+\gamma ((a1)dccda).$$
(41)
The resulting Lie algebra is, again, that of $`SL(2,)`$ with one of the commutators modified, in this case the one giving $`\stackrel{~}{X}_c^L`$ on the r.h.s.:
$`[\stackrel{~}{X}_a^L,\stackrel{~}{X}_b^L]`$ $`=`$ $`2\stackrel{~}{X}_b^L`$
$`[\stackrel{~}{X}_a^L,\stackrel{~}{X}_c^L]`$ $`=`$ $`2(\stackrel{~}{X}_c^L+\gamma \mathrm{\Xi })`$ (42)
$`[\stackrel{~}{X}_b^L,\stackrel{~}{X}_c^L]`$ $`=`$ $`\stackrel{~}{X}_a^L.`$
The 2-form
$$d\mathrm{\Theta }=2\gamma dadc$$
(43)
defines a presymplectic structure on $`\stackrel{~}{G}`$. The characteristic subalgebra is $`𝒢_C=<\stackrel{~}{X}_b^L>`$. In this case, there is essentially one polarization, given by:
$$𝒫=<\stackrel{~}{X}_b^L,\stackrel{~}{X}_a^L>,$$
(44)
and this provides, by solving the equation $`\stackrel{~}{X}_a^L\mathrm{\Psi }=\stackrel{~}{X}_b^L\mathrm{\Psi }=0`$, the wave functions $`\mathrm{\Psi }=\zeta e^{i\gamma c}\mathrm{\Phi }(\tau )`$, where again $`\tau \frac{c}{a}`$. The action of right-invariant vector fields on polarized wave functions is:
$`\stackrel{~}{X}_a^R\mathrm{\Psi }`$ $`=`$ $`\zeta e^{i\gamma c}[2\tau {\displaystyle \frac{d}{d\tau }}]\mathrm{\Phi }(\tau )`$
$`\stackrel{~}{X}_b^R\mathrm{\Psi }`$ $`=`$ $`\zeta e^{i\gamma c}[\tau ^2{\displaystyle \frac{d}{d\tau }}]\mathrm{\Phi }(\tau )`$ (45)
$`\stackrel{~}{X}_c^R\mathrm{\Psi }`$ $`=`$ $`\zeta e^{i\gamma c}[{\displaystyle \frac{d}{d\tau }}i\gamma ]\mathrm{\Phi }(\tau ).`$
The redefinition of the right-invariant generators $`\stackrel{~}{X}_{g^i}^R\stackrel{~}{X}_{g^i}^R{}_{}{}^{}=\stackrel{~}{X}_{g^i}^R+\lambda _i^0\mathrm{\Xi }`$ in order to obtain the representation of $`G`$, affects only to the $`\stackrel{~}{X}_c^R`$ generator, which changes to $`\stackrel{~}{X}_c^R{}_{}{}^{}=\stackrel{~}{X}_a^R+\gamma \mathrm{\Xi }`$. Its action on polarized wave functions turns out to be:
$$\stackrel{~}{X}_c^R{}_{}{}^{}\mathrm{\Psi }=\zeta e^{i\gamma c}[\frac{d}{d\tau }]\mathrm{\Phi }(\tau ).$$
(46)
The representation of $`SL(2,)`$ here constructed, as in the case of the 1-sheet hyperboloid, is irreducible but not unitary.
The reason for this lack of unitarity is the same as before, that is, the lack of an invariant measure on the support manifold $`G/G_𝒫`$. In fact, the polarization $`𝒫`$ is the same as in the case of the 1-sheet hyperboloid, only the vector fields are slightly different, since they come from different pseudo-extensions. Therefore, the wave functions are essentially the same as before, and consequently $`G/G_𝒫`$ is the same as in the case of the 1-sheet hyperboloid.
The measure on $`G/G_𝒫`$ is again $`\mathrm{\Omega }_𝒫^L=i_{\stackrel{~}{X}_b^L}i_{\stackrel{~}{X}_a^L}\mathrm{\Omega }^L=adccda=\kappa ^2d\tau `$, which does not fall down to the quotient.
Thus, we keep $`\mathrm{\Omega }_𝒫^L`$ as the measure on $`G/G_𝒫`$, and we impose the polarization conditions $`\stackrel{~}{X}_i^L\stackrel{~}{\mathrm{\Psi }}=\frac{1}{2}k_i^{G/G_𝒫}\stackrel{~}{\mathrm{\Psi }}`$, instead of $`\stackrel{~}{X}_i^L\mathrm{\Psi }=0,\stackrel{~}{X}_i^L𝒫`$. In other words, we impose $`\stackrel{~}{\mathrm{\Psi }}`$ to be a $`\frac{1}{2}`$-$`\rho `$-function in such a way that $`\stackrel{~}{\mathrm{\Psi }}^{}\stackrel{~}{\mathrm{\Psi }}^{}`$ is a $`\rho `$-function, $`\stackrel{~}{\mathrm{\Psi }}`$ and $`\stackrel{~}{\mathrm{\Psi }^{}}`$ being two $`\frac{1}{2}`$-$`\rho `$-functions. Now, $`\stackrel{~}{\mathrm{\Psi }}^{}\stackrel{~}{\mathrm{\Psi }}^{}\mathrm{\Omega }_𝒫^L`$ is a well-defined quantity on $`G/G_𝒫`$ and can be integrated with respect to $`\tau `$.
Modular constants $`k_i^{G/G_𝒫}=k_i^Gk_i^{G_𝒫},i=1,\mathrm{},p`$, are the same as before, since $`G_𝒫`$ is the same group. Therefore, $`k_a^{G/G_𝒫}=2`$ and $`k_b^{G/G_𝒫}=0`$.
The new polarization equations are:
$$\stackrel{~}{X}_a^L\stackrel{~}{\mathrm{\Psi }}=\stackrel{~}{\mathrm{\Psi }},\stackrel{~}{X}_b^L\stackrel{~}{\mathrm{\Psi }}=0.$$
(47)
with solutions:
$$\stackrel{~}{\mathrm{\Psi }}(g)=a^1\mathrm{\Psi }(g),$$
(48)
where $`\mathrm{\Psi }(g)`$ is a solution of the previous (horizontal) polarization equations. Thus, the form of the solutions is:
$$\stackrel{~}{\mathrm{\Psi }}=\zeta \kappa ^1e^{i\gamma \kappa \tau }\mathrm{\Phi }(\tau ).$$
(49)
The right-invariant vector fields, when acting on $`\frac{1}{2}`$-$`\rho `$-functions, acquire extra terms restoring the unitarity of the representation:
$$\stackrel{~}{X}_i^R\stackrel{~}{\mathrm{\Psi }}=\kappa ^1\stackrel{~}{X}_i^R\mathrm{\Psi }+\kappa ^1(\kappa \stackrel{~}{X}_i^R.\kappa ^1)\mathrm{\Psi }.$$
(50)
This way, the final representation restricted to its action on $`\mathrm{\Phi }(\tau )`$ has the form :
$`\stackrel{~}{X}_a^R\mathrm{\Phi }(\tau )`$ $`=`$ $`[12\tau {\displaystyle \frac{d}{d\tau }}]\mathrm{\Phi }(\tau )`$
$`\stackrel{~}{X}_b^R\mathrm{\Phi }(\tau )`$ $`=`$ $`[\tau \tau ^2{\displaystyle \frac{d}{d\tau }}]\mathrm{\Phi }(\tau )`$ (51)
$`\stackrel{~}{X}_c^R{}_{}{}^{}\mathrm{\Phi }(\tau )`$ $`=`$ $`[{\displaystyle \frac{d}{d\tau }}]\mathrm{\Phi }(\tau ).`$
Again, we can readily verify that these operators are self-adjoint with respect to the quasi-invariant measure $`d\tau `$ (what remains of $`\mathrm{\Omega }_𝒫^L`$ after multiplication by the factor $`\kappa ^2`$ contained in the wave functions). Therefore, the representation is now unitary.
This representation can be seen as the limit $`\alpha 0`$ of the Principal series of representations. We should stress at this point that the representation does not depend on $`\gamma `$, nor even on its sign. Therefore, we obtain the same representation for both cones, which are clearly equivalent. The reason for this equivalence is that the group isomorphism $`(a,b,c)(a,b,c)`$ induces a unitary transformation between the two representations. This representation (up to equivalence) is called the Mock representation and is associated with the two cones.
### 4.5 Representations associated with the 2-sheets hyperboloids: Discrete Series
According to Sec. 4.2, we can choose the point $`\stackrel{}{\lambda }^0=(0,\beta >0,\gamma =\beta )`$ in the upper sheet and $`\stackrel{}{\lambda }^0=(0,\beta <0,\gamma =\beta )`$ in the lower sheet of the 2-sheets hyperboloid, to define the pseudo-extension of $`SL(2,)`$ by $`U(1)`$. Let us consider $`\stackrel{}{\lambda }^0=(0,\beta ,\beta )`$, keeping the sign of $`\beta `$ undetermined for the time being.
The easiest function $`\lambda `$ on $`SL(2,)`$ satisfying $`\frac{}{g^i}\lambda (g)|_{g=e}=\lambda _i^0`$ is the function linear on the coordinate $`(bc)`$, since here $`b`$ and $`c`$ are true canonical coordinates. Therefore, we fix $`\lambda (g)=\beta (bc)`$.
The representation obtained when applying GAQ to the resulting group will be associated with one of the coadjoint orbits for which the Casimir is $`C=\beta ^2<0`$. The resulting group law for $`SL(2,)`$, pseudo-extended by $`U(1)`$ by means of the two-cocycle $`\xi _\lambda `$, is:
$`a^{\prime \prime }`$ $`=`$ $`a^{}a+b^{}c`$
$`b^{\prime \prime }`$ $`=`$ $`a^{}b+b^{}{\displaystyle \frac{1+bc}{a}}`$
$`c^{\prime \prime }`$ $`=`$ $`c^{}a+{\displaystyle \frac{1+b^{}c^{}}{a^{}}}c`$ (52)
$`\zeta ^{\prime \prime }`$ $`=`$ $`\zeta ^{}\zeta e^{i\beta ((a^{}1)b(a1)c^{}+\frac{1+bca}{a}b^{}\frac{1+b^{}c^{}a^{}}{a^{}}c)}.`$
Left- and right-invariant vector field are:
$$\begin{array}{ccc}\stackrel{~}{X}_a^L\hfill & =& a\frac{}{a}+c\frac{}{c}b\frac{}{b}\beta (b+c)\mathrm{\Xi }\hfill \\ \stackrel{~}{X}_b^L\hfill & =& a\frac{}{b}+\beta (a1)\mathrm{\Xi }\hfill \\ \stackrel{~}{X}_c^L\hfill & =& \frac{1+bc}{a}\frac{}{c}+b\frac{}{a}\beta (\frac{1+bca}{a})\mathrm{\Xi }\hfill \\ \stackrel{~}{X}_\zeta ^L\hfill & =& \mathrm{\Xi }\hfill \end{array}\begin{array}{ccc}\stackrel{~}{X}_a^R\hfill & =& a\frac{}{a}+b\frac{}{b}c\frac{}{c}+\beta (b+c)\mathrm{\Xi }\hfill \\ \stackrel{~}{X}_b^R\hfill & =& \frac{1+bc}{a}\frac{}{b}+c\frac{}{a}+\beta (\frac{1+bca}{a})\mathrm{\Xi }\hfill \\ \stackrel{~}{X}_c^R\hfill & =& a\frac{}{c}\beta (a1)\mathrm{\Xi }\hfill \\ \stackrel{~}{X}_\zeta ^R\hfill & =& \mathrm{\Xi }.\hfill \end{array}$$
(53)
The left-invariant 1-form associated with the variable $`\zeta `$ is:
$`\mathrm{\Theta }\theta ^{L(\zeta )}`$ $`=`$ $`{\displaystyle \frac{d\zeta }{i\zeta }}+\beta (\theta ^{L(b)}db\theta ^{L(c)}+dc)={\displaystyle \frac{d\zeta }{i\zeta }}+\beta [{\displaystyle \frac{1a}{a}}db(1+a+{\displaystyle \frac{b^2}{a}})dc+`$ (54)
$`({\displaystyle \frac{b}{a^2}}(1+bc)c)da].`$
The resulting Lie algebra is, as in the other cases, the one of $`SL(2,)`$ with some of the commutators modified, in this case those giving $`\stackrel{~}{X}_b^L`$ and $`\stackrel{~}{X}_c^L`$ on the r.h.s.:
$`[\stackrel{~}{X}_a^L,\stackrel{~}{X}_b^L]`$ $`=`$ $`2(\stackrel{~}{X}_b^L+\beta \mathrm{\Xi })`$
$`[\stackrel{~}{X}_a^L,\stackrel{~}{X}_c^L]`$ $`=`$ $`2(\stackrel{~}{X}_c^L\beta \mathrm{\Xi })`$ (55)
$`[\stackrel{~}{X}_b^L,\stackrel{~}{X}_c^L]`$ $`=`$ $`\stackrel{~}{X}_a^L.`$
The 2-form defining a presymplectic structure on $`\stackrel{~}{G}`$is
$$d\mathrm{\Theta }=2\beta \left[\frac{b}{a}dbdc+\frac{1+bc}{a^2}dadb+dadc\right].$$
(56)
The characteristic subalgebra turns out to be $`𝒢_C=<\stackrel{~}{X}_b^L\stackrel{~}{X}_c^L>`$. Looking for a polarization subalgebra containing the characteristic subalgebra, we get into trouble, since there is no such real subalgebra. We are forced to complexify the algebra, and then we find (essentially) two complex polarizations:
$$𝒫=<\stackrel{~}{X}_b^L\stackrel{~}{X}_c^L,\stackrel{~}{X}_b^L+\stackrel{~}{X}_c^L\pm i\stackrel{~}{X}_a^L>.$$
(57)
Clearly, the solution to these polarization equations are complex functions defined on a complex submanifold of the complexification of $`SL(2,)`$. These will be holomorphic or anti-holomorphic, depending on the choice of sign in (57). The explicit construction of the representations in the discrete series, according to the group quantization framework, was firstly given in Ref. in connection to the quantum dynamics of a free particle on Anti-de Sitter space-time. Higher-order, real polarizations were used in Ref. in the study of the relativistic harmonic oscillator. They have also been considered in conformal field theory as factor of $`SO(2,2)SL(2,)SL(2,)`$ representations .
## Acknowledgements
We thank G. Marmo for very useful comments on Sec. 3.
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# Gravitational energy of simple bodies: the method of negative density
##
### Introduction
In Newtonian theory of gravitation, there are suprisingly few exact solutions for the proper energy of the homogeneous self-gravitating bodies: a) ellipsoid (including sphere as a particular case) , b) concave bispherical lens (including spherical segment as a particular case) (KA), and c) rectangular parallelepiped (including cube as a particular case) . We present here a method to obtain the formulas for potential energy of homogeneous bodies. The method includes a notion of negative density. Here we use this method, first, to check the formulas of KA and, second, to get new formula for the homogeneous concavo-convex lens. To this end, we widely used Mathematica .
### Homogeneous sphere and segment
Here we use the method of negative density (MND) to check the formula by KA for potential energy of homogeneous spherical segment.
In Fig. 1 we show the setting of the problem. First, we have the homogeneous sphere with radius $`R=OA=OC=OD`$ and density $`\rho _1`$. Gravitational potential, $`\phi `$, inside the homogeneous sphere, as function of radius $`r`$ is:
$$\phi =2\pi G\rho _1(R^2r^2/3);$$
(1)
and the gravitational energy of the sphere is:
$$W_{sphere}=\frac{1}{2}\underset{M}{}\phi 𝑑m=\frac{16}{15}\pi ^2GR^5\rho _1^2.$$
(2)
Now we put into the sphere the additional spherical segment with density $`\rho _2`$, the same radius $`R`$, with height $`h=AB`$ and radius of base $`a=BC`$. Note that the total matter density in the region occupied by this segment is $`\rho _1+\rho _2`$.
Potential energy of interaction of the sphere and the additional segment is obtained by integration of $`\phi \rho _2`$ over the volume of the segment:
$$W_{int}=\underset{V}{}\phi \rho _2𝑑v=\underset{Rh}{\overset{R}{}}\phi [2\pi \rho _2r(rR+h)]𝑑r;$$
(3)
we have:
$$W_{int}=\frac{Gh^2\pi ^2(h^35h^2R+20R^3)\rho _1\rho _2}{15}.$$
(4)
According to KA, the gravitational energy of the homogeneous spherical segment with radius $`R`$, height $`h`$, radius of base $`a`$ and density $`\rho _2`$ is ($`a=(2Rhh^2)^{1/2}`$):
$$W_{segment}(a,h,R,\rho _2)=\frac{1}{9}\pi G\rho _2^2[\frac{3}{2}\pi h^4(\frac{h}{5}R)+8[3R^4aR^2a^3\frac{2}{5}a^56R^4(Rh)\mathrm{arctan}\frac{h}{a}]].$$
(5)
If this formula is correct, then the gravitational energy of the homogeneous spherical segment with radius $`R`$, height $`2Rh`$, base’s radius $`a`$ and density $`\rho _1`$ is: $`W_{ad.seg}=W_{segment}(a,2Rh,R,\rho _1).`$ We shall show that this is indeed correct.
We may look at the body in the Fig. 1 in twofold way: as
a)homogeneous sphere plus homogeneous segment, or as
b)two homogeneous segments with same $`a`$ and $`R`$ but with different heights and densities.
We can not calculate the total gravitational energy of the body in the case b) as we do not know an external potential of spherical segment. However, we can calculate the total gravitational energy of the body in the case a), $`W_a`$, as a sum of three terms:
the proper gravitational energy of homogeneous sphere, Eq. (2)
the proper gravitational energy of homogeneous segment, Eq. (5), and
the potential energy of interaction of sphere and segment, Eq. (4).
$$W_a=W_{sphere}+W_{segment}+W_{int}.$$
(6)
Now the essence of the negative density method: we take $`\rho _2=\rho _1`$!
Then, turning to the case b), the body consisted of two segments gives us one (right) ”segment” with zero density, and the other (left) homogeneous segment with ”ordinary” positive density $`\rho _1`$. In fact, we have only one (left) segment with parameters $`a,\mathrm{\hspace{0.17em}2}Rh,R,\rho _1`$. It means that if we take $`\rho _2=\rho _1`$, we should get , from Eq. (6), the proper potential energy of ”left” segment.
As a result, we have the identity for formula of potential energy of homogeneous spherical segment:
$$W_{segment}(a,2Rh,R,\rho )W_{segment}(a,h,R,\rho )=\frac{G\pi ^2\left(hR\right)\left(h^2+2hR+4R^2\right)^2\rho ^2}{15},$$
(7)
which is correct if we take into account the evident relation: $`\mathrm{arctan}\frac{h}{a}+\mathrm{arctan}\frac{2Rh}{a}=\frac{\pi }{2}`$. This is a rather rigorous check of the formula (5) for potential energy of spherical segment. We conclude that the method of negative density is a powerful method of checking the sophisticated formula of gravitational energy of homogeneous spherical segment.
However our MND may be also applied to check the even more complex formulas by KA for the proper potential energy of symmetric and asymmetric concave lenses, and also to obtain the formula for potential energy of concavo-convex lens.
### Homogeneous sphere and asymmetric convex lens
In Fig. 2 we show the setting of the problem. First, we have the homogeneous sphere with radius $`R=OA=OC=OD`$ and density $`\rho _1`$. Then, we put into this sphere the homogeneous, with density $`\rho _2`$, asymmetrical bispherical lens, comprised of two homogeneous spherical segments: one (right) segment of the same radius $`R`$, of height $`h=AB`$ and radius of base $`a=BC`$; the other (left) segment of radius $`R_1>R`$, of height $`h_1=BE<h=AB`$ and the same radius of base $`a=BC`$. As a result, we have an inhomogeneous spherical body comprised of two homogeneous parts: one asymmetric convex lens (at right) and other asymmetric bispherical concavo-convex lens with surfaces having the radii of curvature $`R`$ (left surface) and $`R_1`$ (right surface). The asymmetric convex lens has density $`\rho _1+\rho _2`$ and central thickness $`h+h_1`$, the other asymmetric concavo-convex lens has density $`\rho _1`$ and central thickness $`h_2=DE=2Rhh_1`$.
This last figure is of a new kind of homogeneous figures and we are going to calculate the proper gravitational potential energy of this lens using MND provided the gravitational energy of asymmetric convex lens is known.
According to KA (see their Eq. (92)) the gravitational energy of homogeneous asymmetric lens is:
$$\begin{array}{c}d=R_2+R_1h_2h_1;a=\sqrt{2R_1h_1h_1^2}=\sqrt{2R_2h_2h_2^2;}\hfill \\ W_{ASL}(\rho _2,a,R_1,R_2)=\frac{\pi G\rho _2^2}{9}\{\pi [a^2(h_2+h_1)+R_2h_2^2+R_1h_1^2][R_2^2+R_1^2+R_2h_2+R_1h_1\frac{h_2^2+h_1^2}{2}]+\hfill \\ \pi [h_1(R_1h_1)(a^2+R_1h_1)h_2(R_2h_2)(a^2+R_2h_2)](h_2h_1+R_1R_2)\hfill \\ \pi [h_2^3(2R_2^2R_2h_2+\frac{h_2^2}{5})+h_1^3(2R_1^2R_1h_1+\frac{h_1^2}{5})]3\pi (R_2^3h_2^2+R_1^3h_1^2)\hfill \\ \frac{\pi (R_2^3+R_1^3)}{d}[h_2^2(3R_2h_2)+h_1^2(3R_1h_1)]\frac{32}{3}a^58a^3(R_2^2+R_1^2)\hfill \\ \frac{8a}{d}[R_2^2(R_2h_2)(R_2^2+\frac{2}{3}a^2)+R_1^2(R_1h_1)(R_1^2+\frac{2}{3}a^2)]+24a(R_2^4+R_1^4)+\hfill \\ 16R_2^4[\frac{R_2^2}{d}3(R_2h_2)]\mathrm{arctan}\frac{h_2}{a}+16R_1^4[\frac{R_1^2}{d}3(R_1h_1)]\mathrm{arctan}\frac{h_1}{a}\};\hfill \end{array}$$
(8)
here $`d`$ is distance between two spheres ”generating” asymmetric lens. Potential energy of ASL in the gravitational field of homogeneous sphere is (we assume here that $`\rho _2=\rho _1`$):
$$\begin{array}{c}W_{intASL}=\frac{G\pi ^2\rho _1^2}{90d}(R_1+R_2d)^2\hfill \\ (d^4+2d^3(R_1+R_2)+5(R_1R_2)^2(7R_1^2R_2^2)\hfill \\ 6d^2(2R_1^2R_1R_2+2R_2^2)2d(13R_1^3+9R_1^2R_2+9R_1R_2^27R_2^3)).\hfill \end{array}$$
(9)
Not repeating the procedure of previous section we present here only the formula for potential energy of concavo-convex lens in the next compact form:
$$W_{CCL}=W_{sphere}+W_{intASL}W_{ASL}(\rho _1,a,R_1,R_2);$$
(10)
see Eqs. (2), (9), (8).
### Homogeneous sphere and concavo-convex lens
In this section we calculate the potential energy of asymmetric convex lens provided the potential energy of concavo-convex lens is known. The setting of the problem is shown in Fig. 3. Here we attach to homogeneous sphere (at left) of radius $`R_1`$ and density $`\rho _1`$, the concavo-convex lens with radii of surface $`R_1`$ and $`R_2`$ and density $`\rho _2`$. Now we should calculate potential energy of CCL in the (external) field of sphere. We have ($`d`$ is the distance $`O_1O_2`$ between centers of spheres, $`R_1R_2dR_1+R_2`$:
$$W_{intCCL}=\frac{4G\pi ^2\rho _1\rho _2R_{1}^{}{}_{}{}^{3}\left(dR_12R_2\right)\left(dR_1+R_2\right)^2}{9d}.$$
(11)
Taking $`\rho _2=\rho _1`$ we get the formula for potential energy of homogeneous figure in Fig. 3 as a sum of three terms (see Eqs. (2), (11), (10)):
$$W_{sphere}+W_{intCCL}+W_{CCL},$$
(12)
which gives the potential energy of ASL coinciding with $`W_{ASL}(\rho _1,a,R_1,R_2)`$, see Eq. (8).
### Conclusion
In conclusion, we presented here the method of negative density which is a powerful tool to check the complicated formulas by KA for potential energy of homogeneous bispheric lenses. The method allows to obtain new formulas for the potential energy of homogeneous bodies as well. Though here we used MND only to study homogeneous bispherical lenses, the MND allows to obtain many new formulas, e.g. potential energies of spheres (and ellipsoids) and rectangles with spheres, ellipsoids and rectangles laying within (or adjacent to) them. The density $`\rho _2`$ of ”additional” bodies may be different from density $`\rho _1`$ of the ”parent” bodies or, according to MND, may be equal to $`\rho _1`$. We leave the study of these so called ”simple” bodies for next communications.
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# Mean parameter model for the Pekar-Fröhlich polaron in a multilayered heterostructure
## I Introduction
Quasi-two-dimensional (2D) systems have attracted a lot of attention during the last decade because of their practical realization. If a heterostructure is made of polar materials such as $`GaAs`$, $`InAs`$ etc., the electron-phonon interaction modifies the properties of the electron confined to a 2D-structure resulting in a shift of the binding energy and the effective band mass.
The polaron effects in a 2D electron gas have extensively been studied. At earlier stages the attention was paid to the properties of a polaron confined to an infinite thin 2D layer. The binding energy and the effective mass were calculated for a $`GaAs/Al_xGa_{1x}As`$ infinitely deep quantum well of a finite length. In these papers only the interaction with the bulk LO-phonon mode has been taken into account. Actually, LO-phonon modes are modified in a 2D layer (the so-called confined slab LO-phonon modes). Besides, there exist interface optical-phonon modes as well as half-space LO-phonon modes in a barrier material. For the review of these modes (also in complicated multi-layer structures) see the book by Pokatilov, Fomin and Beril and also more recent publications of this group. The influence of the mentioned modes on polarons were studied in Refs. .
While different phonon modes were studied in details, the quantum well potential was supposed to be infinitely deep in the cited papers. On the other hand, the properties of the system would be quite different if a confining potential had a finite depth. Indeed, for an infinitely deep confining potential the binding energy is the monotonous function of the potential width which varies between limiting values $`E_{3D}^{(in)}=\alpha _{in}\mathrm{}\omega _{in}`$ for the three-dimensional (3D) space and $`E_{2D}^{(in)}=(\pi /2)\alpha _{in}\mathrm{}\omega _{in}`$, where $`\alpha _{in}`$ is the standard Fröhlich electron-phonon coupling constant and $`\omega _{in}`$ is the LO-phonons frequency for the quantum well material. If a particle is confined to a finite potential well, the limiting value of the binding energy should be the same at large width of the well. But when the width becomes too small, the energy level approaches the edge of the well, so that effectively the particle is spread over the 3D space. Thus, the limiting value of the binding energy should coincide with the 3D limiting value rather than with the 2D one. This means, the binding energy takes the value $`E_{3D}^{(out)}=\alpha _{out}\mathrm{}\omega _{out}`$ at small widths where the parameters $`\omega _{out}`$ and $`\alpha _{out}`$ are now related to the barrier material. The binding energy evidently has a peak at some intermediate value of the width if $`E_{3D}^{(out)}E_{3D}^{(in)}`$. If this is not the case, the existence of the peak should be checked in more detail.
Different rectangular quantum wells of a finite height have been investigated by Hai, Peeters and Devreese and Shi, Zhu et al. in the scope of the second order perturbation theory in powers of the electron-phonon coupling constant $`\alpha `$ with all phonon modes being incorporated. Peaks of the phonon induced energy shift and the polaron effective mass were found for some values of the confining potential widths.
In principle, the same approach can be used while dealing with a quantum well constructed of layers of different materials. But the problem becomes then too complicated because one has to take into account interface phonon modes at each frontier of different materials as well as quantized phonon modes inside each of the layers. The main goal of the present paper is to formulate a simplified model to take these effects into account and to deal with the effective confining potential and only one bulk phonon mode. We calculate polaron characteristics for the same rectangular quantum well as in Refs. to compare the results. Another example is given of a quantum well of the finite depth for which the second-order correction due to the electron-phonon interaction can be calculated explicitly. Namely, we take the Rosen-Morse potential to confine electrons to a 2D-multilayered heterostructure and calculate the shift of the ground-state energy and the effective mass perturbatively, that is, in the weak-coupling limit. In contrast with the rectangular potential we should not worry about the correct including of all virtual states because the Green function is known analytically for this system.
## II Formulation of the model
Let us consider a quantum well in the $`z`$ direction constructed of the $`xy`$ plane layers of $`GaAs/Al_xGa_{1x}As`$. That is, the $`AlAs`$ mole fraction $`x`$ depends on the coordinate $`z`$: $`x=x(z)`$. The energy gap between different materials forms the confining potential $`V(z)`$ which serves us as the main entity. Given the potential $`V(z)`$, one can find the corresponding mole fractions $`x(z)`$ and a dependence on $`z`$ of any of the medium parameters (such as the electron band mass $`m(z)`$, phonon frequencies $`\omega (z)`$, dielectric constants $`\epsilon _0(z),\epsilon _{\mathrm{}}(z)`$, Fröhlich coupling constants $`\alpha (z)`$, etc.).
To avoid difficulties with mass mismatch in different layers we suggest to use a mean band mass $`m`$ which is common for all layers. Then we start with the following expression for the electronic part of the Hamiltonian:
$`H_{el}`$ $`=`$ $`H_{el,}+H_{el,},`$ (1)
$`H_{el,}`$ $`=`$ $`{\displaystyle \frac{\stackrel{}{p}_{}^{\mathrm{\hspace{0.17em}2}}}{2m}},H_{el,}={\displaystyle \frac{p_z^{\mathrm{\hspace{0.17em}2}}}{2m}}+V(z),`$ (2)
where the electron mean band mass is defined by the relation
$`{\displaystyle \frac{1}{m}}={\displaystyle 𝑑z\frac{|\psi _1(z)|^2}{m(z)}}`$ (3)
and the ground state wave function $`\psi _1(z)`$ for the electron motion in $`z`$ direction is a solution to the Schrödinger equation
$`H_{el,}\psi _1=E_1\psi _1`$ (4)
with $`E_1`$ being a ground state energy. As the wave function $`\psi _1`$ also depends on the mean band mass $`m`$, the latter can be found as a self-consistent solution of Eqs. (3), (4).
In a similar way we define the free LO-phonon Hamiltonian
$`H_{ph}=\mathrm{}\omega _{\text{LO}}{\displaystyle \underset{\stackrel{}{k}}{}}a_\stackrel{}{k}^{}a_\stackrel{}{k},`$ (5)
where $`a_\stackrel{}{k}^{}(a_\stackrel{}{k})`$ are creation (annihilation) operators of a phonon with a wave vector $`\stackrel{}{k}`$, and mean frequency $`\omega _{\text{LO}}`$ can be found from the expression
$`\omega _{\text{LO}}={\displaystyle 𝑑z\omega (z)|\psi _1(z)|^2}.`$ (6)
Evidently, we have to address why the free phonon Hamiltonian is averaged with respect to the electron wave function. Our motivation is based on the fact that we are going to apply our model to calculate polaron effects. That is, our effective phonons will be considered only in a cloud around the electron, and the properties of this cloud depend on the electron position. So, in our model the effective phonons replace numerous phonon modes whose frequencies depend on the coordinate $`z`$ of the electron.
Finally, we accept the conventional form of the Hamiltonian describing the interaction of the electron with effective phonons:
$`H_{int}={\displaystyle \underset{\stackrel{}{k}}{}}\left(a_\stackrel{}{k}V_\stackrel{}{k}e^{i\stackrel{}{k}\stackrel{}{r}}+a_\stackrel{}{k}^{}V_\stackrel{}{k}^{}e^{i\stackrel{}{k}\stackrel{}{r}}\right),`$ (7)
where the Fourier transforms of the electron-phonon interaction potential are specified as follows:
$`V_\stackrel{}{k}=i\mathrm{}\omega _{\text{LO}}\left({\displaystyle \frac{4\pi \alpha }{Vk^2}}\sqrt{{\displaystyle \frac{\mathrm{}}{2m\omega _{\text{LO}}}}}\right)^{1/2}.`$ (8)
Here the use is made of a mean Fröhlich coupling constant $`\alpha `$ which can be found from the relation
$`\sqrt{\alpha }={\displaystyle 𝑑z|\psi _1(z)|^2\frac{\omega (z)}{\omega _{\text{LO}}}\left(\alpha (z)\sqrt{\frac{m\omega _{\text{LO}}}{m(z)\omega (z)}}\right)^{1/2}}.`$ (9)
Note that we define the mean parameters in Eqs. (3), (6), (9) according to the way they enter the Hamiltonian.
Thus, we describe a complicated multilayered heterostructure by the Hamiltonian
$`H=H_{el}+H_{ph}+H_{int}`$ (10)
with the bulk phonon mode only which inhabits an effective medium with mean characteristics defined above. The details of the heterostructure are taken into account by the confining potential $`V(z)`$.
Performing a unitary transformation $`HH^{}=U^1HU`$ with the operator
$`U=\mathrm{exp}\left[i\stackrel{}{r}_{}{\displaystyle \underset{\stackrel{}{k}}{}}\stackrel{}{k}_{}a_\stackrel{}{k}^{}a_\stackrel{}{k}\right],`$ (11)
we arrive at the Hamiltonian
$`H^{}=H_{el,}^{}+H_{el,}+H_{ph}+H_{int}^{},`$ (13)
$`H_{el,}^{}`$ $`=`$ $`{\displaystyle \frac{1}{2m}}\left(\stackrel{}{p}_{}\mathrm{}{\displaystyle \underset{\stackrel{}{k}}{}}\stackrel{}{k}_{}a_\stackrel{}{k}^{}a_\stackrel{}{k}\right)^2,`$ (14)
$`H_{int}^{}={\displaystyle \underset{\stackrel{}{k}}{}}\left(a_\stackrel{}{k}V_\stackrel{}{k}e^{ik_zz}+a_\stackrel{}{k}^{}V_\stackrel{}{k}^{}e^{ik_zz}\right).`$ (15)
The quantity $`\stackrel{}{p}_{}`$ is a c-number corresponding to the conserved momentum in the $`xy`$ plane and the Hamiltonians $`H_{el,},H_{ph}`$ are defined by Eqs. (2), (5), respectively.
Keeping in mind the smallness of the electron-phonon coupling constant $`\alpha `$ for most of the materials, we calculate the second-order correction to the unperturbed Hamiltonian $`H_0^{}=H_{el,}^{}+H_{el,}+H_{ph}`$ (note that the quantum-mechanical first-order correction is equal to zero). The unperturbed energy levels are given by the expression
$`E(\stackrel{}{p}_{},n_\stackrel{}{k},N)`$ $`=`$ $`{\displaystyle \frac{1}{2m}}\left(\stackrel{}{p}_{}\mathrm{}{\displaystyle \underset{\stackrel{}{k}}{}}\stackrel{}{k}_{}n_\stackrel{}{k}\right)^2+`$ (17)
$`\mathrm{}\omega _{\text{LO}}{\displaystyle \underset{\stackrel{}{k}}{}}n_\stackrel{}{k}+E_N,`$
where $`n_\stackrel{}{k}`$ is the number of phonons with the wave vector $`\stackrel{}{k}`$. The energy $`E_N`$ is the $`N`$-th energy level of the one-dimensional system $`H_{el,}`$ of Eq. (2). Here $`N`$ is the corresponding quantum number not necessarily a discrete one: it stands for both the quantum number $`n`$ which varies from 1 to $`n_{max}`$ and the wave vector $`q`$ of the continuous spectrum states.
The wave functions of the unperturbed Hamiltonian $`H_0^{}`$ are given by the direct product
$`|\stackrel{}{p}_{};n_\stackrel{}{k},N=|n_\stackrel{}{k}|N`$ (18)
of the corresponding wave functions of different terms in $`H_0^{}`$.
Because of the structure of the interaction term $`H_{int}^{}`$ only intermediate states with one phonon contribute to the second order correction to the ground-state energy. The latter is then given by the expression
$`\mathrm{\Delta }_2E(\stackrel{}{p}_{})=`$ (19)
$`{\displaystyle \underset{N,\stackrel{}{k}}{}}{\displaystyle \frac{|V_\stackrel{}{k}|^2|G(N,k_z)|^2}{E_N+\mathrm{}\omega _{\text{LO}}+{\displaystyle \frac{(\stackrel{}{p}_{}\mathrm{}\stackrel{}{k}_{})^2\stackrel{}{p}_{}^{}{}_{}{}^{\mathrm{\hspace{0.17em}2}}}{2m}}E_1}},`$ (20)
where
$`G(N,k_z)={\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}𝑑z\psi _N(z)\psi _1(z)e^{ik_zz}`$ (21)
and $`\psi _N(z)`$ are the wave functions of the Hamiltonian $`H_{el,}`$ in Eq. (2). The concrete application of these formulae is given in the following section.
## III Rectangular Potential
### A Medium mean characteristics
As an example we now consider the rectangular confining potential
$`V(z)=\{\begin{array}{cc}0,\hfill & |z|L/2\hfill \\ V_0,\hfill & |z|>L/2\hfill \end{array}`$ (24)
and
$`m(z)=\{\begin{array}{cc}m_{in},\hfill & |z|L/2\hfill \\ m_{out},\hfill & |z|>L/2\hfill \end{array}`$ (27)
with $`m_{in}(m_{out})`$ being the electron band masses in the well (barrier) material, respectively. For concreteness we assume $`GaAs`$ to be the quantum well material and $`Al_xGa_{1x}As`$ to be the barrier material.
Symmetrical wave functions of the discrete spectrum in the rectangular quantum well with the mean band mass $`m`$ take the form
$`\psi _{s,n}=N_{s,n}\{\begin{array}{cc}\mathrm{cos}q_nz,\hfill & |z|L/2\hfill \\ \mathrm{cos}{\displaystyle \frac{q_nL}{2}}e^{p_n(|z|L/2)},\hfill & |z|>L/2,\hfill \end{array}`$ (30)
where
$`p_n=\sqrt{q_{max}^2q_n^2},q_{max}^2={\displaystyle \frac{2mV_0}{\mathrm{}^2}}`$ (31)
and the normalization constant is given by
$`N_{s,n}=\sqrt{{\displaystyle \frac{2p_n}{p_nL+2}}}.`$ (32)
Antisymmetrical wave functions of the discrete spectrum take the form
$`\psi _{a,n}=N_{s,n}\{\begin{array}{cc}\mathrm{sin}q_nz,\hfill & |z|L/2\hfill \\ \text{sgn}(z)\mathrm{sin}{\displaystyle \frac{q_nL}{2}}e^{p_n(|z|L/2)},\hfill & |z|>L/2\hfill \end{array}`$ (35)
(36)
with the same normalization constant given by Eq. (32).
The total number $`n_{max}`$ of the discrete energy levels is given by the expression
$`n_{max}=1+\left[{\displaystyle \frac{q_{max}L}{\pi }}\right],`$ (37)
where $`[A]`$ is an integer part of a number $`A`$. The expression for the discrete energy levels reads as follows:
$`{\displaystyle \frac{q_nL}{2}}=\text{atan}\sqrt{{\displaystyle \frac{q_{max}^2}{q_n^2}}1}+{\displaystyle \frac{\pi (n1)}{2}},n=1,2,\mathrm{}.`$ (38)
Energies with odd (even) $`n`$ correspond to the symmetrical (antisymmetrical) wave functions.
The energy $`E_q=\mathrm{}^2q^2/2m`$ of the continuous spectrum state depends on the wave vector $`q`$. The corresponding symmetrical wave functions are as follows:
$`\psi _{s,q}={\displaystyle \frac{N_{s,q}}{\sqrt{L_z}}}\{\begin{array}{cc}p\mathrm{cos}qz,\hfill & |z|L/2,\hfill \\ p\mathrm{cos}{\displaystyle \frac{qL}{2}}\mathrm{cos}p(|z|L/2)\hfill & \\ q\mathrm{sin}{\displaystyle \frac{qL}{2}}\mathrm{sin}p(|z|L/2),\hfill & |z|>L/2,\hfill \end{array}`$ (42)
(43)
where
$`p=\sqrt{q^2q_{max}^2}`$ (44)
and $`L_z`$ is the (infinite) size of the system in the $`z`$ direction. The normalization constant is given by the expression
$`N_{s,q}=\sqrt{{\displaystyle \frac{2}{p^2\mathrm{cos}^2{\displaystyle \frac{qL}{2}}+q^2\mathrm{sin}^2{\displaystyle \frac{qL}{2}}}}}.`$ (45)
The antisymmetrical wave functions are as follows:
$`\psi _{a,q}={\displaystyle \frac{N_{a,q}\text{sgn}(z)}{\sqrt{L_z}}}\{\begin{array}{cc}p\mathrm{sin}q|z|,|z|L/2,\hfill & \\ p\mathrm{sin}{\displaystyle \frac{qL}{2}}\mathrm{cos}p(|z|L/2)+\hfill & \\ q\mathrm{cos}{\displaystyle \frac{qL}{2}}\mathrm{sin}p(|z|L/2),\hfill & \\ |z|>L/2,\hfill & \end{array}`$ (50)
(51)
where the normalization constant is given by the expression
$`N_{a,q}=\sqrt{{\displaystyle \frac{2}{p^2\mathrm{sin}^2{\displaystyle \frac{qL}{2}}+q^2\mathrm{cos}^2{\displaystyle \frac{qL}{2}}}}}.`$ (52)
The electron mean band mass is defined as
$`{\displaystyle \frac{1}{m}}`$ $`=`$ $`{\displaystyle \frac{W_{in}}{m_{in}}}+{\displaystyle \frac{W_{out}}{m_{out}}}`$ (53)
$`m`$ $`=`$ $`{\displaystyle \frac{m_{in}m_{out}}{W_{in}m_{out}+(1W_{in})m_{in}}},`$ (54)
where $`W_{in}`$ and $`W_{out}=1W_{in}`$ are probabilities to find the electron inside (outside) the quantum well. The expression for $`W_{in}`$ follows from Eq. (30)
$`W_{in}=2N_{s,1}^2{\displaystyle \underset{0}{\overset{L/2}{}}}𝑑z\mathrm{cos}^2q_1z=1{\displaystyle \frac{(q_1/q_{max})^2}{1+p_1L/2}},`$ (55)
where $`q_1`$ is a solution to Eq. (38) for the ground-state ($`n=1`$).
To finish this subsection, we note that the exact energy levels in the rectangular potential with different masses $`m_{in}`$ and $`m_{out}`$ calculated for the $`GaAs/Al_xGa_{1x}As`$ heterostructure practically coincide with the levels obtained with the electron mean band mass $`m`$. To obtain an inner criterion of the validity of the anzatz concerning the mean band mass we notice that the particle being on lowest energy levels is located mostly inside the well which means that its band mass is almost coincide with $`m_{in}`$.
One can await the largest discrepancy for a level near the potential edge. The $`n`$-th discrete level appears at the width $`L=L_n^{(av)}`$, where
$`L_n^{(av)}=\pi (n1){\displaystyle \frac{\mathrm{}}{\sqrt{2mV_0}}}={\displaystyle \frac{\pi (n1)}{q_{\mathrm{max}}}},`$ (56)
and the analogous width for the exact solution reads as follows:
$`L_n^{(ex)}=\pi (n1){\displaystyle \frac{\mathrm{}}{\sqrt{2m_{in}V_0}}}.`$ (57)
Thus, the ratio
$`{\displaystyle \frac{L_n^{(av)}}{L_n^{(ex)}}}=\sqrt{{\displaystyle \frac{m_{in}}{m}}}=\sqrt{W_{in}+(1W_{in}){\displaystyle \frac{m_{in}}{m_{out}}}}.`$ (58)
can serve us as the numerical criterion of the validity of the anzatz. The largest discrepancy happens at $`n=2`$ and in this case Eqs. (31), (38) (55) lead to the following expression:
$`{\displaystyle \frac{L^{(av)}}{L^{(ex)}}}=\sqrt{0.844+0.156{\displaystyle \frac{m_{in}}{m_{out}}}}.`$ (59)
Note that numerical coefficients here do not depend on the material parameters. For the $`GaAs/Al_{0.3}Ga_{0.7}As`$ quantum well we have $`m_{in}/m_{out}0.7`$ and the discrepancy is about $`2\%`$; in the worst possible case when $`m_{in}/m_{out}1`$ the discrepancy is still not large: $`100\%\sqrt{0.844}8\%`$.
### B Electron-phonon correction to the polaron energy and the effective mass
Summation over the wave vector $`\stackrel{}{k}`$ in Eq. (20) can be reduced to integration in a conventional way
$`{\displaystyle \underset{\stackrel{}{k}}{}}|V_\stackrel{}{k}|^2(\mathrm{})={\displaystyle \frac{V}{(2\pi )^3}}{\displaystyle 𝑑\stackrel{}{k}|V_\stackrel{}{k}|^2(\mathrm{})}=`$ (60)
$`(\mathrm{}\omega _{\text{LO}})^2\sqrt{{\displaystyle \frac{\mathrm{}}{2m\omega _{\text{LO}}}}}{\displaystyle \frac{\alpha }{2\pi ^2}}{\displaystyle \frac{d\stackrel{}{k}_{}dk_z}{k_{}^2+k_z^2}(\mathrm{})}.`$ (61)
Then, the integration over $`\stackrel{}{k}_{}`$ in Eq. (20) can be performed explicitly. As we are interested in corrections to the ground-state energy and the effective mass $`m_{eff}m+\mathrm{\Delta }_2m`$ of the polaron motion in the $`xy`$ plane, we expand $`\mathrm{\Delta }_2E(\stackrel{}{p}_{})\mathrm{\Delta }_2E{\displaystyle \frac{\mathrm{\Delta }_2m}{2m^2}}\stackrel{}{p}_{}^2`$ in powers of the conserved momentum $`\stackrel{}{p}_{}`$. Doing this the use is made of the integral
$`{\displaystyle \frac{d^2\stackrel{}{k}_{}}{(\stackrel{}{k}_{}^2+k_z^2)[\stackrel{}{k}_{}^{}{}_{}{}^{\mathrm{\hspace{0.17em}2}}2\stackrel{}{k}_{}\stackrel{}{p}_{}/\mathrm{}+b^2]}}\pi {\displaystyle \frac{\mathrm{ln}(k_z^2/b^2)}{k_z^2b^2}}+`$ (62)
$`\left({\displaystyle \frac{\stackrel{}{p}_{}}{\mathrm{}}}\right)^2\pi {\displaystyle \frac{k_z^4b^42k_z^2b^2\mathrm{ln}(k_z^2/b^2)}{b^2(k_z^2b^2)^3}}.`$ (63)
As the next step we use dimensionless “polaronic” units performing the scaling $`k_zk_z\sqrt{2m\omega _{\text{LO}}/\mathrm{}},zz\sqrt{\mathrm{}/2m\omega _{\text{LO}}}`$ and using the notation
$`l=L\sqrt{{\displaystyle \frac{2m\omega _{\text{LO}}}{\mathrm{}}}},\epsilon _N={\displaystyle \frac{E_N}{\mathrm{}\omega _{\text{LO}}}}.`$ (64)
In these units the correction to the ground-state energy takes the form
$`{\displaystyle \frac{\mathrm{\Delta }_2E}{\mathrm{}\omega _{\text{LO}}}}=`$ (65)
$`{\displaystyle \frac{\alpha }{\pi }}{\displaystyle \underset{N}{}}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑k_z{\displaystyle \frac{\mathrm{ln}(k_z^2/b_N^2)}{k_z^2b_N^2}}\left(|G_s(N,k_z)|^2+|G_a(N,k_z)|^2\right),`$ (66)
(67)
where
$`b_N=\sqrt{\epsilon _N+1\epsilon _1}.`$ (68)
The correction to the effective mass reads as follows:
$`{\displaystyle \frac{\mathrm{\Delta }_2m}{m}}=`$ (69)
$`{\displaystyle \frac{\alpha }{\pi }}{\displaystyle \underset{N}{}}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}dk_z{\displaystyle \frac{k_z^4b_N^42k_z^2b_N^2\mathrm{ln}(k_z^2/b_N^2)}{b_N^2(k_z^2b_N^2)^3}}\times `$ (70)
$`\left(|G_s(N,k_z)|^2+|G_a(N,k_z)|^2\right).`$ (71)
Quantities $`G_j(N,k_z)`$ in Eqs. (67), (71) are given in dimensionless units by the same Eq. (21); the indices $`(a)s`$ are related to (anti)symmetrical wave functions being used in Eq. (21):
$`G_s(N,k_z)=2{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑z\psi _{s,N}(z)\psi _{s,1}(z)\mathrm{cos}k_zz,`$ (72)
$`G_a(N,k_z)=2{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑z\psi _{a,N}(z)\psi _{s,1}(z)\mathrm{sin}k_zz,`$ (73)
Evidently, the replacement $`Ll`$ should be done in the definition of the wave functions and their normalization constants; in addition $`L_zl_z`$ in Eqs. (43), (51) as well as in Eq. (38) for the energy levels of the discrete spectrum. Eq. (31) now reads as follows:
$`p_n=\sqrt{v_0q_n^2},v_0={\displaystyle \frac{V_0}{\mathrm{}\omega _{\text{LO}}}},q_{max}^2=v_0.`$ (74)
Eq. (37) takes the form
$`n_{max}=1+\left[{\displaystyle \frac{\sqrt{v_0}l}{\pi }}\right].`$ (75)
The relation of dimensionless energies of the discrete and continuous spectra with subsequent wave vectors takes the form $`\epsilon _n=q_n^2,\epsilon _q=q^2`$. All the changes mentioned should also be done in Eq. (55).
The final note of this section concerns summation over $`N`$ in Eqs. (67), (71):
$`{\displaystyle \underset{N}{}}(\mathrm{})={\displaystyle \underset{n=1}{\overset{n_{max}}{}}}(\mathrm{})+\underset{l_z\mathrm{}}{lim}{\displaystyle \frac{l_z}{2\pi }}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑p(\mathrm{}).`$ (76)
The replacing of the sum over the continuous spectrum by the integration over the wave vector $`p`$ follows from Eqs. (43, 51) in the limit $`L_z\mathrm{}`$. The wave vectors $`q`$ and $`p`$ are related to each other because of Eq. (44) which now takes the form $`q=\sqrt{p^2+v_0}`$. Note also that only $`G_s(N,k_z)`$ ($`G_a(N,k_z)`$) has to be taken into account for odd (even) $`n`$ in the sum over the discrete quantum number $`n`$.
The numerical results obtained are plotted in Fig. 1 for $`\mathrm{\Delta }_2E`$ and in Fig. 2a for $`\mathrm{\Delta }_2m/m`$. Because the mean effective mass $`m`$ depends on the potential width we also plotted in Fig. 2b the ratio of the mass shift to those in the well material, that is, the ratio
$`\delta _2m={\displaystyle \frac{\mathrm{\Delta }_2m}{\mathrm{\Delta }_2m_{in}}},\mathrm{\Delta }_2m_{in}=m_{in}{\displaystyle \frac{\alpha _{in}}{6}}.`$ (77)
The discussion of the numerical results is given in the last section.
## IV Rosen-Morse potential
### A Energy-dependent Green function
In this section we present another example — a multilayered heterostructure described by a confining potential $`V(z)`$ which is chosen in the form of the Rosen-Morse potential
$`V(z)`$ $`=`$ $`V_0\mathrm{tanh}^2\left({\displaystyle \frac{z}{L_{RM}}}\right),`$ (78)
$`V_0`$ $`=`$ $`{\displaystyle \frac{\mathrm{}^2}{2mL_{RM}^2}}\kappa (\kappa +1).`$ (79)
where $`L_{RM}`$ is the parameter close to the half-width of the Rosen-Morse quantum well and $`\kappa `$ is the dimensionless parameter to govern the strength of the potential.
The summation (20) over the quantum number $`N`$ can be represented through the Green function which is known analytically for the Rosen-Morse potential. Namely, the second-order correction to the ground-state energy can be written in the form
$`\mathrm{\Delta }_2E`$ $`=`$ $`\mathrm{}\omega _{\text{LO}}\alpha {\displaystyle \frac{l_{RM}}{\sqrt{2}}}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑k_{}{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}𝑑z_a{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}𝑑z_be^{k_{}|z_az_b|}`$ (81)
$`\psi _1^{}(z_a)\psi _1(z_b)G(z_a,z_b;E),`$
where we made a scaling $`zzL_{RM},\stackrel{}{k}\stackrel{}{k}/L_{RM}`$ to use dimensionless variables $`z,\stackrel{}{k}`$ and integrated over $`k_z`$ and angles of $`\stackrel{}{k}_{}`$. The dimensionless parameter
$`l_{RM}=L_{RM}\sqrt{{\displaystyle \frac{m\omega _{\text{LO}}}{\mathrm{}}}}`$ (82)
is the width of the confining potential in polaronic units while the potential strength can now be written as follows:
$`V_0=\mathrm{}\omega _{\text{LO}}{\displaystyle \frac{\kappa (\kappa +1)}{2l_{RM}^2}}.`$ (83)
The quantity $`G(z_a,z_b;E)`$ is the Green function of the dimensionless Hamiltonian (2) which takes the form
$`H_{el,}^{\prime \prime }={\displaystyle \frac{1}{2}}{\displaystyle \frac{d^2}{dz^2}}+{\displaystyle \frac{\kappa (\kappa +1)}{2}}\mathrm{tanh}^2z,`$ (84)
that is $`G(z_a,z_b;E)=z_a|(H_{el,}^{\prime \prime }E)^1|z_b`$, while $`\psi _1(z)`$ is the ground-state wave function of the potential (84)
$`\psi _1(z)=\left[{\displaystyle \frac{\mathrm{\Gamma }(\kappa +1/2)}{\sqrt{\pi }\mathrm{\Gamma }(\kappa )}}\right]^{1/2}{\displaystyle \frac{1}{\mathrm{cosh}^\kappa z}}.`$ (85)
The ground-state energy of the Hamiltonian (84) is given by
$`E_1={\displaystyle \frac{\kappa }{2}}.`$ (86)
The energy $`E`$ in Eq. (81) reads as follows
$`E={\displaystyle \frac{k_{}^2}{2}}l_{RM}^2+{\displaystyle \frac{\kappa }{2}}.`$ (87)
The energy-dependent Green function of the system can be represented in the form:
$`G(z_a,z_b;E)=`$ (88)
$`{\displaystyle \frac{\mathrm{\Gamma }(\nu +\kappa +1)\mathrm{\Gamma }(\nu \kappa )}{\mathrm{\Gamma }^2(\nu +1)}}{\displaystyle \frac{1}{(4\mathrm{cosh}z_a\mathrm{cosh}z_b)^\nu }}\times `$ (89)
$`{}_{2}{}^{}F_{1}^{}(\nu \kappa ,\nu +\kappa +1;\nu +1;{\displaystyle \frac{1\mathrm{tanh}z_>}{2}})\times `$ (90)
$`{}_{2}{}^{}F_{1}^{}(\nu \kappa ,\nu +\kappa +1;\nu +1;{\displaystyle \frac{1+\mathrm{tanh}z_<}{2}}),`$ (91)
where $`z_>(z_<)`$ denotes the maximum (minimum) of $`z_a`$ and $`z_b`$. The parameter $`\nu `$ is defined by the relation
$`\nu =\sqrt{2\left(E{\displaystyle \frac{\kappa (\kappa +1)}{2}}\right)}=\sqrt{k_{}^2+\kappa ^2+2l_{RM}^2}.`$ (92)
The polaron effective mass can be represented in a similar way
$`{\displaystyle \frac{\mathrm{\Delta }_2m}{m}}`$ $`=`$ $`\alpha {\displaystyle \frac{l^3}{2\sqrt{2}}}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑k_{}k_{}^2{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}𝑑z_a{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}𝑑z_be^{k_{}|z_az_b|}`$ (94)
$`\psi _1^{}(z_a)\psi _1(z_b){\displaystyle \frac{^2}{E^2}}G(z_a,z_b;E).`$
To simplify numerical calculations we may replace the derivative with respect to $`E`$ by the derivative with respect to $`\nu `$
$`{\displaystyle \frac{^2}{E^2}}={\displaystyle \frac{1}{\nu ^2}}{\displaystyle \frac{^2}{\nu ^2}}{\displaystyle \frac{1}{\nu ^3}}{\displaystyle \frac{}{\nu }}`$ (95)
and perform once the integration by parts. As the result, we arrive at the following representation equivalent to Eq. (94):
$`{\displaystyle \frac{\mathrm{\Delta }_2m}{m}}=\alpha {\displaystyle \frac{l^3}{2\sqrt{2}}}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑k_{}{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}𝑑z_a{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}𝑑z_b(1k_{}|z_az_b|)`$ (96)
$`e^{k_{}|z_az_b|}\psi _1^{}(z_a)\psi _1(z_b)\left[{\displaystyle \frac{1}{\nu }}{\displaystyle \frac{}{\nu }}\right]G(z_a,z_b;E).`$ (97)
Note that $`m,\alpha ,\omega _{\text{LO}}`$ in all these formulae stand for the mean characteristics of the medium. The wave function in their definitions is given by Eq. (85). The numerical results are plotted in Fig. 3 and discussed in the last Section.
### B Effective width
If we decide to compare the results for the rectangular and Rosen-Morse potentials, we have to define a parameter which plays the role of the effective width of the Rosen-Morse potential. That is, this parameter (for which we use a notation $`L`$) should be close to $`2L_{RM}`$ of Eq. (79) being also related to the rectangular potential. We accept the following definition: let us call the effective width of the Rosen-Morse potential the width $`L`$ of the rectangular well of the same height with the same ground-state energy in the absence of the electron-phonon interaction (that is, at $`\alpha =0`$). The advantage of this definition is that while calculating the polaron binding energy for the Rosen-Morse and rectangular potentials we subtract the same quantity in both the cases and can compare only energy shifts due to the electron-phonon interaction.
The ground-state energy of a rectangular potential with the height $`V_0`$ and width $`L`$ is given by the relations
$`E_{RC}`$ $`=`$ $`{\displaystyle \frac{\mathrm{}^2k^2}{2m}},`$ (98)
$`\text{tan}{\displaystyle \frac{kL}{2}}`$ $`=`$ $`\sqrt{{\displaystyle \frac{V_0}{E_{RC}}}1},`$ (99)
while the RM ground-state energy looks like
$`E_{RM}={\displaystyle \frac{\mathrm{}^2}{mL_{RM}^2}}{\displaystyle \frac{\kappa }{2}}`$ (100)
and the height $`V_0`$ of the potential is given by Eq. (79). With the equality $`E_{RM}=E_{RC}`$ we arrive at the relation between the parameter $`L_{RM}`$ of the Rosen-Morse potential and its effective width defined as has been discussed:
$`{\displaystyle \frac{L}{L_0}}`$ $`=`$ $`2\sqrt{\lambda }\text{arctg}\sqrt{\lambda 1},`$ (101)
$`\lambda `$ $`=`$ $`\kappa +1={\displaystyle \frac{1}{2}}\left[1+\sqrt{1+\left(2L_{RM}/L_0\right)^2}\right].`$ (102)
Here we introduce the distance scale
$`L_0=\sqrt{{\displaystyle \frac{\mathrm{}^2}{2mV_0}}}.`$ (103)
The relation to the other dimensionless parameter $`l_{RM}`$ of Eq. (82) is given by
$`{\displaystyle \frac{L_{RM}}{L_0}}=l\sqrt{{\displaystyle \frac{2V_0}{\mathrm{}\omega _{\text{LO}}}}}.`$ (104)
At small $`L_{RM}L_0`$ we obtain $`L2L_{RM}`$ from Eq. (102), that is indeed the parameter $`L_{RM}`$ plays a role of the half-width of the Rosen-Morse potential in this case. When $`L_{RM}L_0`$, it follows from Eq. (102) that $`L\pi \sqrt{L_{RM}L_0}`$.
The effective width $`L`$ defined in this subsection allows us to apply the results for the rectangular potential to the Rosen-Morse quantum well. The example is given in Fig. 3 where we plotted also the energy and the mass shifts for the rectangular potential vs. the parameter $`L_{RM}`$ related to $`L`$ as is described.
## V Numerical results and discussion
To proceed to the numerical calculations we need now the dependence of medium parameters on the $`AlAs`$ mole fraction $`x`$. At first we present the parametrization from the review by Adachi:
$`\alpha (z)=0.068+0.058x,`$ (106)
$`m(z)=m_e(0.0665+0.0835x),`$ (107)
$`\mathrm{}\omega (z)=(36.25+1.83x+17.12x^25.11x^3)\mathrm{meV},`$ (108)
which was used in numerical calculations by Hai, Peeters and Devreese. Here $`m_e`$ is the electron mass in vacuum and $`m(z)`$ — its band mass in the subsequent layer; the values of the electron-phonon coupling constant $`\alpha (z)`$ and the LO-phonon frequency $`\omega (z)`$ are also related to this layer.
Some comments are to the point. The expression for the electron band mass is nothing else but the linear interpolation between the values $`m=0.0665m_e`$ for $`GaAs`$ and $`m=0.150m_e`$ for $`AlAs`$. As to the LO-phonon frequency there are two phonon modes with different frequencies $`\omega ^{(G)}(z)`$ and $`\omega ^{(A)}(z)`$ for the $`GaAs`$-like and $`AlAs`$-like modes in $`Al_xGa_{1x}As`$ crystal. Experimental results of Ref. are interpolated by the following formulae:
$`\mathrm{}\omega ^{(G)}(z)=(36.256.55x+1.79x^2)\mathrm{meV},`$ (110)
$`\mathrm{}\omega ^{(A)}(z)`$ $`=`$ $`(44.63+8.78x3.32x^2)\mathrm{meV}.`$ (111)
Because the exact theory of the two-phonon interaction in alloys where there are two-mode phonons present has not been reported, Adachi suggested to use the effective phonon frequency $`\omega =(1x)\omega ^{(G)}+x\omega ^{(A)}`$, that is the linear interpolation between these two modes. Inserting here the expressions (V) one arrives at the result (108).
As to the interpolation formula (106) for the Fröhlich coupling constant $`\alpha `$, the situation seems to be a bit inconsistent. Indeed, $`\alpha `$ depends on the values of the static $`\epsilon _0`$ and the high-frequency $`\epsilon _{\mathrm{}}`$ dielectric constants:
$`\alpha `$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{}\omega }}{\displaystyle \frac{\overline{e}^2}{\sqrt{2}}}\sqrt{{\displaystyle \frac{m\omega }{\mathrm{}}}}\left({\displaystyle \frac{1}{\epsilon _{\mathrm{}}}}{\displaystyle \frac{1}{\epsilon _0}}\right)`$ (112)
$`=`$ $`116.643\left({\displaystyle \frac{1}{\epsilon _{\mathrm{}}}}{\displaystyle \frac{1}{\epsilon _0}}\right)\sqrt{{\displaystyle \frac{m}{m_e}}}\sqrt{{\displaystyle \frac{1\mathrm{meV}}{\mathrm{}\omega }}}.`$ (113)
Earlier measurements of $`\epsilon _0`$ of $`GaAs`$ have yielded widely different values ranging from 9.8 to 13.3 (see Ref. and references therein). For instance, Kartheuser reports the values $`\epsilon _{\mathrm{}}=10.9`$ and $`\epsilon _0=12.83`$ and $`\mathrm{}\omega =36.75\mathrm{meV}`$ for $`GaAs`$. This leads to the result $`\alpha =0.068`$, which is widely known and used by many people.
On the other hand, Adachi used the more recent results for $`GaAs`$: $`\epsilon _0=13.18\pm 0.40`$ and $`\epsilon _{\mathrm{}}=10.89`$, and for $`AlAs`$: $`\epsilon _0=10.06\pm 0.04`$ and $`\epsilon _{\mathrm{}}=8.16\pm 0.02`$. This gives birth to his interpolation formulae:
$`\epsilon _0=13.183.12x,`$ (115)
$`\epsilon _{\mathrm{}}=10.892.73x,`$ (116)
Inserting formulae (107), (108) and (V) into Eq. (113) Adachi declared the result $`\alpha =0.126`$ for $`AlAs`$. Together with the value $`\alpha =0.068`$ reported in Ref. this leads to the interpolation formulae (106). The problem is that both these values for $`\alpha `$ do not follow from the parametrizations mentioned above.
Taking the same values for $`AlAs`$ as Adachi did take ($`m=0.150m_e,\mathrm{}\omega =50.09\mathrm{meV},\epsilon _0=8.16,\epsilon _{\mathrm{}}=10.06`$) we arrive at the result $`\alpha =0.1477`$. Moreover, if one takes the same interpolation formulae (V) at $`x=0`$ one obtains the value $`\alpha =0.0797`$ for GaAs. That is, Adachi had to obtain the formula
$`\alpha (z)=0.0797+0.0680x`$ (117)
as a linear interpolation between the values of $`\alpha `$ in $`GaAs`$ and $`AlAs`$. Note, that this formulae can be presented in the form $`\alpha (z)=1.172(0.068+0.058x)`$. The expression between the brackets coincide (probably occasionally) with the Adachi interpolation formulae for $`\alpha `$ (cf. Eq. (106)). That is, the discrepancy of (106) and of our interpolation (117) is about 17% and do not depend on $`x`$. To be consistent we have to accept the parametrization (117) in what follows.
For the confining potential we take the expression derived from the band-gap energy fit of Ref. and used in Ref. :
$`V(z)=600(1.155x+0.37x^2)\mathrm{meV}.`$ (118)
Thus, we use the parametrization (107), (108), (115), (116), (117) and the potential (118) in our numerical calculations.
The results of our study for the rectangular potential (which is formed by a layer of $`GaAs/Al_xGa_{1x}As`$) are shown in Fig. 1 for the polaronic energy shift and in Fig. 2 for the polaron effective mass at the $`AlAs`$ mole fraction $`x=0.3`$. The contribution of the discrete and continuous spectra are plotted separately for this potential. In Fig. 2a the relative mass shift $`\mathrm{\Delta }_2m/m`$ is shown where the mean mass $`m`$ also depends on the potential width $`L`$. Thus, the ratio $`\delta _2m=\mathrm{\Delta }_2m/\mathrm{\Delta }_2m_{in}`$ of the mass shifts in the potential and in $`GaAs`$ is presented also in Fig. 2b for the same $`AlAs`$ mole fraction. Evidently, the asymptotics of this curve is equal to the unity at large $`L`$ and to the ratio $`m_{out}\alpha _{out}/m_{in}\alpha _{in}`$ at $`L0`$.
We may conclude that the continuous spectrum dominates at small potential widths. At large widths its contribution could also be significant although it is smaller than the contribution of the discrete spectrum (especially in deep potential wells). We also confirm the conclusion of the preceding papers that the leading term approximation is not adequate to describe this system and leads to wrong asymptotics at both small and large potential widths (see the dashed lines in Figs 1, 2).
An example of a multilayered heterostructure is presented. The results for the energy and the effective mass for the polaron in the Rosen-Morse potential well are shown in Fig. 3. For the numerical calculations we fix the value $`V_0=227.9`$ meV in Eqs. (79), (83) which corresponds to the limiting mole fraction at large distances $`x_{\mathrm{}}=lim_z\mathrm{}x(z)=0.3`$. Thus, we obtain the dependence of the mole fraction $`x`$ on the coordinate $`z`$:
$`600(1.155x+0.37x^2)=227.9\mathrm{tanh}^2z.`$ (119)
Now Eqs. (107), (108), (117) allow one to define the dependence of parameters on the coordinate $`z`$ and to calculate the mean characteristics of the heterostructure.
The calculations were completely different in comparison with the rectangular potential: instead of the direct summation over all intermediate states we used the analytical expression for the Green function of the Rosen-Morse potential. The results obtained demonstrate a similar behavior which is also close numerically to the results for the rectangular potential. The polaronic energy and mass shifts for the rectangular quantum well are also plotted here (dashed line) vs. the Rosen-Morse width $`L_{RM}`$ obtained from $`L`$ as is described above. We see that both the energies almost coincide, which gives the opportunity to approximate different quantum wells by the rectangular potential. The discrepancy in the effective mass is larger but not so crucial. This serves also as an additional internal criterion of the validity of our calculations.
Thus, we obtained a monotonous behavior of $`\mathrm{\Delta }_2E`$ between the correct 3D limiting values $`\alpha _{in}\mathrm{}\omega _{in}`$ and $`\alpha _{out}\mathrm{}\omega _{out}`$ both for the rectangular and the Rosen-Morse potentials (see Fig. 1 and Fig. 3a). Actually the peaks are “hidden” and they reveal themselves if we plot the dimensionless energy shift $`\mathrm{\Delta }_2E/(\mathrm{}\omega _{\text{LO}}\alpha )`$ which has the same 3D limit (the unity) at both small and large potential widths. But in the “real” units (meV) the peaks are smoothed.
To compare our results with the calculations performed for the one-layer heterostructure we refer to the papers where the authors took into account the contributions of different phonon modes as well as mass and dielectric constant mismatches in the materials of the barrier and the well. Note that the analytical formulas of Ref. contain a mistake — the wrong expression for the density of states. Namely, in some parts of the continuous spectrum contribution the integration is performed not over the wave vector $`p`$ but over the wave vector $`q`$ (that is $`_{V_0}^{\mathrm{}}𝑑E_z/\sqrt{E_z}(\mathrm{})`$ in the notations of that paper instead of the correct integration $`_{V_0}^{\mathrm{}}𝑑E_z/\sqrt{E_zV_0}(\mathrm{})`$). It is clear that this mistake results in lowering of the resulting curve for the energy, and the discrepancy is larger when the energy is closer to the potential edge, that is at small widths. This is just what we see in Fig. 4a comparing the result of Ref. (the curve $`\mathrm{\Delta }_{HPD}E`$) with the new calculations of the same authors (the curve $`\mathrm{\Delta }_{HPD}E`$) which came to our knowledge when the present paper was already submitted for the publication.
Thus, our model does not reproduce the more complicated structure with the peak and the dip which was obtained in Ref. . Some hints on the existence of peaks can also be seen in our plots but the maximal values are so close to the asymptotics that the peaks are almost invisible. Probably, the dip appears because of the presence of several phonon modes (bulk, interface, etc.). At widths $`L50\AA `$ our results for the energy practically coincide with those of Ref. . The discrepancy at smaller widths seems to be more crucial. But the difference between the values in the peak and in the dip for the curve $`\mathrm{\Delta }_HE`$ in Fig. 4a is about 0.1 eV (3%). This phenomena hardly can be seen experimentally and this discrepancy is in the limits of the accuracy of our model estimated above. This gives indeed a strong support to our model and we may conclude that the latter provides us with the rather accurate approximation and can be used for more complicated calculations in multilayered heterostructures.
As to the shift of the electron band mass we found clear peaks for both the rectangular and the Rosen-Morse potentials (see Figs. 2, 3). As is seen in Fig. 2 the effective mass shift for the polaron in the rectangular quantum well has a peak at $`L20\AA (x=0.3)`$. Calculations show also that the larger is $`x`$ the smaller is the potential width corresponding to the peak. For the Rosen-Morse potential at $`x_{max}=0.3`$ the peaks in the effective mass occur at $`2L_{RM}20\AA `$. Note, that again the authors of Ref. obtained curves with peaks and dips in contrast with our results (see Fig. 4b). The maximal discrepancy for the mass is about 11% at $`L3\AA `$ which is beyond the region available for experiments. Our results are very close to those of Ref. at $`L10\AA `$ and practically coincde with them at $`L20\AA `$.
To compare our results with those of Ref. we need now another parametrization used by these authors (although they refer also to the paper by Adachi). Namely, they took a slightly different expression for the confining potential:
$`V(z)=600(1.266x+0.26x^2)\mathrm{meV},`$ (120)
which follows from the band gap of Ref. . Furthermore, instead of the effective LO-phonon frequency they used the expression (110) for the energy of the $`GaAs`$-like phonons. The Fröhlich coupling constant $`\alpha `$ was calculated then using also the parametrization (107) and (V). Note, that these numerical calculation, as we found, can be approximated by the interpolation formula
$`\alpha (z)=0.0797+0.0772x+0.0295x^2.`$ (121)
The results of the comparison are shown in Fig. 5 (we used in our calculations for this plot the same parametrization as was used in Ref. ).
The curve $`\mathrm{\Delta }_{Ch}E`$ in Fig. 5a for the energy shift taken from Ref. has also a small dip (qualitatively similar to this of Ref. ). But the discrepancy between energy shifts is much more drastic in this case, and we have no explanation for this.
It is clear that at large potential width only a bulk phonon mode inside the quantum well contributes so these curves should have the same limiting value $`\alpha _{in}\omega _{in}`$. Numerically we found $`\alpha _{in}=0.0797`$ and $`\omega _{in}=36.25`$ meV, so $`\alpha _{in}\omega _{in}=2.89`$ meV. Moreover, the behavior of the curves at large $`L`$ should be qualitatively and quantitatively the same which was the case when we compared our model with Refs. . In contrast with our model and the cited results by Hai, Peeters and Devreese the curve $`\mathrm{\Delta }_{Ch}E`$ in Fig. 5b approaches the asymptotics from below and the subsequent mechanism remains unclear. On the other hand there are some reasons why the curve have to approach its asymptotics from above. Indeed, at large potential width the particle does not feel yet the finite height of the potential, and the energy shift takes the same value as in the infinitely high potential which is a bit larger than the free polaron energy.
As to the opposite limit of the small width of the confining potential, it is surprising that the asymptotic value is not reached even at $`L0.3\AA ,`$ as is found in Ref. . Numerically we obtained $`\alpha _{out}=0.1014`$ in this scheme and $`\omega _{out}=34.72`$ meV, so $`\alpha _{out}\omega _{out}=3.52`$ meV.
Both asymptotic values coincide with what was obtained by the authors of Ref. . Looking at the behavior of the mass shift, we see that both curves coincide at large widths as it should be. At widths smaller than 100 Å the discrepancy becomes evident. But we may conclude that something is wrong with the numerical job of Ref. because their curve approaches the wrong limit at $`L0`$. Indeed, in this limit the asymptotical value of the plotted ratio should be equal to $`\alpha _{out}m_{out}/\alpha _{in}m_{in}`$. As it follows from our analysis of the energy shift, we obtained the same values for the Fröhlich coupling constants. The values for the band masses follow from Eq. (107): $`m_{in}=0.0665m_e`$ and $`m_{out}=0.0874m_e`$ at $`x=0.25`$. Then, the asymptotical value of the plotted ratio should be equal to $`1.67`$, instead of 1.83 what was found in Ref. That is, the discrepancy is about 10% in this limit and we cannot explain its origin as well.
It would be highly desirable to include a comparison of our results and the results by other authors with corresponding experiments. To the best of our knowledge, no such experiments do exist at the moment.
## VI Conclusions
To conclude, we suggested an approximate model to describe a multilayered $`GaAs/Al_xGa_{1x}As`$ heterostructure as an effective medium with one (bulk) phonon mode. The fundamental entity is the confining potential generated by these layers which we take into account explicitly. Then we calculate the mean characteristics of the electron in the effective medium (such as its band mass, phonon frequencies etc.) which depend on the form of the confining potential. With these parameters we calculated the energy and the effective mass of a polaron confined to a quasi-2D quantum well $`GaAs/Al_xGa_{1x}As`$ for different $`AlAs`$ mole fractions. The calculations include the full energy spectrum as intermediate states. Peaks are found for the effective mass at some potential widths while the energy demonstrates rather monotonous behavior between the correct 3D-limits. Finally, some discrepancies in the interpolation formulae for the experimental results are discussed as well as discrepancies with the previously obtained theoretical results. We demonstrated that our model gives practically the same (or very close) results as the explicit calculations of Ref. for potential widths $`L10\AA `$.
###### Acknowledgements.
We thank J. T. Devreese, V. N. Gladilin, G. Q. Hai, H. Leschke, V. M. Fomin, F. M. Peeters, E. P. Pokatilov, and J. Wüsthoff for useful discussions and valuable remarks and advices. Special thanks are to the authors of the paper for making their results available to us prior to publication and to the authors of the paper who kindly provided us with their data-files. M.O.D. and M.A.S. are grateful to Dortmund University for the kind hospitality during their visits to Germany. Financial support of the Heisenberg-Landau program (Germany-JINR collaboration in theoretical physics) and Deutsche Forschungsgemeinschaft (Graduiertenkolleg GKP 50/2) is gratefully acknowledged.
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# Schematic phase diagram and collective excitations in the collisional regime for trapped boson-fermion mixtures at zero temperature
## 1 Introduction
The achievement of Bose-Einstein condensation in trapped atomic gases \[1 - 3\] has given new impulse to the study of dilute quantal fluids. Experimental attention has extended to the properties of double Bose condensates and more recently to the attainment of quantal degeneracy in gases of fermionic atoms \[6 - 10\] and in dilute mixtures of bosons and fermions . A specific motivation for interest in boson-fermion mixtures is that the collisions between fermions and bosons can foster thermalization of the fermionic component and induce so-called sympathetic cooling . Elastic collisions between fermions in a single hyperfine state are inhibited by the Pauli principle and even in multi-component Fermi gases the collisions between atoms in different hyperfine states become ineffective at very low temperature because of Fermi factors .
A number of theoretical studies have already been addressed to properties of dilute boson-fermion mixtures confined in harmonic traps. They have concerned the equilibrium density profiles of the two species both at zero and at finite temperature \[14 - 17\], the kinetic energy of the fermionic component and collective excitations from sum rules in the collisionless regime . Our main focus here is on small-amplitude collective excitations for a mixture under harmonic confinement in the collisional regime. For a detailed discussion on how this regime may be achieved in dilute Fermi gases we refer to the work of Amoruso et al. .
We first set up a general background for our discussion of collective excitations in boson-fermion mixtures by constructing a schematic phase diagram from the earlier studies of the equilibrium state \[14 - 18\]. This helps to identify various dynamical behaviours, which may later be examined in detail by explicit numerical solution of the equations of motion. We then report analytic results for the eigenfrequencies of the mixture in limiting situations, corresponding to mixing of Bogolubov sound and fermionic first sound in a nearly homogeneous mixture and to surface modes in a regime of weak fermion-boson coupling.
We assume throughout a positive boson-boson scattering length, as is applicable to Bose condensates of <sup>87</sup>Rb and <sup>23</sup>Na. Mixtures of <sup>6</sup>Li and <sup>7</sup>Li should have special interest, once condensation is realized for <sup>7</sup>Li in a hyperfine state where the scattering length is positive. According to calculations on the collisional properties of ultracold potassium , the <sup>40</sup>K-<sup>41</sup>K mixture should also be an interesting system for study from the present viewpoint, if condensation of the rare isotope <sup>41</sup>K can be achieved.
## 2 Equilibrium density profiles and phase diagram
We consider a system of $`N_f`$ fermions of mass $`m_f`$ and $`N_b`$ bosons of mass $`m_b`$ at $`T=0`$, confined by spherically symmetric external potentials $`V_{f,b}(r)=m_{f,b}\omega _{f,b}^2r^2/2`$ with frequencies $`\omega _f`$ and $`\omega _b`$. We assume that a single spin state is trapped for each component of the mixture and that all bosons are in the condensate. We can then omit the fermion-fermion interaction, which is inhibited by the Pauli exclusion principle. The boson-boson and boson-fermion interactions are described by contact potentials with scattering lengths $`a_{bb}>0`$ and $`a_{bf}`$, yielding coupling-strength parameters $`g=4\pi \mathrm{}^2a_{bb}/m_b>0`$ and $`f=2\pi \mathrm{}^2a_{bf}/m_r`$ where $`m_r`$ is the reduced boson-fermion mass.
The equations of motion for the partial particle densities $`n_\sigma (𝐫,t)`$ and partial current densities $`𝐣_\sigma (𝐫,t)`$, where $`\sigma `$ is an index for the atomic species ($`\sigma =f`$ or $`b`$) follow from those for the one-body density matrix upon projection on the main diagonal (see for instance ). They take the form of continuity equations,
$$_tn_\sigma (𝐫,t)=𝐣_\sigma (𝐫,t)$$
(2.1)
and of generalized hydrodynamic equations which, in a mean-field approximation and neglecting coupling to energy fluctuations, can be written as
$$m_f_t𝐣_f(𝐫,t)=\mathrm{\Pi }^{(f)}(𝐫,t)n_f(𝐫,t)\left[V_f(𝐫)+fn_b(𝐫,t)\mu _f\right]$$
(2.2)
and
$`m_b_t𝐣_b(𝐫,t)`$ $`=`$ $`\mathrm{\Pi }^{(b)}(𝐫,t)`$ (2.3)
$`n_b(𝐫,t)\left[V_b(𝐫)+gn_b(𝐫,t)+fn_f(𝐫,t)\mu _b\right]`$
In these equations $`\mathrm{\Pi }^{(\sigma )}(𝐫,t)`$ are the kinetic stress tensors and $`\mu _\sigma `$ are the chemical potentials of the two atomic species. In the following we treat the kinetic stress tensors by a local density approximation for fermions,
$$\mathrm{\Pi }_{ij}^{(f)}(𝐫,t)=\frac{2}{5}A[n_f(𝐫,t)]^{5/3}\delta _{ij}$$
(2.4)
with $`A=\mathrm{}^2(6\pi ^2)^{2/3}/2m_f`$ and by a Thomas-Fermi approximation for bosons,
$$\mathrm{\Pi }^{(b)}(𝐫,t)=0.$$
(2.5)
In these equations we have dropped terms which are quadratic in the velocity field.
The treatment given above assumes that local equilibrium has been established for the ground-state density profiles and that it is maintained during dynamic fluctuations of the particle densities. The dynamical behaviour that we shall discuss in Sect.3 will therefore correspond to a collisional regime .
### 2.1 Equilibrium density profiles
The equilibrium density profiles ($`n_\sigma ^{eq}(r)`$, say) follow from Eqs. (2.2) - (2.5) by setting the current densities to zero. The results are well known :
$$n_b^{eq}(r)=g^1\left[\mu _bV_b(r)fn_f^{eq}(r)\right]\theta (\mu _bV_b(r)fn_f^{eq}(r))$$
(2.6)
and
$$n_f^{eq}(r)=A^{3/2}\left[\mu _fV_f(r)fn_b^{eq}(r)\right]^{3/2}\theta (\mu _fV_f(r)fn_b^{eq}(r)).$$
(2.7)
The chemical potentials are to be determined by self-consistency conditions on the particles numbers, $`N_\sigma =n_\sigma ^{eq}(r)𝑑𝐫`$. On comparing the two density profiles in the limit of vanishing boson-fermion coupling, it is easily seen that the quantity $`(2A/3)[n_f^{eq}(r)]^{1/3}`$ may be viewed as an effective fermion-fermion coupling arising from the kinetic pressure of the Fermi gas . It is also useful to introduce the radii $`R_b=(2\mu _b/m_b\omega _b^2)^{1/2}`$ and $`R_f=(2\mu _f/m_f\omega _f^2)^{1/2}`$ for the two clouds in the absence of interactions, but taken at the true values of the chemical potentials.
For values of the boson-boson coupling and of the confinement frequencies which are typical of current experiments, and assuming comparable magnitudes for the boson-boson and the boson-fermion coupling, the fermionic cloud is considerably more dilute than the bosonic cloud. In this situation the boson cloud is essentially unaffected by the interactions with the fermionic component and the latter can be treated as confined inside a double-parabola effective potential $`V_f^{eff}(r)`$ . This is given by
$$V_f^{eff}(r)=\{\begin{array}{cc}m_f\omega _f^2(1\gamma )r^2/2+f\mu _b/g\hfill & \hfill \mathrm{for}r<R_b\\ m_f\omega _f^2r^2/2\hfill & \hfill \mathrm{for}r>R_b\end{array}$$
(2.8)
where
$$\gamma =(fm_b\omega _b^2)/(gm_f\omega _f^2).$$
(2.9)
Evidently, if $`\gamma >1`$ the effective potential (2.8) has negative concavity inside the Bose radius and the fermions will tend to be expelled from the centre of the trap. The result at strong coupling is phase separation in the mixture, with the fermions forming a shell around the boson cloud .
A second type of phase separation for the boson-fermion mixture was demonstrated in the work of Mølmer . In this case one takes $`N_fN_b`$ and considers very large values of the boson-boson coupling $`g`$, such that the two clouds acquire similar densities. One ultimately finds that the bosons are expelled from the centre of the trap and form a shell around the fermion cloud.
Finally, Eq. (2.8) is valid also in the case $`\gamma <0`$, i.e. for attractive boson-fermion interactions, as long as the fermionic cloud is more dilute than the boson cloud. However, the fermions now tend to draw the bosons towards the centre of the trap, ultimately leading to collapse once $`\gamma `$ becomes sufficiently negative .
### 2.2 Schematic phase diagram at zero temperature
As shown by Vichi et al. , the thermodynamics of the boson-fermion mixture at $`T=0`$ is well characterized by means of two scaling parameters, which are $`\gamma `$ in Eq. (2.9) and
$$x=\left(\frac{R_b}{R_f}\right)^{1/2}=\left(\frac{m_f\omega _f}{2m_b\omega _b}\right)^{1/2}\left[15\frac{a_{bb}}{a_{ho}}\frac{N_b}{(6N_f)^{5/6}}\right]^{1/5},$$
(2.10)
with $`a_{ho}=(\mathrm{}/m_b\omega _b)^{1/2}`$. The description of the system with only two scaling parameters, instead of the eight original ones entering the expression of the equilibrium density profiles, is a major simplification. The ratio $`R_b/R_f`$ determines the deviation of the kinetic energy of the fermion cloud from its ideal-gas value .
For the purposes of the present discussion we make the simplifying assumptions $`m_f=m_b`$, $`\omega _f=\omega _b`$ and $`N_f=N_b`$. In this case the two scaling parameters take the simple expressions $`x=(15a_{bb}k_f^{(0)}/48)^{1/5}`$ and $`\gamma =a_{bf}/a_{bb}`$, with $`k_f^{(0)}=(48N_f)^{1/6}/a_{ho}`$ being the Fermi wave number at the trap centre for a non-interacting Fermi gas. As remarked under Eq. (2.7) above, the quantity $`a_{ff}3\pi /4k_f^{(0)}`$ may be viewed as an effective fermion-fermion scattering length arising from the Pauli kinetic pressure. We can now proceed to construct a schematic phase diagram for the boson-fermion mixture in the plane defined by the variables $`\gamma =a_{bf}/a_{bb}`$ and $`y=a_{bb}k_f^{(0)}`$. Figure 1 shows the result for the case $`\gamma >0`$ (i.e. repulsive boson-fermion interactions).
We search for the boundary of phase separation by means of a simple condition of linear stability on the two-by-two matrix of scattering lengths $`a_{bb}`$, $`a_{bf}`$ and $`a_{ff}`$ (see also ). A mixed state is stable if the inequality $`a_{bb}a_{ff}a_{bf}^2>0`$ holds, i.e. when
$$k_f^{(0)}a_{bb}\frac{3\pi }{4}\left(\frac{a_{bb}}{a_{bf}}\right)^2$$
(2.11)
Within this region we can distinguish a region of appreciable overlap between the two clouds (region I in Figure 1), bounded by two regions of diminishing overlap: a region at $`\gamma >1`$ where the fermions are being expelled from the centre of the trap (region II in Figure 1) and by a region at $`y>4/3\pi `$ where the bosons are being expelled from the centre of the trap (region III in Figure 1). If the spherical symmetry of the trapping potential is preserved in the demixed state, these two regions end into spherical-shell structures with the fermions at large $`\gamma `$ (or the bosons at large $`y`$) forming the outer shell.
## 3 Small amplitude oscillations around the equilibrium profiles
The equations of motion for small-amplitude density fluctuations $`\stackrel{~}{n}_f(𝐫,t)`$ and $`\stackrel{~}{n}_b(𝐫,t)`$ in the collisional regime are found by combining Eq. (2.1) with the linearized form of Eqs. (2.2) and (2.3). They are
$$_t^2\stackrel{~}{n}_f(𝐫,t)=\frac{1}{m_f}\left\{n_f^{eq}(r)\left[\frac{2}{3}A(n_f^{eq}(r))^{1/3}\stackrel{~}{n}_f(𝐫,t)+f\stackrel{~}{n}_b(𝐫,t)\right]\right\}$$
(3.1)
and
$$_t^2\stackrel{~}{n}_b(𝐫,t)=\frac{1}{m_b}\left\{n_b^{eq}(r)\left[g\stackrel{~}{n}_b(𝐫,t)+f\stackrel{~}{n}_f(𝐫,t)\right]\right\}.$$
(3.2)
Again the quantity $`(2A/3)[n_f^{eq}(r)]^{1/3}`$ enters Eq. (3.1) as an effective fermion-fermion coupling.
It is easily checked that the form of Eqs. (3.1) and (3.2) is such as to satisfy the generalized Kohn theorem for the centre-of-mass coordinate $`𝐱(t)`$ of the whole fluid when the two confinement frequencies coincide. In this mode of motion we have $`n_{b,f}(𝐫,t)=n_{b,f}^{eq}(𝐫𝐱(t))`$ and hence
$$\stackrel{~}{n}_{b,f}(𝐫,t)=𝐱(t)n_{b,f}^{eq}(r)$$
(3.3)
for small amplitude oscillations. By substituting Eq. (3.3) in Eqs. (3.1) and (3.2) and using Eqs. (2.6) and (2.7), we find
$$m_{b,f}_t^2𝐱(t)n_{b,f}^{eq}(r)=\underset{i,j}{}x_j(t)_i\left[n_{b,f}^{eq}(r)_i_jV_{b,f}(r)\right].$$
(3.4)
Therefore, the centre of mass of the system oscillates at the bare trap frequency.
Evidently, a general solution of Eqs. (3.1) and (3.2) can only be obtained numerically, since the equilibrium density profiles are not known analytically from Eqs. (2.6) and (2.7). The dynamics of the fluid as a function of the scaling parameters will be especially rich in the regions of the phase diagram in which the nature of the partial density profiles is changing on the approach to phase separation. Of course, this will be signalled by a softening of the modes associated with concentration fluctuations. Below we restrict ourselves to analytic results which are applicable to special regimes.
### 3.1 Homogeneous limit
The dynamics of a spatially homogeneous mixture is relevant to the situation in which the density profiles of both components are slowly varying in space, on the length scale set by the mean interparticle distances. Taking the equilibrium densities as constant, it is easily seen from Eqs. (3.1) and (3.2) that the eigenmodes are two sound waves, with linear dispersion relation $`\omega _{1,2}(q)=c_{1,2}q`$. The sound velocities are given by
$$c_{1,2}^2=\frac{1}{2}\left[c_b^2+c_f^2\pm \sqrt{(c_b^2c_f^2)^2+4f^2n_b^{eq}n_f^{eq}/m_bm_f}\right].$$
(3.5)
In Eq. (3.5) $`c_b=(gn_b^{eq}/m_b)^{1/2}`$ is the speed of Bogolubov sound and $`c_f=[2A(n_f^{eq})^{2/3}/3m_f]^{1/2}`$ is that of first sound in a Fermi gas.
This result may be useful in regard to measurements of the speed of sound in elongated systems, in experiments such as already carried out by Matthews et al. on a single Bose condensate.
### 3.2 Surface modes for weak boson-fermion coupling
Assuming equal confinement frequencies for the two components of the mixture ($`\omega _{ho}`$, say), it is known from earlier results of Stringari for a Bose condensate and of Amoruso et al. for a one-component Fermi gas that in the limit $`f=0`$ the frequencies of the surface modes ($`n`$= 0) with angular momentum number are equal and given by $`\omega _l=l^{1/2}\omega _{ho}`$. A perturbative treatment of the effect of fermion-boson coupling is therefore especially simple for these modes.
We define the unperturbed equilibrium density profiles $`n_{b,f}^{(0)}(r)`$ and the $`𝐫`$-dependent amplitudes $`\stackrel{~}{n}_{b,f}^{(l,m)}(𝐫)`$ of the unperturbed density fluctuations, for each value of the angular momentum $`l`$ and of its $`z`$-component $`m`$. After Fourier-transforming Eqs. (3.1) and (3.2) and linearizing them in the coupling parameter $`f`$, we multiply Eq. (3.1) by $`[n_f^{(0)}(r)]^{1/3}\stackrel{~}{n}_f^{(l,m)}(𝐫)`$ and Eq. (3.2) by $`\stackrel{~}{n}_b^{(l,m)}(𝐫)`$, and then integrate both equations over $`𝐫`$. We obtain in this way a determinantal equation for the shift of the eigenmode frequencies $`\mathrm{\Omega }_l`$ due to the fermion-boson coupling. Writing
$$\mathrm{\Omega }_l^2=\omega _{ho}^2(l+\mathrm{\Delta }_l),$$
(3.6)
we find
$$\mathrm{\Delta }_l=\frac{a_{bf}}{2a_{ho}}\left[(E+B)\pm \mathrm{sign}(a_{bf})\sqrt{(EB)^2+4CD}\right],$$
(3.7)
this result being valid for either sign of $`a_{bf}`$. The quantities entering Eq. (6) are defined in the Appendix, where it is also shown that they can all be expressed in terms of the Gauss hypergeometric function .
The result given in Eq. (3.7) should be better applicable in the case $`R_b<R_f`$, since in deriving it we have assumed that the equilibrium density profiles for both components are perturbed in a symmetric manner. This corresponds to working in region I of the phase diagram (and in the corresponding region at negative $`\gamma `$).
## 4 Summary and discussion
The main features of the zero-temperature phase diagram that we have sketched in Figure 1 for boson-fermion mixtures arise from the competition between the kinetic energy of the Fermi gas and the repulsive boson-boson and boson-fermion interactions. The former disfavours localization and phase separation, while the latter favour the spatial separation of the two components. On varying these system parameters while preserving the spherical symmetry of the trapped fluid mixture, a spontaneous symmetry breaking occurs towards a situation where the two components are spatially separated along the radial direction. Evidently, axial separation of the two components would instead be observed in an anisotropic trap. One may think of exploiting Feshbach resonances to tune the scattering lengths towards such phase-separation regime.
Given the complex behaviour of the density profiles on the approach to phase separation, a complete investigation of the oscillatory eigenmodes of a boson-fermion mixture can only be carried out by numerical means. This will be well worth doing once a mixture of specific experimental relevance and with reasonably known scattering lengths is identified. Of course, the dynamics in the phase-separated region is simply related to that of the two pure components, aside for the presence of interfacial modes. On the other hand, we have seen that suggestive analytic results can be obtained in the quasi-homogeneous limit and in the weak boson-coupling regime, for either sign of the boson-fermion coupling. With increasingly large attractions between boson and fermions, however, the mixture is expected to undergo collapse .
We thank Dr Ilaria Meccoli for many useful discussions and for her help in the early stages of this work.
## Appendix A Calculation of frequency shifts for surface modes
We introduce the notations $`|F=\stackrel{~}{n}_f^{(l,m)}(𝐫)`$ and $`|B=\stackrel{~}{n}_b^{(l,m)}(𝐫)`$, define the operators $`P_{b,f}=(f/m_{b,f})[n_{b,f}^{(0)}(r)]`$ and $`D_f=(f/m_f)\{(n_f^{(0)}(r))^{1/3}[(n_f^{(0)}(r))^{1/3}]\}`$, and introduce the scalar-product notation $`F|\widehat{O}|F=\stackrel{~}{n}_f^{(l,m)}(𝐫)\widehat{O}n_f^{(l,m)}(𝐫)𝑑𝐫`$ etcetera. With these notations we find the following expressions for the quantities entering Eq. (3.7):
$`B`$ $`=`$ $`F|(n_f^{(0)}(r))^{1/3}D_f|F/F|(n_f^{(0)}(r))^{1/3}|F`$ (A.1)
$`=`$ $`{\displaystyle \frac{2\pi ^{1/2}l}{\alpha _{bb}X_f^{2l+2}}}{\displaystyle \frac{\mathrm{\Gamma }(l+3)}{\mathrm{\Gamma }(l+3/2)}}{\displaystyle _0^{X_b}}𝑑xx^{2l+2}(X_f^2x^2)^{1/2}(2X_f^2+X_b^23x^2),`$
$`C`$ $`=`$ $`F|(n_f^{(0)}(r))^{1/3}P_f|B/F|(n_f^{(0)}(r))^{1/3}|F`$ (A.2)
$`=`$ $`{\displaystyle \frac{2l}{\pi ^{5/2}X_f^{2l+2}}}{\displaystyle \frac{\mathrm{\Gamma }(l+3)}{\mathrm{\Gamma }(l+3/2)}}{\displaystyle _0^{X_b}}𝑑xx^{2l+2}(X_f^2x^2)^{1/2},`$
$`D`$ $`=`$ $`B|P_b|F/B|B`$ (A.3)
$`=`$ $`{\displaystyle \frac{(2l+3)}{2\alpha _{bb}X_b^{2l+3}}}{\displaystyle _0^{X_b}}dxx^{2l+2}\{2l(X_f^2x^2)^{1/2}+X_f^2(X_f^2x^2)^{3/2}(X_b^2x^2)`$
$`+(X_f^2x^2)^{1/2}[x^2+2(l+1)(X_b^2x^2)]\},`$
and
$`E`$ $`=`$ $`B|P_f|B/B|B`$ (A.4)
$`=`$ $`{\displaystyle \frac{l(2l+3)}{2\pi ^2X_b^{2l+3}}}{\displaystyle _0^{X_b}}𝑑xx^{2l+2}(X_f^2x^2)^{1/2}.`$
In these equations $`\mathrm{\Gamma }(n)`$ is the Gamma function and we have set $`X_{b,f}=R_{b,f}/a_{ho}`$ and $`\alpha _{bb}=a_{bb}/a_{ho}`$.
All the integrals in Eqs. (A.1) -(A.4) can be expressed through the Gauss hypergeometric function $`{}_{2}{}^{}F_{1}^{}(a,b;c;z)`$, using the relation
$${}_{2}{}^{}F_{1}^{}(\frac{1+\alpha }{2},\frac{n}{2};\frac{3+\alpha }{2};\frac{X_b^2}{X_f^2})=\frac{1+\alpha }{X_f^nX_b^{\alpha +1}}_0^{X_b}𝑑xx^\alpha (X_f^2x^2)^{n/2}.$$
(A.5)
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# 1 The law of conservation of 𝜎-number.
## 1 The law of conservation of $`\sigma `$-number.
It is known that upper and lower quarks are positioned asymmetrically on the ”charge axis”. With the aim to bring in some symmetry to this asymmetric disposition we will introduce the new additive quantum number $`\sigma `$ determined so as to result, in combination with quark electric charge $`q`$, the charge of the respective lepton.
Then for $`u`$ and $`d`$ quarks we will have, respectively, $`\sigma _u=1/3`$ and $`\sigma _d=2/3`$. In general, $`\sigma `$-numbers for all quarks and anti-quarks are determined by the following formula:
$$\sigma =q1/3,\stackrel{~}{\sigma }=\stackrel{~}{q}+1/3.$$
(1)
Now the quarks (and leptons) are positioned symmetrically with respect to $`\sigma =q`$ axis on $`(q,\sigma )`$ plane (see Fig.1). This brings about an idea that together with the law of conservation of electric charge there is realized the second law of conservation, the law of conservation of $`\sigma `$-number. As we will see later, this idea results in extremely valuable findings. In particular, the conservation of $`\sigma `$ -number results in conservation of the difference between baryon and lepton numbers obtained earlier by Georgi and Glashow from $`SU(5)`$ symmetry.
Let us first show that $`\sigma `$-numbers for baryons, mesons, leptons and photons are determined by the following expressions:
$$\sigma _B=Q_B1,\sigma _M=Q_M,\sigma _L=Q_L+1,\sigma _\gamma =0,\stackrel{~}{\sigma }=\sigma $$
(2)
where $`Q`$ is an electric charge of a particle. Indeed, formula (2) is obvious for baryons and mesons due to their quark structure. Similarly, $`\sigma `$-numbers can be obtained for exotic adrons. Using (2), we will obtain, in particular, $`\sigma (p)=0,\sigma (n^0)=1,\sigma (\pi ^0)=0,\sigma (\stackrel{~}{\mathrm{\Lambda }}^0)=1,\sigma (^{})=2`$ and so on. Taking into account that $`\sigma (\pi ^0)=0`$, we will obtain $`\sigma _\gamma =0`$from $`\pi ^02\gamma `$decay.
In order to find $`\sigma _L`$, we will give the same quantum number (by analogy with the electric charge) to all charged leptons $`(\sigma _e=\sigma _\mu =\sigma _\tau )`$and another quantum number to neutral leptons $`(\sigma _{\nu _e}=\sigma _{\nu _\mu }=\sigma _{\nu _\tau })`$ . These numbers are related to each other by the formula $`\sigma _{\nu _e}=1+\sigma _e`$that can be derived from $`n^0p+e^{}+\stackrel{~}{\nu }`$decay. Then from $`uuxe^+\stackrel{~}{d}`$reaction where the same boson may decay into antilepton + antiquark or quark pair we derive $`\sigma _{e^+}=\sigma _e^{}=0`$ , and $`\sigma _\nu =1`$. As we can see, these results are in accordance with expressions (2). After we have derived formula (2), we can check out that the conservation of $`\sigma `$-number takes place for all the reactions observed so far . It does not prohibit also possible proton decay through channels $`pe^+\pi ^0;e^+\pi ^+\pi ^{}`$ where conservation of baryon and lepton numbers is violated.
The general formula for conservation of $`\sigma `$-number can be written as follows:
$$Q_BN_B+\stackrel{~}{Q}_B+\stackrel{~}{N}_B+Q_M+\stackrel{~}{Q}_M+Q_L+N_L+\stackrel{~}{Q}_L\stackrel{~}{N}_L=const$$
(3)
where $`N`$ is the number of particles. Here, taking into account the conservation of electric charge, we will finally obtain the desired result:
$$(N_B\stackrel{~}{N}_B)(N_L\stackrel{~}{N}_L)=const$$
(4)
It should be pointed out that the baryon number in formula (4) could be replaced by total number of quarks$`N_k`$ comprised of baryon and meson quark numbers. Indeed, taking into account that for baryons and mesons $`(N_k\stackrel{~}{N}_k)=3(N_B\stackrel{~}{N}_B)`$and $`(N_k\stackrel{~}{N}_k)=0`$and using formula (4), we will obtain the following relation between quarks and leptons
$$(N_k\stackrel{~}{N}_k)3(N_L\stackrel{~}{N}_L)=const$$
(5)
As we will see later, the law of conservation of $`\sigma `$-number predicts that electron type neutrino mass is exactly zero. It should be mentioned in this relation that two contradictory theories trying to resolve the problem of neutrino mass have been developed so far by Dirak-Weyl in 1929 and Majoran in1936. According to the former electron neutrino mass is exactly zero and the conservation of lepton number takes place. According to the latter the electron neutrino mass is not zero and therefore the conservation of lepton number does not take place. The Majoran theory allows neutrinoless double beta decay where two $`d`$-quarks decay through channel $`du+e^{}`$with violation of lepton number.
The contemporary theories basing on $`SU(5)`$ and $`SO(10)`$ symmetry result in the same contradictory predictions ”The $`SO(10)`$ symmetry allows occurrence of some phenomena prohibited by $`SU(5)`$ symmetry. In particular, $`SU(5)`$ theory predicts conservation of difference $`BL`$ while conservation of baryon number $`B`$ and lepton number $`L`$does not take place. In $`SO(10)`$ theory, the difference $`BL`$ may not hold constant if sufficient number of Higgs fields are involved”. Approximately the same results are mentioned in the Gell-Mann report.
Thus, the main question relating to neutrino mass remains open. There have been a number of attempts for the past 60 years to find theoretical or experimental solution but the problem still exists. The recent experiments in many scientific centers worldwide have been aimed at observing of neutrinoless double beta decay and neutrino oscillations. Currently several new projects are proposed and huge scientific potential and funds mobilized to get a breakthrough on this fundamental problem. Twenty-three scientific centers are involved in the project MINOS (USA). The underground experiment is aimed at neutrino beams detecting after passing some 730km from Fermi-lab (Viskonsin) to Soudan (Minnesota). It is expected that some of the detected neutrinos will have changed their flavor due to neutrino oscillations. In Japanese project KEK neutrino beam will pass some 230km from an accelerator to the detector Super-Kamiokande. The new Heidelberg project GENIUS is aimed at increasing the sensitivity of Majorana neutrino mass from the present level of at best 0.1ev down to 0.01 or even 0.001ev.
Summarizing, we can say that the newly proposed law of conservation of $`\sigma `$-number prohibits neutrinoless double beta decay and, in compliance with Dirac-Weyl and $`SU(5)`$ theories, predicts that electron type neutrino mass is exactly zero.
## 2 Quark-lepton symmetry.
Now we will focus on the problem of quark-lepton symmetry. The essence of the problem is to find such a symmetry among the elementary particles, which would explain why each quark-lepton generation is necessarily comprised of the pair of leptons and pair of quarks. Let us shift to the new system of coordinates $`(q^{^{}},\sigma ^{^{}})`$ with the point of origin in the center of a rectangle $`e\nu _eud`$ and axes directed along the axes of symmetry (see Fig1.). The relation with an original system of coordinates is given by the following expressions:
$$q^{^{}}=\frac{1}{\sqrt{2}}q+\frac{1}{\sqrt{2}}\sigma ,\sigma ^{^{}}=\frac{1}{\sqrt{2}}q+\frac{1}{\sqrt{2}}\sigma \frac{1}{3\sqrt{2}}$$
(6)
The coordinates of the vertices are now as follows:
$$e(\frac{1}{\sqrt{2}},\frac{\sqrt{2}}{3}),\nu _e(\frac{1}{\sqrt{2}},\frac{\sqrt{2}}{3}),u(\frac{1}{\sqrt{2}},\frac{\sqrt{2}}{3}),d(\frac{1}{\sqrt{2}},\frac{\sqrt{2}}{3});$$
(7)
The group that represents the symmetry transformation for the above mentioned rectangle, consists of the following elements: $`E`$-identity transformation; turn around axis $`z`$ for angle $`2\pi `$ ; $`A`$-turn around axis $`\sigma ^{^{}}`$ for angle $`\pi `$; $`B`$-turn around axis $`q^{^{}}`$ for angle $`\pi `$; $`C`$-turn around axis $`z`$ for angle $`\pi `$.
For these elements we will obtain the group multiplication table (see Table 1).
Under multiplication we understand the subsequent execution of the corresponding operations. The elements of the group have order 2 (except for identity element ), since $`\chi ^2=E`$and $`\chi ^1=\chi `$, where $`\chi `$is an arbitrary element of the group. Thus, the totality of the elements $`E,A,B,C`$ makes up the Abelian group. The ”self-transformation” of a regular polygon is expressed by means of the following matrices
$$D_k=\left(\begin{array}{cc}\mathrm{cos}\frac{2\pi k}{n},\hfill & \mathrm{sin}\frac{2\pi k}{n}\hfill \\ \mathrm{sin}\frac{2\pi k}{n},\hfill & \mathrm{cos}\frac{2\pi k}{n}\hfill \end{array}\right),U_k=\left(\begin{array}{cc}\mathrm{cos}\frac{2\pi k}{n},\hfill & \mathrm{sin}\frac{2\pi k}{n}\hfill \\ \mathrm{sin}\frac{2\pi k}{n},\hfill & \mathrm{cos}\frac{2\pi k}{n}\hfill \end{array}\right).$$
(8)
where $`k=0,1,2,,n1`$. These $`2n`$-dimensional matrices make up the group of order $`2n`$ known as the group of diedr. In case of n=2 we have the simplest case of the group with elements
$$E=\left(\begin{array}{cc}1\hfill & 0\hfill \\ 0\hfill & 1\hfill \end{array}\right),A=\left(\begin{array}{cc}1\hfill & 0\hfill \\ 0\hfill & 1\hfill \end{array}\right),B=\left(\begin{array}{cc}1\hfill & 0\hfill \\ 0\hfill & 1\hfill \end{array}\right),C=\left(\begin{array}{cc}1\hfill & 0\hfill \\ 0\hfill & 1\hfill \end{array}\right).$$
(9)
Taking into account that the group (9) is isomorphic to the above mentioned group, we come to conclusion that (9) is the matrix representation of the group. Thus, if particles of the same generation are located in the vertices of rectangle (7) on the plane $`(q^{^{}},\sigma ^{^{}})`$, this distribution is subject to the symmetry transformation described by the group (9). It is also easy to show that if distribution of particles on the plane $`(q^{^{}},\sigma ^{^{}})`$ is subject to the symmetry transformation described by the group (9) and if coordinates of any arbitrary particle from (7) coincide with one of the vertices, there should be three more particles (and only three, without taking into account the color of the quarks) whose coordinates coincide with the remaining vertices.
Summarizing this paragraph, we can say that the nature of the quark-lepton symmetry could be described by $`q\sigma `$-symmetry. It is reflected by the fact that particles of the same generation are subject to the symmetry transformation represented by 4-group of diedr.
## 3 Symmetry among colorless elementary particles.
It is natural to expect that there is some symmetry among colorless particles as well. We will consider all the baryons and mesons which can be comprised of eighteen known quarks as well as the leptons. Each point on the plane $`(q,\sigma )`$ is assumed to ”contain” the whole family of the particles (see fig.2). For example, point $`(1;0)`$ contains, besides electron, other charged leptons $`\mu ^{}`$and $`\tau ^{}`$. It is clear from the picture that there is symmetry with respect to $`\sigma ^{^{}}`$axis. The number of particles in symmetrical points is identical. It is easy to show that the same symmetry exists for multi-quark baryons as well. For example, if baryons with quantum numbers $`q=3,\sigma =2,B=5`$ consist of eight upper and seven lower quarks, then symmetrical baryons with quantum numbers $`q=2,\sigma =3,B=5`$ consist of seven upper and eight lower quarks. It is obvious that the number of colorless combinations is identical.
For baryons the same symmetry exists also with respect to other axis $`q^{^{}}`$ which is determined by equation $`\sigma =q3`$. It can be proved by direct count of the particles in corresponding points (see fig.3). The dashed lines parallel to $`q^{^{}}`$ axis are determined by the following formula:
$$\sigma =q(BL)$$
(10)
It should be pointed out that formula (2) can be represented by formula (10). For baryons, instead of (10), we have $`\sigma =qB`$. The point $`q=3,\sigma =3,B=6`$contains only one baryon comprised of all the eighteen known quarks. The distribution of colorless particles shows that, in order to have completed the symmetry, it is necessary to assume some particles not observed so far in the points marked by crosses. For example, there should be three particles in the point with coordinates $`q=4,\sigma =3`$because there are three particles (leptons) in the symmetrical point. Now the distribution of colorless particles will be subject to the symmetry transformation described by 4-group of diedr (9). The predicted particles should be located in five points symmetrically positioned to those of leptons and mesons. These particles, like leptons, should not have quark structure because for them $`B=6;7`$. They are supposedly ten times heavier than baryons. We will call them transbaryons. The number of transbaryons is less by one of the total number of leptons and mesons. If, by any selection rule, it is prohibited for some particles to be located in some arbitrary points, then, in compliance with above mentioned symmetry, the same number of particles should be missing out in symmetrically positioned points. But this does not pertain to transbaryons because symmetrical to them leptons and mesons (most of them) have already been observed. The newly proposed symmetry for colorless particles does not depend on generation. Based on this symmetry alone, we cannot judge about forth generation.
Concluding this paragraph, we can say that within the frame of the proposed symmetry for colorless elementary particles the existence of transbaryons seems undisputable .
Fig 1. Distribution of quarks and leptons of the first generation on $`q,\sigma )`$ plane.
Fig 2. Distribution of ordinary particles on $`(q,\sigma )`$ plane. The distance between neighbouring points along axes $`q`$ and $`\sigma `$ is taken as a measurement unit.
Fig 3. Distribution of colorless particles on $`(q,\sigma )`$ plane. Transbaryons are marked by crosses. The figures indicate the number of particles in each point.
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# Gluino Condensation in an Interacting Instanton Ensemble
## I Introduction
In order to improve our understanding of non-perturbative phenomena in QCD it is useful to view QCD from a larger perspective, as a member of a family of QCD-like theories with different matter contents. In this context we would like to understand the phase structure of QCD-like theories with $`N_f`$ fermions in the fundamental representation of the gauge group and $`N_{ad}`$ fermions in the adjoint representation. Theories with adjoint fermions are special because the action may display a symmetry that connects bosonic and fermionic degrees of freedom, supersymmetry. Supersymmetry imposes powerful restrictions on the structure of the low energy effective action. These constraints have been used, for example, to determine the phase structure of $`N=1`$ supersymmetric QCD with $`N_c`$ colors and $`N_f`$ flavors of fundamental quarks .
In addition to that, supersymmetry provides the opportunity to isolate certain non-perturbative effects, in particular instantons. This idea has been used in order to calculate the gluino condensate in the simplest supersymmetric gauge theory, SUSY gluodynamics. The strategy behind the so called weak coupling instanton (WCI) calculation is to add to the theory a fundamental fermion together with its scalar superpartners and consider the regime where the expectation value of the scalar field is large. In this case there is a unique non-perturbative superpotential induced by instantons. Since the scalar vev is large, instantons are semi-classical and the superpotential can be calculated reliably. The superpotential determines the gluino condensate in the theory with additional matter. Finally, the matter fields can be decoupled by sending their mass to infinity. The result for the gluino condensate in $`SU(2)`$ supersymmetric gluodynamics is
$$\lambda \lambda =6\mathrm{\Lambda }^3,$$
(1)
where $`\mathrm{\Lambda }`$ is the scale parameter defined in
$$\mathrm{\Lambda }=M_{PV}\left(\frac{16\pi ^2}{bg^2}\right)^{1/3}\mathrm{exp}\left(\frac{8\pi ^2}{bg^2}\right).$$
(2)
Here, $`b=3N_c`$ is the first coefficient of the beta function in supersymmetric gluodynamics and $`M_{PV}`$ is a Pauli-Vilars regulator.
There is an old puzzle concerning this result. The puzzle is connected with the fact that there is an alternative method for calculating the gluino condensate, usually referred to as the strong coupling instanton (SCI) method , for a review see . In the case of $`N_c=2`$ one considers a four fermion correlation function. This correlation function is a topological quantity. Not only can it be saturated with one instanton, but supersymmetry implies that the correlator is just a constant. At short distance, one expects that this constant is saturated by small instantons and can be calculated reliably. The gluino condensate is then extracted by using clustering. The puzzle is that the result differs from the WCI calculation by a factor 4/5.
Several suggestions have been put forward in order to resolve the puzzle . We do not wish to discuss these possibilities in detail. Instead, we would like to employ a somewhat different, more qualitative approach. Even though there is no direct instanton contribution to the gluino condensate, one would still expect configurations with instantons and anti-instantons to contribute to the gluino condensate. Here we have in mind that the theory is studied in a finite volume and in the presence of a non-zero mass term. The thermodynamic limit is approached by taking the volume to infinity before we let the mass go to zero. The mechanism for gluino condensation is similar to the instanton liquid picture of quark condensation in ordinary QCD . For $`N_f>1`$ there is no direct instanton contribution to the quark condensate but chiral symmetry breaking may take place in an ensemble of instantons and anti-instantons in the thermodynamic limit. The Banks-Casher relation $`\overline{q}q=\pi \rho (0)`$ connects the quark condensate with the density of eigenvalues of the Dirac operator at zero virtuality. For simplicity let us consider an ensemble with an equal number of instantons and anti-instantons. In this case the Dirac operator no longer has any exact zero modes. However, if the interaction between the instantons is sufficiently weak, the approximate zero modes associated with individual instantons and anti-instantons form a zone around zero virtuality and lead to spontaneous chiral symmetry breaking. This quark condensation mechanism has been investigated in some detail, both analytically and on the lattice , and the results seem to support the instanton picture.
In the present work we wish to extend these studies to theories with fermions in the adjoint representation. Since we are dealing with a strongly coupled theory, our calculations are necessarily approximate. In particular, we will have to restrict ourselves to the contribution of small instantons for which the semi-classical description is appropriate. On the other hand, the methods we are using are applicable also to non-supersymmetric theories with several flavors of adjoint fermions. In addition to that, we can use these methods to study non-constant correlation functions that determine the spectrum of the theory.
The paper is organized as follows. In section II we discuss some general aspects of chiral symmetry breaking in theories with fermions in the adjoint representation. In section III we describe the structure of the instanton zero mode wave functions and in IV we calculate matrix elements of the Dirac operator between zero mode states. These results are used in order to determine the fermion determinant in the field of an instanton-anti-instanton pair (section V) and to calculate the gluino condensate in a random instanton ensemble (section VI). In section VII we describe simulations of an interacting instanton ensemble with different numbers of fermions in the fundamental and adjoint representation.
## II QCD with adjoint fermions
QCD with adjoint fermions is defined by the lagrangian
$$=\underset{i=1}{\overset{N_{ad}}{}}\frac{1}{2}\overline{\lambda }_M^{(i)a}(iD/)^{ab}\lambda _M^{(i)b}\frac{1}{4g^2}G_{\mu \nu }^aG_{\mu \nu }^a,$$
(3)
where $`\lambda _M^a`$ is a Majorana fermion in the adjoint representation of the gauge group and $`G_{\mu \nu }^a`$ is the usual field strength tensor. The covariant derivative in the adjoint representation is given by
$$D_\mu ^{ab}=_\mu \delta ^{ab}+f^{abc}A_\mu ^c.$$
(4)
For several Majorana flavors the theory (3) possesses a $`SU(N_{ad})`$ chiral symmetry. A non-zero gluino condensate
$$\overline{\lambda }_M^{(i)}\lambda _M^{(j)}=\delta ^{ij}\sigma $$
(5)
breaks this symmetry to $`SO(N_{ad})`$ . This fact can be seen most easily by considering the conserved vector and axial-vector currents . There are $`\frac{1}{2}N_{ad}(N_{ad}1)`$ conserved vector currents $`V_\mu ^{ij}=\overline{\lambda }_M^{(i)}\gamma _\mu \lambda _M^{(j)}`$ and $`\frac{1}{2}N_{ad}(N_{ad}+1)`$ classically conserved axial-vector currents $`A_\mu ^{ij}=\overline{\lambda }_M^{(i)}\gamma _\mu \gamma _5\lambda _M^{(j)}`$. The singlet axial current $`A_\mu ^{ii}`$ is anomalous. At the quantum level this leaves $`N_{ad}^21`$ conserved charges that generate the $`SU(N_{ad})`$ chiral symmetry. Gluino condensation breaks the axial symmetries and leads to the appearance of $`\frac{1}{2}(N_{ad}^2+N_{ad}2)`$ Goldstone bosons. The unbroken $`\frac{1}{2}N_{ad}(N_{ad}1)`$ vector charges generate the residual $`O(N_{ad})`$ symmetry.
In the case of supersymmetric gluodynamics, $`N_{ad}=1`$, there is no continuous symmetry. Instantons break the axial $`U(1)_A`$ symmetry but leave a discrete $`Z_{N_c}`$ symmetry intact. This discrete symmetry is spontaneously broken by gluino condensation. As discussed above, the value of the gluino condensate is known from a weak coupling instanton calculation. There are no predictions for the spectrum of the theory, but we expect the lowest states to fill out a chiral supermultiplet containing a scalar and a pseudoscalar meson as well as a Majorana fermion. These results can be summarized in terms of an effective lagrangian . This is not an effective lagrangian in the Wilsonean sense. The effective action does not generate the low momentum scattering amplitudes of the theory. Instead, it mainly serves as a generating functional for the anomalous Ward identities of the theory.
## III Instanton gauge potential and fermionic zero modes
In theories with adjoint fermions it is convenient to employ a spinor notation for spin, vector, and color indices . We can convert vectors to spinors using
$$V_{\alpha \dot{\alpha }}=V_\mu (\sigma _\mu )_{\alpha \dot{\alpha }}.$$
(6)
The euclidean spinor conventions used in this paper are summarized in Appendix A. The instanton gauge potential couples spin to color degrees of freedom. A field $`A^a`$ in the adjoint representation of $`SU(2)`$ can be represented by a symmetric tensor $`A^{\alpha \beta }`$
$$A^a=A^{\alpha \beta }ϵ_{\alpha \gamma }(\tau ^a)_\beta ^\gamma .$$
(7)
In spinor notation, the instanton gauge potential in regular gauge is given by
$$A_{\gamma \dot{\delta }}^{\alpha \beta }=2i(\delta _\gamma ^\alpha x_{\dot{\delta }}^\beta +\delta _\gamma ^\beta x_{\dot{\delta }}^\alpha )\frac{1}{x^2+\rho ^2}.$$
(8)
We can transform the gauge potential to singular gauge using the gauge transformation
$$U^{\dot{\alpha }\alpha }=\widehat{x}_\mu (\overline{\sigma }_\mu )^{\dot{\alpha }\alpha }.$$
(9)
Note that this matrix transforms an undotted color index into a dotted one. We can perform a ‘fake’ conversion of the dotted spinor back to an undotted one using the fact that $`(\sigma _0)^{\alpha \dot{\alpha }}`$ is just the unit matrix.
In the case of $`SU(2)`$, the Dirac operator in the background field of an instanton has four zero modes. The first two are conventionally called the supersymmetric (ss) zero modes
$$\lambda _{\alpha (\beta )}^{\gamma \delta }=(\delta _\alpha ^\gamma \delta _\beta ^\delta +\delta _\beta ^\gamma \delta _\alpha ^\delta )\frac{\rho ^2}{\pi }\frac{1}{(x^2+\rho ^2)^2},$$
(10)
where $`\beta =1,2`$ enumerates the zero modes. The other two are referred to as the superconformal (sc) zero modes
$$\lambda _{\alpha (\dot{\beta })}^{\gamma \delta }=(\delta _\alpha ^\gamma x_{\dot{\beta }}^\delta +\delta _\alpha ^\delta x_{\dot{\beta }}^\gamma )\frac{\rho }{\sqrt{2}\pi }\frac{1}{(x^2+\rho ^2)^2}.$$
(11)
In singular gauge, the zero modes are given by
$`\lambda _{\alpha (\beta )}^{\dot{\gamma }\dot{\delta }}=(x_\alpha ^{\dot{\gamma }}x_\beta ^{\dot{\delta }}+x_\beta ^{\dot{\gamma }}x_\alpha ^{\dot{\delta }}){\displaystyle \frac{\rho ^2}{\pi }}{\displaystyle \frac{1}{x^2(x^2+\rho ^2)^2}}(ss),`$ (12)
$`\lambda _{\alpha (\dot{\beta })}^{\dot{\gamma }\dot{\delta }}=(x_\alpha ^{\dot{\gamma }}\delta _{\dot{\beta }}^{\dot{\delta }}+x_\alpha ^{\dot{\delta }}\delta _{\dot{\beta }}^{\dot{\gamma }}){\displaystyle \frac{\rho }{\sqrt{2}\pi }}{\displaystyle \frac{1}{(x^2+\rho ^2)^2}}(sc).`$ (13)
Analogously, we can construct the zero modes of the Dirac operator in the background field of an anti-instanton. The regular gauge supersymmetric zero mode has the structure $`\lambda _{\dot{\gamma }\dot{\delta }}^{\dot{\alpha }(\dot{\beta })}(\delta _{\dot{\gamma }}^{\dot{\alpha }}\delta _{\dot{\delta }}^{\dot{\beta }}+\delta _{\dot{\delta }}^{\dot{\alpha }}\delta _{\dot{\gamma }}^{\dot{\beta }})`$, etc.
The effect of the zero modes on the propagation of fermions can be summarized in terms of an effective lagrangian . The ’t Hooft effective interaction in the case of one Majorana fermion in the adjoint representation of $`SU(2)`$ was determined in . The result is
$``$ $`=`$ $`{\displaystyle \frac{4\pi ^4}{3}}\left({\displaystyle \frac{2\pi }{\alpha _s}}\right)^4\mathrm{exp}({\displaystyle \frac{2\pi }{\alpha _s}})\rho ^3d\rho \{\overline{\lambda }_M^a\lambda _M^a_\mu \overline{\lambda }_M^b^\mu \lambda _M^b+\overline{\lambda }_M^a\gamma _5\lambda _M^a_\mu \overline{\lambda }_M^b\gamma _5^\mu \lambda _M^b`$ (15)
$`{\displaystyle \frac{1}{2}}\overline{\lambda }_M^a\sigma _{\alpha \beta }\lambda _M^b_\mu \overline{\lambda }_M^b\sigma ^{\alpha \beta }^\mu \lambda _M^a\}.`$
This result has to be interpreted with some care. The notion of an effective interaction induced by instantons of some fixed size is incompatible with supersymmetry. In order to derive manifestly supersymmetric results we always have to integrate over the collective coordinates of the instanton. Nevertheless, it is instructive to compare the result (15) with the effective interaction in the case of $`N_f=2`$ Dirac fermions in the fundamental representation. The structure of the two interactions is quite similar, which suggests that instantons may lead to similar physical effects. The most important difference between the two effective lagrangians is the presence of derivatives acting on two of the four Majorana spinors in (15). This difference is connected with the asymptotic behavior of the supersymmetric zero modes, which is not $`1/z^3`$, but $`1/z^4`$.
## IV Matrix elements of the Dirac operator
In the following, we wish to study the spectrum of the Dirac operator in an instanton ensemble. For this purpose, we have to calculate matrix elements of the Dirac operator between the zero modes of individual instantons and anti-instantons
$$T_{IA}=d^4x\overline{\lambda }_I^a(iD/)^{ab}\lambda _A^b.$$
(16)
An ensemble of instantons and anti-instantons is only an approximate saddle point of the action. If the system is sufficiently dilute then the instantons and anti-instantons are well separated and the approximate saddle point solution for the gauge potential is given by a simple sum of the gauge potentials of the individual instantons. For this purpose, the gauge potential of the individual instantons has to be put in singular gauge. In the sum ansatz, we can use the equations of motion of the fermion fields in order to replace the covariant derivative in (16) by an ordinary derivative
$$T_{IA}=d^4x\overline{\lambda }_I^a(i/)\lambda _A^a.$$
(17)
The structure of the Dirac operator is dictated by the form of the zero modes. In the background field of an instanton-anti-instanton pair we have
$$T_{IA}=\left(\begin{array}{cc}0& \begin{array}{cc}T_{IA}^{ssss}& T_{IA}^{sssc}\\ T_{IA}^{scss}& T_{IA}^{scsc}\end{array}\\ \begin{array}{cc}T_{AI}^{ssss}& T_{AI}^{sssc}\\ T_{AI}^{scss}& T_{AI}^{scsc}\end{array}& 0\end{array}\right),$$
(18)
where the matrix elements $`T_{AI}^{ss},\mathrm{}`$ are real quaternions. These quaternions can be decomposed as
$`(T_{AI}^{ssss})_{\dot{\beta }\beta ^{}}`$ $`=`$ $`T_\mu ^{ss}(\sigma _\mu )_{\beta ^{}\dot{\beta }}`$ (19)
$`(T_{AI}^{scsc})_{\beta \dot{\beta }^{}}`$ $`=`$ $`T_\mu ^{sc}(\sigma _\mu )_{\beta \dot{\beta }^{}}`$ (20)
$`(T_{AI}^{sssc})_{\dot{\beta }\dot{\beta }^{}}`$ $`=`$ $`T^{sssc}ϵ_{\dot{\beta }\dot{\beta }^{}}+T_{\mu \nu }^{sssc}ϵ_{\dot{\beta }\dot{\gamma }}(\overline{\sigma }_{\mu \nu })_{\dot{\beta }^{}}^{\dot{\gamma }},`$ (21)
$`(T_{AI}^{scss})_{\beta \beta ^{}}`$ $`=`$ $`T^{scss}ϵ_{\beta \beta ^{}}+T_{\mu \nu }^{scss}(\sigma _{\mu \nu })_\beta ^\gamma ϵ_{\gamma \beta ^{}}.`$ (22)
Here, $`T_\mu ^{ss}`$ and $`T_\mu ^{sc}`$ are real vectors, $`T^{sssc}`$ and $`T^{scss}`$ are real scalars and $`T_{\mu \nu }^{sssc}`$ and $`T_{\mu \nu }^{scss}`$ are self-dual and anti-self-dual tensors, respectively. Chiral symmetry implies that the diagonal blocks of $`T_{IA}`$ are zero. The upper right and lower left blocks are related by hermitean conjugation. For example, we find that
$$(T_{IA}^{ssss})_{\beta \dot{\beta }^{}}=T_\mu ^{ss}(\overline{\sigma }^\mu )^{\dot{\beta }^{}\beta }.$$
(23)
In general, we have $`(T^{})_{AI}=(T)_{IA}`$. The eigenvalues of (18) come in quartets $`(\xi ,\xi ,\xi ,\xi )`$. These results are in agreement with the general arguments presented in .
The functions $`T_\mu ^{ss},\mathrm{}`$ depend on the collective coordinates of the instanton and anti-instanton. We will characterize the relative color orientation by the four vector $`u_\mu =1/2\mathrm{tr}(U_A^{}U_I\sigma _\mu )`$. Here, $`U_{I,A}`$ are $`SU(N_c)`$ matrices that characterize the color orientation of the instanton and anti-instanton. For color $`SU(2)`$ $`u_\mu `$ is a real vector with $`u^2=1`$. Using rotational symmetry and the fact that $`T_{AI}`$ is quadratic in $`u_\mu `$ we have
$`T_\mu ^{ss}`$ $`=`$ $`\widehat{z}_\mu T_1^{ss}+u_\mu (u\widehat{z})T_2^{ss}+\widehat{z}_\mu (u\widehat{z})^2T_3^{ss},`$ (24)
$`T_\mu ^{sc}`$ $`=`$ $`\widehat{z}_\mu T_1^{sc}+u_\mu (u\widehat{z})T_2^{sc}+\widehat{z}_\mu (u\widehat{z})^2T_3^{sc},`$ (25)
$`T^{sssc}`$ $`=`$ $`T_1^{sssc}+T_2^{sssc}(u\widehat{z})^2,`$ (26)
$`T_{\mu \nu }^{sssc}`$ $`=`$ $`(u_\mu \widehat{z}_\nu u_\nu \widehat{z}_\mu )(u\widehat{z})T_3^{sssc},`$ (27)
where $`z_\mu =z_\mu ^Az_\mu ^I`$ and the functions $`T_1^{ss},\mathrm{}`$ depend on $`(|z_\mu |,\rho _I,\rho _A)`$. For simplicity, we will assume that the dependence on $`\rho _{I,A}`$ only enters through their geometric mean $`\overline{\rho }=\sqrt{\rho _I\rho _A}`$. The fact that this assumption is valid to fairly good accuracy was checked in the case of fundamental fermions. In a more sophisticated treatment of the instanton-anti-instanton gauge configuration the dependence of the overlap matrix element on $`(z,\rho _I,\rho _A)`$ is restricted by conformal invariance .
In the following we shall outline the calculation of the invariant functions $`T_1^{ss},\mathrm{}`$. We describe the case $`T_1^{ss}`$ in some detail but relegate the results for the other functions to appendix B. Using the expression (12) for the wave function of the supersymmetric zero modes in singular gauge we find
$`T_\eta ^{ss}`$ $`=`$ $`\left\{\mathrm{tr}\left(\overline{\sigma }_\mu \sigma _\rho \overline{\sigma }_\beta \sigma _\alpha \right)\mathrm{tr}\left(\overline{\sigma }_\nu \sigma _\sigma \overline{\sigma }_\gamma \sigma _\eta \right)+\mathrm{tr}\left(\overline{\sigma }_\mu \sigma _\rho \overline{\sigma }_\beta \sigma _\eta \overline{\sigma }_\nu \sigma _\sigma \overline{\sigma }_\gamma \sigma _\alpha \right)\right\}`$ (29)
$`u_\rho u_\sigma {\displaystyle d^4x\varphi _{\mu \nu }(xz)_\alpha \varphi _{\beta \gamma }(x)},`$
where $`\varphi _{\mu \nu }(x)`$ is the profile function of the supersymmetric zero mode
$$\varphi _{\mu \nu }(x)=\frac{\rho ^2}{\pi }\frac{\widehat{x}_\mu \widehat{x}_\nu }{(x^2+\rho ^2)^2}.$$
(30)
The integral in (29) is most easily calculated in Fourier space. The Fourier transform of $`\varphi _{\mu \nu }`$ is given by
$$\varphi _{\mu \nu }(k)=\delta _{\mu \nu }\varphi _1(k)+\widehat{k}_\mu \widehat{k}_\nu \varphi _2(k)$$
(31)
with
$`\varphi _1(k)`$ $`=`$ $`{\displaystyle \frac{2\pi \rho ^2}{y}}\left\{{\displaystyle \frac{4}{y^3}}\left({\displaystyle \frac{4}{y^2}}+1\right)K_1(y){\displaystyle \frac{2}{y}}K_0(y)\right\},`$ (32)
$`\varphi _2(k)`$ $`=`$ $`2\pi \rho ^2\left\{{\displaystyle \frac{16}{y^4}}\left({\displaystyle \frac{16}{y^3}}+{\displaystyle \frac{4}{y}}\right)K_1(y)\left({\displaystyle \frac{8}{y^2}}+1\right)K_0(y)\right\},`$ (33)
and $`y=k\rho `$. $`K_n(y)`$ is the modified Bessel function of the first kind and order $`n`$. We can now calculate the overlap integral and perform the traces. In momentum space the result is given by
$$T_\eta ^{ss}(k)=(i)\left(2\widehat{k}_\eta 8u_\eta (u\widehat{k})+16(u\widehat{k})^2\right)|\varphi _2(k)|^2.$$
(34)
Finally, we can determine the functions $`T_{1,2,3}^{ss}`$ by performing the inverse Fourier transform. In the $`d^4k`$ integral all integrations except for the one over the absolute magnitude of $`k`$ can be performed analytically. We find
$`T_1^{ss}(z)`$ $`=`$ $`{\displaystyle \frac{1}{8\pi ^2}}{\displaystyle 𝑑k\left\{2k^4j_1(kz)16\frac{k^3}{z}j_2(kz)\right\}|\varphi _2(k)|^2}`$ (35)
$`T_2^{ss}(z)`$ $`=`$ $`{\displaystyle \frac{1}{8\pi ^2}}{\displaystyle 𝑑k\left\{8k^4j_1(kz)32\frac{k^3}{z}j_2(kz)\right\}|\varphi _2(k)|^2}`$ (36)
$`T_3^{ss}(z)`$ $`=`$ $`{\displaystyle \frac{1}{8\pi ^2}}{\displaystyle 𝑑k\mathrm{\hspace{0.17em}\hspace{0.33em}16}k^4j_3(kz)|\varphi _2(k)|^2},`$ (37)
where $`j_n(x)`$ is the spherical Bessel function of order $`n`$. The integrals (35) have to be performed numerically. The results are shown in Figure 1. In the following, we will use a simple parameterization of the numerical results. In the case of the supersymmetric overlaps, we use
$`\overline{\rho }T_1^{ss}(z)`$ $`=`$ $`{\displaystyle \frac{1.26\overline{z}}{1.0+2.34\overline{z}^2+0.35\overline{z}^4+0.24\overline{z}^6}},`$ (38)
$`\overline{\rho }T_2^{ss}(z)`$ $`=`$ $`{\displaystyle \frac{1.05\overline{z}}{(1.0+0.38\overline{z}^2)^3}}+{\displaystyle \frac{6.36\overline{z}^3}{(1.0+0.68\overline{z}^2)^4}},`$ (39)
$`\overline{\rho }T_3^{ss}(z)`$ $`=`$ $`{\displaystyle \frac{15.8\overline{z}^3}{(1.0+0.84\overline{z}^2)^4}},`$ (40)
where $`\overline{z}=z/\overline{\rho }`$. These parameterizations respect the asymptotic behavior of the overlap integrals. In particular, we have $`T^{ss}(z)1/z^5`$, $`T^{sc}(z)1/z^3`$ and $`T^{sssc}(z)1/z^4`$.
For completeness, let us compare these results to the corresponding expressions in the case of fundamental fermions. In this case, there is only one zero mode per instanton. The overlap matrix element $`T_{IA}`$ is a real number in the case of $`SU(2)`$, and complex for $`SU(N_c>2)`$. $`T_{IA}`$ satisfies the symmetry relation $`T_{IA}=T_{AI}^{}`$. As a consequence, the eigenvalues are real and occur in pairs $`(\xi ,\xi )`$. We can extract the dependence of $`T_{IA}`$ on the collective coordinates. The result is
$$T_{AI}^{fund}=(u\widehat{z})T^f(z,\rho _I,\rho _A),T^f(z,\rho _I,\rho _A)\frac{1}{\rho _I\rho _A}\frac{4z}{(2.0+z^2/(\rho _I\rho _A))^2}.$$
(41)
We note that the fundamental overlap matrix element only depends on one $`SU(2)`$ angle $`\mathrm{cos}\theta (u\widehat{z})`$. From the asymptotic form of the zero mode solution one finds $`T^f(z)1/z^3`$.
## V The fermion determinant in the field of an instanton-anti-instanton pair
Before we study gluino condensation we would like to make a brief digression and discuss the gluino induced instanton-anti-instanton interaction. This interaction will play an important role in the calculation of the gluino condensate in an interacting instanton ensemble.
The probability to find an instanton-anti-instanton pair characterized by the collective coordinates $`(z_{I,A},\rho _{I,A},U_{I,A})`$ is controlled by the weight factor $`\mathrm{exp}(S)det(D/+m)`$. The first factor is the well known gluonic interaction. If the instanton and anti-instanton are well separated it has the dipole form
$$S=2S_0S_0\frac{4\rho _I^2\rho _A^2}{z^4}\left(14\mathrm{cos}^2\theta \right),$$
(42)
where $`S_0=(8\pi ^2)/g^2`$ is the single instanton action and $`\mathrm{cos}\theta `$ is the $`SU(2)`$ angle introduced above. We note that the interaction is attractive if the color orientation is aligned with the spatial orientation, $`\mathrm{cos}\theta =\pm 1`$. The second factor is the fermion determinant. In the case of fundamental fermions, it is also well known. We have
$$det(D/)=\mathrm{cos}^2\theta \frac{16}{\rho _I^2\rho _A^2}\frac{z^2}{\left(2.0+z^2/(\rho _I\rho _A)\right)^4},$$
(43)
which is also attractive for $`\mathrm{cos}\theta =\pm 1`$. We also note that the interaction peaks at $`z^2\rho _I\rho _A`$.
Using the results of the last section we can calculate the fermion induced interaction with fermions in the adjoint representation. In order to calculate the determinant for one Majorana fermion we take the square root of the corresponding expression for a Dirac fermion in the adjoint representation. For simplicity, let us begin with the determinant in the basis of the supersymmetric zero modes only. We find
$`det(D/)^{ss}`$ $`=`$ $`|(T_1^{ss})^2+((T_2^{ss})^2+2T_1^{ss}T_2^{ss}+2T_1^{ss}T_3^{ss})\mathrm{cos}^2\theta `$ (45)
$`+((T_3^{ss})^2+2T_2^{ss}T_3^{ss})\mathrm{cos}^4\theta |.`$
The result for the superconformal zero modes is even more simple,
$$det(D/)^{sc}=|(T_1^{sc})^2+((T_2^{sc})^2+2T_1^{sc}T_2^{sc})\mathrm{cos}^2\theta |.$$
(46)
This expression is quite similar to the determinant for fundamental fermions. The supersymmetric determinant (45) is somewhat more complicated, but also peaked for $`\mathrm{cos}\theta =\pm 1`$. When the mixing between supersymmetric and superconformal zero modes is included the fermion determinant depends on other $`SU(2)`$ angles in addition to $`\mathrm{cos}\theta `$. We show numerical results for $`\mathrm{log}(det(D/))`$ as a function of $`z`$, $`\mathrm{cos}\theta `$ and $`\mathrm{cos}\varphi `$ in Fig.2. Here we have taken $`\widehat{z}_\mu =z\delta _{\mu 4}`$ and defined $`\mathrm{cos}\theta =u_4`$ and $`\mathrm{sin}\theta \mathrm{cos}\varphi =u_2`$. We observe that again the determinant peaks for $`z^2\rho _I\rho _A`$. For large $`z`$, the determinant behaves as $`1/z^{16}`$. More importantly, we find that the interaction is again most attractive for $`\mathrm{cos}\theta =\pm 1`$. There is some dependence on $`\mathrm{cos}\varphi `$, but it is not as pronounced as the dependence on $`\mathrm{cos}\theta `$. This means that the gluino induced interaction for one Majorana fermion is qualitatively similar to the quark induced interaction with an effective number of quark flavors between $`N_f=2`$ (which gives $`det1/z^{12}`$) and $`N_f=3`$ (corresponding to $`det1/z^{18}`$).
## VI Gluino condensation in a random instanton ensemble
In this section we study gluino condensation in a random instanton ensemble. This means that we will assume that the collective coordinates of the instantons and anti-instanton are distributed randomly. In particular, we shall neglect the effect of the fermion determinant on the distribution of instantons. This is not a good approximation even in ordinary QCD and it certainly cannot be correct in a supersymmetric theory. Nevertheless, using the approximation of randomness we can get some analytic understanding of the dependence of the gluino condensate on the parameters characterizing the instanton liquid. We can also get an estimate of the relative size of the quark and gluino condensates in theories with both fundamental and adjoint fermions.
The simplest model of the spectrum of the Dirac operator is based on the assumption that the non-zero matrix elements of the Dirac operator are Gaussian random numbers . The distribution is characterized by the first moment
$$\sigma ^2=\frac{2}{N}\mathrm{tr}\left(T^{}T\right).$$
(47)
The eigenvalue distribution for the Gaussian ensemble is given by a semi-circle where the density of eigenvalues at zero virtuality is $`\rho (0)=(N/V)(\pi \sigma )^1`$. Here, $`(N/V)`$ is the number of eigenstates per unit volume. The first moment of the overlap matrix can be estimated by averaging $`|T_{AI}|^2`$ over the collective coordinates of the instantons. Using (41) we find the first moment of the Dirac operator for fermions in the fundamental representation of $`SU(2)`$
$$\sigma =\left(\frac{1}{3}\frac{N}{V}\right)^{1/2}\overline{\rho }\pi ,$$
(48)
where $`\overline{\rho }`$ is the average size of the instanton. This parameter, just like the density of instantons, cannot be determined in the semi-classical approximation. In the instanton liquid model of the QCD vacuum it is assumed that $`\overline{\rho }=1/3`$ fm and $`(N/V)=1\mathrm{fm}^4`$ . Using these values we find
$$\overline{q}q=\frac{1}{\pi \overline{\rho }}3^{1/2}\left(\frac{N}{V}\right)^{1/2}(230\mathrm{MeV})^3,$$
(49)
in very good agreement with the phenomenological value (which, of course, applies to color $`SU(3)`$).
The same arguments can be applied to gluino condensation in a random instanton ensemble. In this case we need to determine the first moment of a quaternionic matrix with the matrix elements determined in section IV. We find
$$\sigma _{ad}=\left(\frac{N}{V}\right)^{1/2}0.43\overline{\rho }\pi ,$$
(50)
which is somewhat smaller than (48). There are four times as many eigenstates per unit volume but for a Majorana fermion the Banks-Casher relation has an additional factor 1/2, $`\overline{\lambda }\lambda =\pi /2\rho (0)`$. We finally get the following estimate of the gluino condensate
$$\overline{\lambda }\lambda =\frac{1}{\pi \overline{\rho }}4.6\left(\frac{N}{V}\right)^{1/2}.$$
(51)
Here and in what follows we have dropped the subscript M on the Majorana spinors. This result is a little more than twice as large as the corresponding result for a Dirac fermion in the fundamental representation. As we saw, this is mainly due to the effective number of zero modes in both cases. We emphasize that the gluino condensate is proportional to the square root of the instanton density, which is also what one would expect if the condensate is extracted from the four-point function using clustering.
We have checked these estimates by performing a numerical calculation of the spectrum of the Dirac operator in a random instanton ensemble. This means that instead of assuming the matrix elements of the Dirac operator to be random we take the collective coordinates of the instantons and anti-instantons to be random. We calculate the spectrum of the Dirac operator and determine the gluino condensate using
$$\overline{\lambda }\lambda =\frac{1}{2}𝑑\lambda \rho (\lambda )\frac{2m}{\lambda ^2+m^2}.$$
(52)
The results are shown in Fig. 3. We observe that the spectrum is not a semi-circle but is peaked towards zero virtuality . This non-analyticity is smoothed out when we calculate the condensate for a non-zero quark or gluino mass. Again using $`(N/V)=1\mathrm{fm}^4`$ and $`\overline{\rho }=1/3\mathrm{fm}`$ as well as $`m_q=m_g=20`$ MeV we find $`\overline{\psi }\psi =(260\mathrm{MeV})^3`$ and $`\overline{\lambda }\lambda =(347\mathrm{MeV})^3`$.
## VII Gluino condensation in an unquenched instanton ensemble
As we stressed in the previous section, the assumption of randomness is not expected to be very useful. The fermion determinant is given by the product of all eigenvalues of the Dirac operator, while the quark or gluino condensate is determined by the density of small eigenvalues. This implies that the determinant tends to suppress fermion condensates. In particular, we expect that the strength of chiral symmetry breaking is reduced as the number of fermion flavors is increased.
In this section we shall study this problem using simulations of the instanton ensemble in QCD with fundamental and adjoint fermions. We consider the partition function
$$Z=\left(\underset{i}{\overset{N}{}}d\mathrm{\Omega }_id(\rho _i)\right)det(D/_f+m_q)^{N_f}det(D/_a+m_g)^{N_{ad}/2}\mathrm{exp}(S)$$
(53)
Here, $`\mathrm{\Omega }`$ denotes the collective coordinates of the instanton, $`d(\rho )`$ is the single instanton distribution , $`D/_{f,a}`$ are the Dirac operators in the fundamental and adjoint representation, and $`\mathrm{exp}(S)`$ is the gluonic interaction between instantons. In order to study spontaneous symmetry breaking in a finite volume we introduce non-zero quark and gluino masses $`m_{q,g}`$. We will study the limit $`m_{q,g}0`$ in some detail.
The partition function (53) suffers from the usual IR problem connected with large instantons for which the semi-classical approximation does not apply. In practice, we deal with this problem by introducing a short range repulsive core in the gluonic instanton interaction, see section V.C. in for a more detailed discussion. The repulsive core eliminates the contributions of large instantons and very close pairs. This particular method for suppressing objects that are not semi-classical has the virtue that it respects the classical scale invariance of Yang-Mills theory.
The instanton ensemble is characterized by two numbers, the scale parameter $`\mathrm{\Lambda }`$ that enters into the instanton weight $`d(\rho )`$ and a dimensionless parameter $`A`$ which determines the size of the core. Lacking a better theory of topological fluctuations beyond the semi-classical domain we have to fix $`A`$ phenomenologically. This could be done, for example, as soon as lattice information on the spectrum and other properties of theories with adjoint fermions becomes available . In this work we will use the same value that was employed in studies of QCD with fundamental fermions. It leads to a dilute instanton ensemble characterized by the dimensionless parameter $`\overline{\rho }^4(N/V)0.12`$. For simplicity we will concentrate on simulations at a fixed instanton density $`(N/V)=1.0\mathrm{\Lambda }^4`$.
To set the stage, we show results for $`N_f=1,\mathrm{},4`$ flavors of fundamental fermions. Fig. 4 shows the quark condensate as a function of the quark mass from simulations in a euclidean box of size $`V=2.0^4\mathrm{\Lambda }^4`$. The case of only one flavor is special. The chiral condensate persists even if the limit $`m_q0`$ is taken in a finite volume. This is due to the fact that for $`N_f=1`$ the quark condensate is dominated by direct instanton contributions. The result for $`N_f=2`$ is characteristic of spontaneous symmetry breaking. The quark condensate vanishes as the quark mass goes to zero but shows a clear plateau for larger quark masses. One can verify that the onset of chiral symmetry breaking moves towards smaller masses as the volume is increased. For more than two flavors the chiral condensate is significantly reduced. In the case of three flavors the signal is already quite weak. Using simulations in bigger volumes one can verify that chiral symmetry is indeed broken. There is no clear evidence for chiral symmetry breaking in simulations with four or more flavors.
Fig. 4b shows the gluino condensate measured in simulations with one or two flavors of Majorana fermions in the adjoint representation. For $`N_{ad}=1`$ there is clear evidence for spontaneous symmetry breaking. Indeed, the behavior is more reminiscent of the case $`N_f=1`$, where $`\overline{q}q`$ receives direct instanton contributions, than the case $`N_f=2`$, in which chiral symmetry breaking is a collective effect.
These observations can be understood in more physical terms. Supersymmetric gluodynamics has no Goldstone bosons, so finite volume effects are much weaker than in $`N_f=2`$ non-supersymmetric QCD. This means that in a fixed volume, gluino condensation can be observed for gluino masses that are significantly smaller than the quark masses required to produce quark condensation. In the standard picture, there is a discrete chiral symmetry which is broken by gluino condensation. This means that if the gluino mass is too small then chiral symmetry will be restored because of tunneling between the $`Z_2`$ vacua. This is different from $`N_f=1`$ non-supersymmetric QCD where instantons leave no unbroken discrete symmetries.
The value of the gluino condensate is $`\overline{\lambda }\lambda 2\mathrm{\Lambda }^3`$. This result has the correct order of magnitude but it cannot yet be compared directly to the prediction (1). First of all, we use a different definition of the scale parameter. In order to make contact with our work on QCD we use a Pauli-Vilars scale parameter. Second, we have an additional parameter $`A`$ which controls the boundary of the semi-classical regime. Finally, we have performed the simulations at a fixed density of instantons $`(N/V)=1.0\mathrm{\Lambda }^4`$. It is this choice which effectively sets the scale in our calculation.
In Fig. 4b we also show the gluino condensate measured in simulations with $`N_{ad}=2`$ Majorana flavors. The condensate is very small and there is no clear evidence for spontaneous chiral symmetry breaking.
The spectrum of the Dirac operator for $`N_f=2`$ quark flavors and $`N_{ad}=1`$ Majorana flavor is shown in Fig. 5a and b. The spectra were determined in simulations with $`m_{q,g}=0.1\mathrm{\Lambda }^1`$. Again, we observe that in both cases there is a finite density of eigenvalues as $`\lambda 0`$. For $`N_f=2`$ the spectral density near $`\lambda =0`$ is flat<sup>*</sup><sup>*</sup>* Fig. 5 shows that the spectral density is flat for intermediate values of $`\lambda `$. There is a finite volume suppression of the spectral density for small $`\lambda `$ and a $`O(m^2)`$ peak at $`\lambda =0`$. To show that the spectral density is flat at $`\lambda =0`$ in the limit $`V\mathrm{},m0`$ requires more numerical work., whereas in the case $`N_{ad}=1`$ it is growing towards small $`\lambda `$. Again, this is similar to the case of only one fundamental fermion. The results are consistent with the effective field theory prediction
$$\rho ^{}(\lambda =0)=\frac{\mathrm{\Sigma }_0^2}{16\pi ^2f_\pi ^4}\frac{(N_f2)(N_f+\beta )}{\beta N_f}.$$
(54)
Here, $`\beta `$ is the Dyson index of the random matrix ensemble with the appropriate symmetry. We have $`\beta =1`$ for fundamental fermions in $`SU(2)`$, $`\beta =2`$ for fundamental fermions in $`SU(N>2)`$, and $`\beta =4`$ in the case of fermions in the adjoint representation. $`N_f`$ denotes the number of Dirac or Majorana flavors in the cases $`\beta =1,2`$ and $`\beta =4`$, respectively. $`\mathrm{\Sigma }_0`$ is the magnitude of the quark condensate and $`f_\pi `$ the pion decay constant. The expression (54) summarizes the fact that the spectrum is peaked towards small virtuality for both $`N_f=1`$ and $`N_{ad}=1`$ while it is flat for $`N_f=2`$. Effective field theory predicts the slope of the Dirac spectrum under the assumption that chiral symmetry is broken. The theory cannot predict whether chiral symmetry breaking takes place for some given $`N_f`$ or $`N_{ad}`$.
## VIII Conclusions
In summary we have studied gluino condensation and the spectrum of the Dirac operator in an instanton ensemble. We employ the semi-classical approximation and focus on the Dirac operator in the subspace spanned by the zero modes of the individual instantons and anti-instantons. We have shown how the quaternionic structure of the Dirac operator in theories with adjoint fermions emerges naturally from the spin and color structure of the zero modes. The dependence of the matrix elements on the collective coordinates of the instantons is quite complicated but qualitatively similar to the simpler case of fundamental fermions.
We have provided evidence that gluino condensation does take place in an ensemble of instanton and anti-instantons. In a random ensemble, the gluino condensate is proportional to the square root of the instanton density. In supersymmetric gluodynamics we find that gluino condensation persists even if interactions between the instantons are taken into account. We observed that finite volume effects are much weaker than in QCD with two flavors of fundamental fermions. This is consistent with the fact that supersymmetric gluodynamics has a large mass gap. In QCD with more than one adjoint flavor we found no compelling evidence for gluino condensation.
There are many problems that remain to be studied. In particular, it would be interesting to make a systematic study of gluino and gluino-glueball correlation functions. There are two types of correlation functions: Constant correlators that provide information on condensates, and $`x`$-dependent correlators related to the spectrum. These correlation functions will also show to what extent supersymmetry is realized in the limit $`m_g0`$. In addition to that, it would be interesting to search for evidence of $`Z_2`$ domains and to investigate the dependence of the results on the topological sector of the theory. In this work we have used the zero mode wave functions that correspond to trivial holonomy and anti-periodic boundary conditions on the fermions. This suggests the question of how the results are changed if the boundary conditions are modified. In this case, the zero modes discussed in will come into play. Finally, it is important to study the role of very large instantons that were excluded in the present study.
There have been suggestions that objects with fractional topological charge may play a role in theories with adjoint fermions . These objects can give a direct contribution to the gluino condensate. Because of tunneling between the different $`Z_N`$ phases the presence of such objects cannot be inferred from the behavior of the gluino condensate as a function of the quark mass in a finite volume. One should be able, however, to detect the presence of fractionally charged objects in lattice simulations by looking for zero modes of the Dirac operator that do not appear in multiples of $`2N_c`$ . In this context it would also be interesting to study gluino condensation for $`N_c>2`$. For adjoint fermions the number of zero modes per topological charge increases with $`N_c`$. One might therefore doubt that instantons alone are sufficient to trigger gluino condensation in large $`N_c`$ SUSY gluodynamics. It has also been suggested that fractionally charged objects can be thought of as instanton constituents . One might then envision a situation where if $`N_c`$ is small, or instantons are small, fractionally charged objects are bound into instantons while for large $`N_c`$, or for large instantons, topological objects dissociate and the instanton liquid should be replaced by liquid of fractional charges.
Acknowledgments: I would like to thank E. Shuryak, J. Verbaarschot and A. Zhitnitsky for useful discussions. I would also like to thank M. Shifman for many valuable comments and pointing out some errors in an earlier version of this manuscript. This work was supported in part by the US DOE grant DE-FG-88ER40388.
## A Euclidean Spinor Conventions
We use the following euclidean spinor conventions
$`\gamma _\mu `$ $`=\left(\begin{array}{cc}\hfill 0& \hfill \sigma _\mu \\ \hfill \overline{\sigma }_\mu & \hfill 0\end{array}\right)=\gamma _\mu ^{},\gamma _5=\left(\begin{array}{cc}\hfill 1& \hfill 0\\ \hfill 0& \hfill 1\end{array}\right),`$ (A5)
$`\sigma _\mu `$ $`=(i\stackrel{}{\sigma },1),\overline{\sigma }_\mu =(i\stackrel{}{\sigma },1),`$ (A6)
$`(\sigma _{\mu \nu })_\alpha ^\beta `$ $`={\displaystyle \frac{1}{4}}\left[(\sigma _\mu )_{\alpha \dot{\alpha }}(\overline{\sigma }_\nu )^{\dot{\alpha }\beta }(\sigma _\nu )_{\alpha \dot{\alpha }}(\overline{\sigma }_\mu )^{\dot{\alpha }\beta }\right],`$ (A7)
$`(\overline{\sigma }_{\mu \nu })_{\dot{\beta }}^{\dot{\alpha }}`$ $`={\displaystyle \frac{1}{4}}\left[(\overline{\sigma }_\mu )^{\dot{\alpha }\alpha }(\sigma _\nu )_{\alpha \dot{\beta }}(\overline{\sigma }_\nu )^{\dot{\alpha }\alpha }(\sigma _\mu )_{\alpha \dot{\beta }}\right].`$ (A8)
Indices are raised and lowered with $`ϵ^{\alpha \beta }`$ and $`ϵ^{\dot{\alpha }\dot{\beta }}`$ where $`ϵ^{\alpha \beta }ϵ_{\beta \gamma }=\delta _\gamma ^\alpha `$ and $`ϵ^{\dot{\alpha }\dot{\beta }}=ϵ^{\alpha \beta }`$. The euclidean sigma matrices have the following properties
$`(\sigma _\mu \overline{\sigma }_\nu )_\alpha ^\beta `$ $`=`$ $`\delta _{\mu \nu }\delta _\alpha ^\beta +2(\sigma _{\mu \nu })_\alpha ^\beta ,`$ (A9)
$`(\overline{\sigma }_\mu \sigma _\nu )_{\dot{\beta }}^{\dot{\alpha }}`$ $`=`$ $`\delta _{\mu \nu }\delta _{\dot{\alpha }}^{\dot{\beta }}+2(\overline{\sigma }_{\mu \nu })_{\dot{\beta }}^{\dot{\alpha }},`$ (A10)
$`(\overline{\sigma }_\mu )^{\dot{\alpha }\alpha }`$ $`=`$ $`ϵ^{\alpha \beta }ϵ^{\dot{\alpha }\dot{\beta }}(\sigma _\mu )_{\beta \dot{\beta }}`$ (A11)
$`\sigma _{\mu \nu }`$ $`=`$ $`{\displaystyle \frac{1}{2}}ϵ_{\mu \nu \rho \sigma }\sigma _{\rho \sigma },\overline{\sigma }_{\mu \nu }={\displaystyle \frac{1}{2}}ϵ_{\mu \nu \rho \sigma }\overline{\sigma }_{\rho \sigma }.`$ (A12)
## B Matrix elements
In this appendix we collect the remaining matrix elements of the Dirac operator. We define the profile function of the superconformal zero mode
$$\varphi _\mu =\frac{\rho }{\sqrt{2}\pi }\frac{x_\mu }{(x^2+\rho ^2)^2}$$
(B1)
and its Fourier transform $`\varphi _\mu (k)=i\widehat{k}_\mu \varphi _3(k)`$ with
$$\varphi _3(k)=\sqrt{2}\pi \rho ^2K_1(k\rho ).$$
(B2)
The matrix elements of the Dirac operator between superconformal zero modes are determined by
$$T_\eta ^{sc}(k)=(+i)\left(2\widehat{k}_\eta +8u_\eta (u\widehat{k})^2\right)k|\varphi _3(k)|^2,$$
(B3)
and the matrix elements between supersymmetric and superconformal zero modes lead to
$`T^{sssc}(k)`$ $`=`$ $`\left(28(u\widehat{k})^2\right)k\varphi _2(k)\varphi _3(k),`$ (B4)
$`T_{\mu \nu }^{sssc}(k)`$ $`=`$ $`8\left(u_\mu \widehat{k}_\nu u_\nu \widehat{k}_\mu \right)(u\widehat{k})k\varphi _2(k)\varphi _3(k).`$ (B5)
From these results we can extract the invariant functions
$$T_1^{sc}(z)=\frac{1}{8\pi ^2}𝑑k\mathrm{\hspace{0.17em}2}k^4j_1(kz)|\varphi _3(k)|^2$$
(B6)
and $`T_2^{sc}(z)=4T_1^{sc}(z)`$ as well as $`T_3^{sc}(z)=0`$. Also
$`T_1^{sssc}(z)`$ $`=`$ $`{\displaystyle \frac{1}{8\pi ^2}}{\displaystyle 𝑑k\left[2k^4j_0(kz)8\frac{k^3}{z}j_1(kz)\right]\varphi _2(k)\varphi _3(k)},`$ (B7)
$`T_2^{sssc}(z)`$ $`=`$ $`{\displaystyle \frac{1}{8\pi ^2}}{\displaystyle 𝑑k\mathrm{\hspace{0.17em}8}k^4j_2(kz)\varphi _2(k)\varphi _3(k)},`$ (B8)
and $`T_3^{sssc}(z)=T_2^{sssc}(z)`$. Numerical results for these functions are shown in Fig. 1. The results can be parametrized as
$$\overline{\rho }T_1^{sc}(z)=\frac{0.25\overline{z}}{1.0+0.42\overline{z}^2+0.21\overline{z}^4}$$
(B9)
as well as
$`\overline{\rho }T_1^{sssc}(z)`$ $`=`$ $`{\displaystyle \frac{0.17}{1.0+0.05\overline{z}^2+0.08\overline{z}^4}},`$ (B10)
$`\overline{\rho }T_2^{sssc}(z)`$ $`=`$ $`{\displaystyle \frac{1.2\overline{z}^2}{(1.0+0.45\overline{z}^2)^3}}+{\displaystyle \frac{0.014\overline{z}^2}{(1.0+0.21\overline{z}^2)^3}},`$ (B11)
where $`\overline{z}=z/\overline{\rho }`$. The overlap matrix elements $`T^{sssc}`$ are related to the corresponding functions with the supersymmetric and superconformal zero modes interchanged. We find $`T_{1,2}^{scss}=T_{1,2}^{sssc}`$ and $`T_3^{scss}=T_3^{sssc}`$.
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# A new cosmological tracker solution for Quintessence
## I Acknowledgements
This work was partly supported by CONACyT, México 119259 (L.A.U.)
## II captions
FIG. 1. Evolution of the dimensionless density parameters taking $`\mathrm{\Omega }_{oM}=0.30`$ and $`\mathrm{\Omega }_{o\mathrm{\Phi }}=0.70`$. The different models with $`\omega _\mathrm{\Phi }=\{0.6(red),0.7(green),0.8(blue),0.9(pink)\}`$ are also shown.
FIG. 2. The luminosity distance vs redshift for the models shown in fig. I. The horizontal line is the luminosity distance for $`\mathrm{\Lambda }CDM`$. The open circles represent the observational results from SCP and the opaque circles are from HZS.
FIG. 3. Angular power spectrum for the models considered in fig. I and $`\mathrm{\Lambda }CDM`$ (black). The data were taken from.
FIG. 4. Mass power spectrum for the models shown in fig. I. All models with spectral index $`n_s=1`$. The data points are from.
FIG. 5. Evolution of the scalar potential $`V(\mathrm{\Phi })`$ in terms of the scalar field $`\mathrm{\Phi }`$.
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# Are the Laws of Thermodynamics Consequences of a Fractal Properties of Universe?
## I Introduction
It is well known, that the multifractal sets have the characteristics very similar to the characteristics of a physical quantities (a free energy, an entropy, a temperature etc) with which these characteristics can be formally compared. The connections between the characteristics of multifractal set and characteristics of physical quantities formally correspond to connections between the thermodynamic quantities. This surprising correspondence till now is completely inexplicable. In the present paper the multifractal analysis advanced in by Mandelbrot , , Renyi , Halsy , etc. (see for example, Rudolf ) was used for a substantiation of thermodynamics laws on the base of the supposition that the space and the time are multifractal sets. If our Universe state is the state of equilibrium (or the state nearly equilibrium) the connections between the global characteristics of Universe as a whole (multifractal time and space) and their local fractal characteristics will be the same that thermodynamic laws (it is shown on the basis of the fractal theory of time and space (see \- ). From the point of view of the fractal theory of time and space the thermodynamics relations (as well as thermodynamics in whole) are consequence of the multifractal structure (structure of time and space) of our Universe.
In the theory \- the time and the space are treated as a real physical fields. These fields consist of small multifractal subsets of time and space (”elements” of time and space), in turn, approximately treated as ”points”. Multifractal sets of time and space defined on a set of the carrier of a measure $`R^n`$, and contain all characteristics of the real world by reflected it’s in their fractional dimensions. The fractal dimensions $`d_{tr}`$ (or $`d_t`$ and $`d_r`$) in small neighborhood of points $`t,r,`$ for these ”points” (belonging the sets of $`t`$) are global dimensions. At the same time for all space - time continuum these dimensions are local fractal dimensions (Gelder exponents). The purpose of the paper is the establishment of connection between the global and the local characteristics of multifractal space and time on the basis of the multifractal analysis. We suppose that the state of the Universe (consisting from multifractal time and spatial sets) may be described as the state close to a thermodynamic equilibrium. The establishment of such connections enables on to view the new reason of origin of thermodynamic relations existing in our world, reducing it to presence of the fractal properties at time and space. We shall show, that the thermodynamic relations used in physics are a natural consequence of known mathematical connections between the multifractal characteristics of the Universe (Universe is considered as multifractal space - time set described within the framework of the fractal theory of time and space \- ). Thus thermodynamics can be considered as a natural consequence of multifractal characteristics of time and space of the world in which we live.
## II Connection Between the Physical and Multifractal Characteristics in the Multifractal Universe
Let’s consider the Universe as a dynamic system at the state that close to a thermodynamic equilibrium (at the present stage of it’s development), defined on a multifractal set $`X`$. Let the state of the Universe is characterized by fractal dimensions of a space - time continuum and by mean values of an internal energy, a free energy, an entropy, a temperature. If the state of Universe is close to the thermodynamic equilibrium, it’s characteristics is possible to describe by a free energy $`F`$, entropy $`S`$, internal energy $`E`$ and temperature $`T`$. Let the Universe has a multifractal nature stipulated by fractional dimensions of time and space (according to the fractal theory of time and space \- ) and is characterized by multifractal set $`X(r,t)`$ it’s space \- time points. The multifractal set $`X`$ is defined on the carrier of a measure (set $`R^n`$ with topological dimensions), i.e. $`XisthesubsetofR^n`$. Let the set of the carrier of a measure is characterized by the temperature $`T_0`$, by the internal energy $`E_0=T_0`$ (in the system in which the Boltzmann constant is equal to unity), by a free energy $`F_0`$. Let’s define a measure $`\mu `$ on the set $`X`$ and consider connection between invariant scaling characteristics of the multifractal Universe with a measure $`\mu `$ on the basis of the theory \- and hypothesis about the origin of the Universe as a result of explosion (big bang). Because of multifractality of space - time sets, the scaling transformations (at measuring volume of the Universe with the help of covering by four- dimension orbs (or cubes) with radius $`\delta `$), for example, for mean values of probability of the casual mass distribution $`<p^q>`$ (or the random distributions of densities of energy of physical fields), will look like (see, for example, \- )
$$<p^q>\delta ^{\tau (q+1)}$$
(1)
where $`q`$ is scale factor bound with $`q`$-dimensions Renyi $`dim_B^q(X)`$ by relation
$$dim_B^q(X)=\frac{\tau (q)}{q1}$$
(2)
The dimension Renyi characterizes global scailing characteristics of the Universe. For definition of it’s physical sense we shall consider local properties of the Universe near to the point $`r,t`$. The local fractal dimensions in this point (Gelder’s exponent) according to \- looks like
$$\alpha (x)d_{t,r}(\stackrel{}{r}(t),t)=4+\underset{i}{}\beta _iL_{i,t\stackrel{}{r}}(\stackrel{}{r}(t),t)$$
(3)
where $`L_{i,t\stackrel{}{r}}`$ are densities of energy of physical fields in this point and characterized by the densities of Lagrangians. The quantity $`p_i`$ in a neighborhood of a point $`(r,t)`$ is transformed as
$$p_i\delta ^{d_{tr}(\stackrel{}{r}(t),t)}$$
(4)
From definition of q-dimensions Renyi
$$dim_B^q(\stackrel{}{r},t)=\frac{1}{q1}lim_{\delta 0}\frac{log\underset{i}{}p_i^q}{log\delta }$$
(5)
follows
$$dim_B^q(\stackrel{}{r},t)=\frac{qd(\stackrel{}{r}(t),t}{q1}$$
(6)
The fractal dimensions $`d(r,t)`$ in (3) are the dimensionless internal energies (after multiplication on $`E_0`$ the relation (3) and correspond an internal energies of Universe in a point with coordinates $`(r,t))`$ and so, for $`q>>1`$, follows
$$dim_B^q(\stackrel{}{r},t)d(\stackrel{}{r}(t),t)$$
(7)
Therefore the $`dim_B^q(X)`$ should has sense of an energies. For describing of a thermodynamic equilibrium of the Universe there are only two energies (internal $`E`$ and free $`F`$) and $`E_0`$ is bound with $`d_{tr}(r,t)`$, therefore the dimensions Renyi there corresponds to a free energy of the Universe ($`F`$ divided by $`E_0`$ ) in the $`q`$ \- state. Let’s define now $`q`$ \- state. From (2) follows, as the Universe cools down and also it’s temperature is decrease and it’s volume grows, that $`q`$ must depends on temperature and will increase with Universe cooling. The simplest dimensionless function satisfying to this requirement is the function
$$T=T_0/T$$
(8)
Now it is necessary to define a function of state of the Universe - the entropy $`S`$. Let’s consider subsets $`S^{}(\alpha )`$ (of the set $`X`$) with identical Gelder’s exponents $`d_{rt}=\alpha `$ (in our case it corresponds to a selection of an isoenergetic sets of the ”internal energy” of the Universe). In this case joining of subsets $`S^{}(\alpha )`$, stratifying original set $`X`$, will coincide with the original set. Let’s introduce a spectrum of fractal dimensions $`f(\alpha `$ ).The joining of all such subsets makes set $`X`$. Let’s the fractal dimensions of set $`S(\alpha (q))`$ (obtained as a result of such stratifying) is $`f(\alpha (q))`$ (spectrum of singularities). For each value of $`q`$ the state of the Universe is determined as a single-valued state and at alteration $`q`$ (that is decreasing of energy of the Universe because of decreasing of it’s temperature) and expansion of the Universe) function $`f(\alpha (q))`$ , describing scaling properties of set $`S(\alpha (q))`$, will grows. Such behavior corresponds to behavior of an entropy (a $`q`$ \- entropy) which we shall designate by $`S`$. Hence, to the every state of Universe there corresponds a spectrum of singularities $`f(\alpha (q))`$ equal to an $`q`$ \- entropy $`S`$.
## III Connection of a Free Energy and an Entropy as a Consequence of a Multifractal Nature of the Universe
We use now the known relation of the multifractal analysis between q-dimensions Renyi, spectrum of a singularity $`f(\alpha (q))`$ and local fractal dimensions $`q`$, (see, for example, )
$$(q1)dim_B^q(X)=q\alpha (q)f(\alpha (q))$$
(9)
For $`T<<T_0`$, substituting in (9) instead of dimensions Renyi the spectrum of singularities $`f(\alpha (q))`$, the local fractal dimensions $`\alpha (q)`$ and the scaling factor $`q`$ their physical values (that we have received earlier) the relation reads
$$F=ETS$$
(10)
The relation (10) is the basic relation of the thermodynamics. As relation (10) is fulfilled for the Universe as a whole, it will be fulfilled and for it’s parts with the state of thermodynamically equilibrium. Therefore in the Universe with the multifractal time and space the realization of the laws of thermodynamics is a simple consequence of it’s structure.
The analysis of connections of global dimensions and local fractal characteristics of the fractal space - time carried out above allows to make the following statements, that are true for a case of equilibrium (or nearly so by equilibrium) of the state of the Universe:
a)The free energy of the Universe $`F`$ can be viewed as fractal $`q`$ \- dimensions Renyi $`(q=T_0/T)`$ of space - time set $`X`$ that consist the Universe
$$dim_B^{T_0/T}(\stackrel{}{r},t)=\frac{T}{T_0T}lim\frac{log(\underset{i}{}\mu _i^{T_0/T})}{log\delta }=F$$
(11)
where $`\mu _i`$ a measure of $`i`$-th of four-dimensional element of space - time;
b) The inverse temperature of the Universe $`T_0/T`$ corresponds to the $`q`$-characteristics of the scaling transformation of multifractal space - time;
c) The entropy of Universe $`S`$ corresponds the spectrum of fractal dimensions $`f(d_{tr}(T_0/T))`$, defined by dependencies of space-time of dimensions Renyi $`dim_B^{T_0/T}(\stackrel{}{r},t)`$ , mean temperature $`T_0/T`$ and local fractal dimensions of space - time sets with identical energy $`d_{tr}(T_0/T)`$;
d) The knowledge of the fractal spectrum and dimensions $`d_{tr}(q)`$ allows to find dimensions Renyi from (11). If the dimensions Renyi is known, the differentiation (\****6.10) with respect to $`q`$ gives in the equation
$$d_{tr}(q)=\frac{d}{dq}[(q1)dim_B^q(\stackrel{}{r},t)$$
(12)
It is possible to find, using (9), the entropy (i,e. the spectrum of fractal dimensions $`f(T_0/T)=f(q))`$;
e) The thermodynamics in viewed model is a consequence of the multifractality of space - time continuum.
## IV Conclusion
The problem of a substantiation of the thermodynamics within the frame- work of the fractal theory of time and space presented in this paper, (as well as a substantiation of irreversibility of time and spatial events (see ) is reduced to a postulating of multifractal properties of space and time. If model of fractal time and spaces \- is correct, the Universe is open system and exchange it’s energy with the carrier of a measure $`R^n`$ (or with the alien Universe of inflationary model or model ).
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# The Representation of Natural Numbers in Quantum Mechanics
## I Introduction
As is well known numbers play an essential role in physics and in many other disciplines. The results of both experimental work and theoretical computations are often given as numbers. Comparison of these numbers is essential to the validation process for any physical theory such as quantum mechanics. As inputs to or outputs of computations or experiments, numbers correspond to states of physical systems. From an information theoretic viewpoint, this correspondence is essential as these states carry information. As Landauer has emphasized, ”Information is Physical” . This is taken very seriously here.
However, the fact that many states of many physical systems correspond to numbers, has, for the most part, been assumed and used implicitly. There has been little attempt to make explicit the assumptions and conditions involved in representing numbers by states of physical systems.
This paper represents one approach to making some of the assumptions and conditions explicit. The emphasis is on the mathematical and physical aspects in the representation of numbers by states of physical systems. No new models of computation are presented. However making the assumptions explicit does offer some insight into the importance of various conditions that may not have been realized so far. An example (Section V) is the essential role played by the condition that there exist physically realizable dynamical operators that can efficiently implement basic arithmetic operations. The fact that these conditions are satisfied for a wide variety of systems, as shown by the ubiquitous existence of computers, does not detract from their importance.
Such a study is also relevant to the development of a coherent theory of mathematics and physics together, which, in one form or another, is a goal of many physicists . Any such coherent theory must take account in detail of how numbers are represented by states of physical systems.
In this paper considerations will be limited to quantum systems. This is not a serious limitation because of the assumed universal applicability of quantum mechanics or a related theory such as quantum field theory. In this case all physical systems are quantum systems and all states of these systems are (pure or mixed) quantum states. This is the case whether the systems are microscopic or macroscopic or whether macroscopic systems can be described by classical mechanics.
For quantum systems numbers are represented by tensor products of states of different degrees of freedom of a system. Usually the system is composite with each degree of freedom associated with a component system. For microscopic systems one condition systems must satisfy is that they have states for which the switching time, $`t_{sw}`$, is short compared to the decoherence time $`t_{dec}`$ or $`t_{sw}t_{dec}`$ . This is a dynamic condition as it is based on the Hamiltonian for the systems including their interaction with other systems and the environment.
This condition eliminates many state spaces of microscopic quantum systems for representation of numbers. A 2-dimensional example would be the state space based on two highly excited states of nuclei that have halflives short compared to $`t_{sw}`$. On the other hand spin projection states of spin $`1/2`$ ground state nuclei in in molecules in a magnetic field are suitable and are used in NMR quantum computers .
Macroscopic quantum systems are such that $`t_{sw}t_{dec}`$ for all states of interest. In this case the systems are candidates for number representation for classical computation if the systems have states that are stabilized by environmental interactions for times long compared to the switching time. The widespread existence of macroscopic computers shows that both $`t_{sw}t_{dec}`$ and environmental stabilization occurs for many quantum systems.
Because of the recent widespread interest in quantum computing, the emphasis of this paper is on number representation by states of microscopic quantum systems. However most of the material also applies to macroscopic systems.
The first step in giving an exact meaning to the representation of numbers by tensor product states of quantum systems is to specify exactly what natural numbers are. Without such a specification all computations are meaningless physical operations. Here the axiomatic approach is used by defining a nonempty set as a set of all natural numbers if it is a model of the axioms of number theory or arithmetic . These axioms are discussed in the next section along with changes needed to account for the limitation of this paper to tensor product states with an arbitrary but fixed finite number $`L`$ of components, or $`kary`$ representations of length $`L`$. The corresponding arithmetic becomes arithmetic modulo $`k^L`$.
It is possible to model the axioms directly on a physical Hilbert space $`^{phy}`$ describing a composite quantum system with $`L`$ components. However the literature on quantum computing makes much use of product qubit states of the form $`|\underset{¯}{s}`$ where $`\underset{¯}{s}`$ is any function from $`1,2,\mathrm{},L`$ to $`\{0,1\}`$. Since the Hilbert space of these states is a very useful reference base for discussing quantum computation that is independent of any physical model, this approach will be used here.
To this end a purely mathematical model of these axioms is described in Section III that is based on a tensor product Hilbert space $`^{arith}=_{j=1}^L_j`$ of $`L`$ $`k`$ dimensional Hilbert spaces $`_j`$. Unitary operators on this space are defined to correspond to the basic arithmetic operations, successor, plus, and times, whose properties are given by the axioms. The presence of $`L`$ successor operators, one for each power of $`k`$, rather than just one as described by the axioms, is based on the condition of efficient implementation discussed later on.
Tensor product states of physical properties of microscopic composite quantum systems belonging to the Hilbert space $`^{phy}`$ are discussed in Section IV. A priori these states, as products over a label set $`A`$ of $`L`$ physical parameter values, do not correspond to any number. Also operators on these states are meaningless regarding any numerical interpretation.
This is remedied by describing tensor product preserving unitary operators from $`^{arith}`$ to $`^{phy}`$. For each of these operators, $`^{phy}`$, along with induced representations of the operators for the basic arithmetic operations, becomes a model for the axioms of modular arithmetic.
So far nothing has been said about the physical realizability of any of these models of the axioms. This is especially relevant for the operators as they are many system nonlocal operators. This is remedied in Section V where the important condition of efficient implementability of the basic arithmetic operations is described. In essence the condition requires that a composite quantum system be such that there exist physically realizable Hamiltonians that can implement the basic arithmetic operations. In addition the space-time and thermodynamic resources required for implementation must be polynomial in $`L`$. The importance of this condition rests in the fact that it is additional to and independent of the axioms of arithmetic. To see this one notes that there are many models of the axioms that do not satisfy this requirement. A simple physical model is any one based on an unary representation of the numbers as most arithmetic operations are inefficient in this representation.
The question arises if use of $`^{arith}`$ can be bypassed by modeling the axioms directly on $`^{phy}`$ where $`^{phy}`$ has an arbitrary tensor product structure. In general this is possible as any structure satisfying the axioms is acceptable. The discussion of this in Section VI is based on a description of properties of a set of operators indexed by a set of physical parameters. The properties are also defined to address the question of necessary and sufficient conditions to conclude that $`^{phy}`$ must have a tensor product structure suitable for length $`L`$ $`kary`$ representations of numbers.
A final section discusses some other aspects and open questions resulting from this work. The importance of the efficient implementability condition in excluding most models of modular arithmetic on $`^{phy}`$ is noted as are some aspects of the use of numbers to describe $`kary`$ representations of length $`L`$.
It must be emphasized that the work of this paper is one attempt to make explicit the assumptions and conditions that are assumed implicitly in the representation of numbers by states of quantum systems and in work in the the literature on quantum computing. Examples of this work are given in papers by Beckman et. al. and Vedral et.al. that describe networks of quantum gates to carry out basic arithmetic operations. The description is in terms of unitary operators on $`^{arith}`$ (extended to include ancillary qubits) as ordered products of polynomially many elementary gate operators. The distinction between physical models, with the associated requirement of efficient physical implementation, and mathematical models is not maintained (and is not needed) in the papers. Also efficient physical implementability implies more than minimizing the number of ancillary qubits and restriction to polynomially many gate operations. These aspects are discussed more in Sections III and V.
## II The Axiomatic Description of Numbers
The first step in making explicit what is involved in the representation of numbers by quantum states is to define the natural numbers. One method of doing this is to follow mathematical logic and define any nonempty set to be a set of natural numbers if it is a model for the axioms of arithmetic or number theory . A model for any axiom system is a collection of elements in which all the axioms are true.
Here the main interest is models of arithmetic based on Hilbert spaces that are tensor products of an arbitrary but fixed number of component spaces. As a result the axioms to be satisfied are those for arithmetic modulo $`N`$ where $`N`$ is arbitrary but fixed. This arithmetic satisfies some of the axioms for all natural numbers. Others need to be either deleted or modified. It also satisfies axioms for a commutative ring with identity .
The exact form and content of axioms for modular arithmetic is not important here. What is important is that both the arithmetic and ring axioms have in common the required existence of binary operations $`+`$ and $`\times `$ with certain properties. Also an unary successor operation $`S`$ is required by the arithmetic axioms. The properties the binary operations must have include commutativity, associativity, the existence of identities $`0`$ and $`S(0)`$ for $`+`$ and $`\times `$, and the distributivity of $`\times `$ relative to $`+`$. Also $`S`$ commutes with $`+`$ and $`x\times S(y)=x\times y+x`$ for all $`x,y`$.<sup>*</sup><sup>*</sup>*The importance of these axioms lies in the requirement of the existence of binary operations $`+,\times `$ with certain properties. The fact that some of the axioms may be redundant is of no importance here. The arithmetic axioms defining an order relation and the induction schema are not considered as they are not needed for the purposes of this paper. However it is useful to keep in mind that the ordering axioms establish the discreteness of the natural numbers in the sense that there is no number between $`x`$ and its successor $`S(x)`$.
## III Abstract Hilbert Space Models
The next step is the description of a purely mathematical model of these axioms based on whatever mathematical systems are appropriate for the physical systems being considered. Since interest here is in $`kary`$ representations of length $`L`$ of natural numbers for composite quantum systems, a model based on an abstract Hilbert space $`^{arith}`$ is needed. To this end let $`^{arith}=_{j=1}^L_j`$ be an $`L`$ fold tensor product Hilbert space where $`_j`$ is a $`k`$ dimensional Hilbert space.
For each $`j`$, the basis states of interest in $`_j`$ have the form $`|\mathrm{},j`$ where $`j`$ denotes the label or property characterizing a qubyte and $`\mathrm{}=0,1,\mathrm{},k1`$. A product state basis in $``$ can be given in the form $`|\underset{¯}{s}=_{j=1}^L|\underset{¯}{s}(j),j`$ where $`\underset{¯}{s}`$ is any function from $`1,\mathrm{},L`$ to $`0,\mathrm{},k1`$. (Here qubits or qubytes refer to quantum bits or bytes of information for $`k=2`$ or $`k2`$ respectively.)
The presence of the parameter $`j`$ in the state and not as a subscript, as in $`_{j=1}^L|\underset{¯}{s}(j)_j`$, is required as the action of operators corresponding to the basic arithmetic operations depends on the value of $`j`$. It is not possible to express this dependence if $`j`$ appears as a subscript of $``$ and not between $`|`$ and $``$.
An important function of the axioms is to provide properties of the unary operation $`S`$ and the binary operations $`+`$ and $`\times `$. For reasons based on efficient implementation (Section V), it is quite useful to define $`L`$ different successor operators, $`V_j^{+1}`$, for $`j=1,\mathrm{},L`$. These operators are defined to correspond to the addition of $`k^{j1}modk^L`$ where $`V_1^{+1}`$ corresponds to $`S`$ in the axioms. These operators and those for $`+`$ and $`\times `$ correspond to the basic arithmetic operations.
It is to be emphasized that definitions of the $`+,`$ and $`\times `$ are given to show their dependence on the $`V_j^{+1}`$. Also they are required by the axioms of arithmetic. The purpose is definitely not to present the definitions as something new as these operators are widely used.
For instance the widely discussed networks of quantum gates are examples of the abstract models considered here for $`k=2`$. In the networks the states in $`^{arith}`$ are represented by horizontal qubit lines and ordered products of gate operators represent operators in $`B(^{arith})`$. Specific examples of this for the basic arithmetic operations of addition, multiplication and modular exponentiation are described in .
It is also the case that in many physical models space and time directions can be assigned to the abstract networks. In this case the spatial ordering of the qubit lines is part of any mapping of the abstract models based on $`^{arith}`$ to physical models based on $`^{phy}`$ in which the corresponding component systems are distinguished by spatial positions. These mappings are examples of the mappings ”$`g`$” discussed in Section IV. The time ordering of the quantum gates corresponds to mapping the ordering of gate operators in the abstract model to a time ordered product of physically implementable quantum gate operators. This is part of the requirement of physical implementability (Section V).
### A Definitions of the $`V_j^{+1}`$
The definition of the $`V_j^{+1}`$ is straightforward. For each $`j`$ let $`u_j`$ be a cyclic shift of period $`k`$ that acts on the states $`|\mathrm{},j`$ according to $`u_j|\mathrm{},j=|\mathrm{}+1modk,j`$. $`u_j`$ is the identity on all states $`|m,j^{}`$ where $`j^{}j`$. Define $`V_j^{+1}`$ by
$$V_j^{+1}=\{\begin{array}{cc}u_jP_{(k1),j}+V_{j+1}^{+1}u_jP_{(k1),j}\hfill & \text{if }1j<L\hfill \\ u_L\hfill & \text{if }j=L\hfill \end{array}$$
(1)
Here $`P_{(k1),j}=|k1,jk1,j|1_j`$ is the projection operator for finding the $`j`$ component state $`|k1,j`$ and the other components in any state. $`P_{m,j}`$ and $`u_j`$ satisfy the commutation relation $`u_jP_{m,j}=P_{m+1,j}u_jmodk`$ for $`m=0,\mathrm{},k1`$. Also $`P_{(k1),j}=1P_{(k1),j}`$. This follows from the fact that the label spaces for each qubyte are one dimensional so that the operator $`1_j|jj|1_j`$ is the identity on the Hilbert space spanned by the $`k^L`$ states $`|\underset{¯}{s}`$
This definition is implicit in that $`V_j^{+1}`$ is defined in terms of $`V_{j+1}^{+1}`$. An explicit definition is given by
$`V_j^{+1}`$ $`=`$ $`{\displaystyle \underset{n=j}{\overset{L}{}}}u_nP_{(k1),n}{\displaystyle \underset{\mathrm{}=j}{\overset{n1}{}}}u_{\mathrm{}}P_{(k1),\mathrm{}}`$ (3)
$`+{\displaystyle \underset{\mathrm{}=j}{\overset{L}{}}}u_{\mathrm{}}P_{(k1),\mathrm{}}`$
In this equation the unordered product is used because for any $`p,q`$, $`u_mP_{p,m}`$ commutes with $`u_nP_{q,n}`$ for $`mn`$. Also for $`n=j`$ the product factor with $`j\mathrm{}n1`$ equals $`1`$.
There are two basic properties the operators $`V_j^{+1}`$ must have: they are cyclic shifts and, for each $`j<L`$, they satisfy
$$(V_j^{+1})^k=V_{j+1}^{+1}$$
(4)
Also if $`j=L`$ then $`(V_L^{+1})^k=1`$. To show that $`V_j^{+1}`$ is a shift, let $`|\underset{¯}{s}`$ be a product state such that for each $`m=1,2,\mathrm{},L`$ the component states $`|\underset{¯}{s}_m,u_m|\underset{¯}{s},(u_m)^2|\underset{¯}{s},\mathrm{},(u_m)^{k1}|\underset{¯}{s}`$ are pairwise orthonormal. It then follows from Eq. 3 and the properties of the $`u_m`$ that any product state $`|\underset{¯}{s}`$ is orthogonal to the state $`V_j^{+1}|\underset{¯}{s}`$ and that $`V_j^{+1}`$ is norm preserving on these states.
Assume that Eq. 4 is valid. Then for each $`j`$ $`(V_j^{+1})^{k^{Lj+1}}=1`$. This, and the facts that for all tensor product states $`|\underset{¯}{s},V_j^{+1}|\underset{¯}{s}`$ is also a tensor product state which is orthogonal to $`|\underset{¯}{s}`$, show that $`V_j^{+1}`$ is a cyclic shift. The existence of a tensor product basis that is common to all the $`V_j^{+1}`$ follows from Eq. 4.
To prove Eq. 4 it is easiest to use Eq. 1. Since $`V_{j+1}^{+1}`$ commutes with $`u_{\mathrm{}}P_{n,\mathrm{}}`$ for all $`\mathrm{}j`$ and the commutation relations $`P_{n,j}u_j=u_jP_{(n1),j}`$ and $`P_{n,j}u_j=u_jP_{(n1),j}`$ hold, one has for each $`mk`$
$$(V_j^{+1})^m=(u_j)^m\underset{\mathrm{}=1}{\overset{m}{}}P_{(k\mathrm{}),j}+V_{j+1}^{+1}(u_j)^m(\underset{\mathrm{}=1}{\overset{m}{}}P_{(k\mathrm{}),j}).$$
Here $`P_{n,j}=1P_{n,j}`$. For $`m=k`$ the term with the product of the projection operators gives $`0`$ and the sum of the projection operators gives unity. The desired result follows from the fact that $`(u_j)^k=1`$. Also $`(V_L^{+1})^k=1`$ follows directly from the definition of $`V_L^{+1}`$.
The above shows that informally the action of $`V_j^{+1}`$ corresponds to addition $`modk^L`$ of $`k^{j1}`$ on the product basis. This cannot yet be proved as addition $`modk^L`$ has not yet been defined. Also the adjoint $`(V_j^{+1})^{}`$ of $`V_j^{+1}`$ corresponds informally to subtraction $`modk^L`$ of $`k^{j1}`$. This can be seen from the fact that $`(V_j^{+1})^{}V_j^{+1}=1`$ where
$`(V_j^{+1})^{}`$ $`=`$ $`{\displaystyle \underset{n=j}{\overset{L}{}}}P_{(k1),n}u_n^{}{\displaystyle \underset{\mathrm{}=j}{\overset{n1}{}}}P_{(k1),\mathrm{}}u_{\mathrm{}}^{}`$ (6)
$`+{\displaystyle \underset{\mathrm{}=j}{\overset{L}{}}}P_{(k1),\mathrm{}}u_{\mathrm{}}^{}.`$
This result is obtained using the commutativity of the shifts and projection operators for different component systems.
It should be noted that the operators $`V_j^{+1}`$ play an important role in quantum computation. This is the case even though for each product state $`|\underset{¯}{s}`$ the state $`V_j^{+1}|\underset{¯}{s}`$ is also a product state and is not a linear superposition of these states. The importance comes from the fact that these operators along with their efficient implementation are used to define the basic arithmetic operations for a quantum computer and to carry out quantum algorithms. For example in Shor’s factoring quantum algorithm , they are used in the step in which the function $`f_y(s)=y^smodN`$ is calculated for each component state $`|\underset{¯}{s}`$.
### B Plus
It is straightforward to define the plus ($`+`$) operation in terms of the $`V_j^{+1}`$. To ensure unitarity the definition will be based on states of the form $`|\underset{¯}{s,w}=|\underset{¯}{s}|\underset{¯}{w}`$ that describe two $`L`$ qubyte product states.
To define the $`+`$ operation let $`V_j^+\mathrm{}=(V_j^{+1})^{\mathrm{}}`$ represent $`\mathrm{}`$ iterations of $`V_j^{+1}`$. Then $`+`$ is defined by
$`+|\underset{¯}{s}|\underset{¯}{w}`$ $`=`$ $`|\underset{¯}{s}V_L^{+s_L}V_{L1}^{+s_{L1}}\mathrm{}V_2^{+s_2}V_1^{+s_1}|\underset{¯}{w}`$ (8)
$`=|\underset{¯}{s},\underset{¯}{s+w}`$
Here the numeral expression $`|\underset{¯}{s+w}`$ is defined to be that generated from $`\underset{¯}{w}`$ by the action of the product $`_{j=1}^LV_j^{+s_j}`$. Note that the different $`V_j^{+1}`$ commute.
For pairs of product states, which are first used here, the domains of the functions $`\underset{¯}{s}`$ and $`\underset{¯}{w}`$ must be different. This is based on the requirement that an algorithm must be able to distinguish components of $`|\underset{¯}{s}`$ from components of $`|\underset{¯}{w}`$. This can be achieved by setting $`|\underset{¯}{s}|\underset{¯}{w}=|\underset{¯}{sw}`$ where $`\underset{¯}{sw}`$ denotes the concatenation of $`\underset{¯}{w}`$ to $`\underset{¯}{s}`$. That is, $`\underset{¯}{sw}`$ is a function from $`1,\mathrm{},2L`$ to $`0,\mathrm{},k1`$ where $`\underset{¯}{sw}(h)=\underset{¯}{s}(h)`$ for $`hL`$ and $`\underset{¯}{sw}(h)=\underset{¯}{w}(hL)`$ for $`h>L`$.
As defined the $`+`$ operator is unitary on the Hilbert space spanned by all pairs of length $`L`$ numeral expression states. Thus a reversible implementation of it is possible where the procedure makes use of the procedures for implementing the $`V_j^{+1}`$. Eq. 8 shows that the procedure can be carried out by carrying out, for each $`j=1,2,\mathrm{},L`$, $`s_j`$ iterations of $`V_j^{+1}`$ where $`s_j`$ is the number $`\underset{¯}{s}(j)`$ associated with the qubyte state $`|\underset{¯}{s}(j),j`$ in $`|\underset{¯}{s}=_{j=1}^L|\underset{¯}{s}(j),j`$. Since $`+`$ is unitary, so is the adjoint, $`+^{}`$. Since $`+`$ was defined to correspond to addition modulo $`k^L`$, the adjoint corresponds to subtraction modulo $`k^L`$. That is if $`+|\underset{¯}{s}|\underset{¯}{w}=|\underset{¯}{s}|\underset{¯}{s+w}`$ then $`+^{}|\underset{¯}{s}|\underset{¯}{s+w}=|\underset{¯}{s}|\underset{¯}{w}`$.
### C Times
Here a definition of multiplication is given that is based on efficient iteration of $`+`$ and is similar to the method taught in primary school. The method is efficient relative to that for $`+`$.
Reversibility of the operations requires that the operator $`\times `$ be unitary. (Caution: the adjoint of $`\times `$ is not division.) This means that both input product states and the product state with the result must be preserved. It is also convenient to have one extra product state for storing and acting on intermediate results. This state begins and ends as $`|\underset{¯}{0}`$. For initial states of the form, $`|\underset{¯}{s},\underset{¯}{w},\underset{¯}{0},\underset{¯}{0}=|\underset{¯}{s}|\underset{¯}{w}|\underset{¯}{0}|\underset{¯}{0}`$,
$$\times |\underset{¯}{s},\underset{¯}{w},\underset{¯}{0},\underset{¯}{0}=|\underset{¯}{s},\underset{¯}{w},\underset{¯}{0},\underset{¯}{s\times w}$$
(9)
where $`|\underset{¯}{s\times w}`$ is the state resulting from the action of $`\times `$. It is supposed to correspond to the result of multiplying, $`modk^L`$, the numbers corresponding to the states $`|\underset{¯}{s}`$ and $`\underset{¯}{w}`$.
In order to define $`\times `$ explicitly one needs to be able to generate the states $`|\underset{¯}{k^{j1}\times w}`$ corresponding to multiplication of $`w`$ by $`k^{j1}`$. For each $`j=1,\mathrm{},L`$ these states are added to themselves $`s_j`$ times. The final result is obtained by adding all the resulting states so obtained. Details are provided in the Appendix.
### D Required Properties of the $`V_J^{+1}`$, Plus, Times
As was noted the operators $`V_j^{+1},+,\times `$ must satisfy the properties expressed by the axioms for modular arithmetic. These include the axioms for arithmetic modified for modularity and the p\[resence of $`L`$ successors, and possibly axioms for a commutative ring with identity .
Properties that must be satisfied include that expressed by Eq. 4 and the requirements that the successor operations commute with $`+`$, (i.e. $`+(1V_j^{+1})=(1V_j^{+1})+`$), the existence of additive and multiplicative identities, which are the states $`|\underset{¯}{0}`$ and $`|\underset{¯}{1}=V_1^{+1}|\underset{¯}{0}`$, and the distributivity of $`\times `$ over $`+`$. Also $`+`$ and $`\times `$ are associative and commutative.
Proof of these properties from the definitions and Eq. 4, which has already been proved, is straight forward and will not be given here. Note that the proofs of some of the properties do use the corresponding properties of the numbers appearing in the exponents. For example to prove that addition is commutative, $`|\underset{¯}{s+w}=|\underset{¯}{w+s}`$, Eqs. 8 and 3 give $`|\underset{¯}{s+w}=_{h=1}^L(V_j^{+1})^{s_h+w_h}|\underset{¯}{0}`$ and $`|\underset{¯}{w+s}=_{h=1}^L(V_j^{+1})^{w_h+s_h}|\underset{¯}{0}`$. The equality of these two states follows from $`s_h+w_h=w_h+s_h`$ for each $`h`$.
## IV Physical Hilbert Space Models
The Hilbert space models described so far are purely abstract in that they do not refer to any physical properties. They do however, serve as a common reference point for models based on physical properties of physical systems. They also give a useful method to associate numbers with quantum states of these systems.
To begin, let $`A`$ and $`B`$ be sets of $`L`$ and $`k`$ different physical parameters or values of some physical properties or observables $`\widehat{A}`$ and $`\widehat{B}`$. The $`A`$ parameters are used to distinguish or label different components of a composite quantum system and $`B`$ is a set of values of a different physical property associated with each component system. For example $`A`$ could be a set of $`L`$ arbitrary locations of component spin $`1/2`$ systems on a 2 dimensional surface and $`B=\{,\}`$ denoting spin aligned along or opposite some axis of quantization. Another example, representative of NMR quantum computation , has $`A`$ as a set of hyperfine splittings of nuclear spin states and $`B=\{,\}`$. Here the values of $`A`$ must contain sufficient information so the physical process can distinguish between the different nuclear spins.
Let $`\underset{¯}{t}`$ be any function from $`A`$ to $`B`$ and $`|\underset{¯}{t}=_{aϵA}|\underset{¯}{t}(a),a`$ be the corresponding tensor product state. Let $`^{phy}=_{aϵA}_a`$ be the $`k^L`$ dimensional Hilbert space spanned by all the states $`|\underset{¯}{t}`$. Each $`_a`$ is a $`k`$ dimensional Hilbert space spanned by states of the form $`|h,a`$ where $`hϵB`$.
The presence of $`a`$ as a separate part in each component state $`|\underset{¯}{t}(a),a`$, and not as a state subscript as in $`|\underset{¯}{t}(a)_a`$, is essential as an algorithm uses the value of $`a`$ to distinguish the different component systems. This is based on the view that the state of the composite quantum system contains all the quantum information available to the algorithm. In particular the states must contain sufficient information so that the algorithm can distinguish among the component systems. This is especially the case for any algorithm whose dynamics is described by a Hamiltonian that is selfadjoint and time independent. This is an example of Landauer’s dictum ”Information is Physical” .
This description can be generalized in that the physical property observable $`\widehat{B}`$ of the component systems can depend on the values of $`a`$ in $`A`$. An example this, which also has different component systems replaced by different degrees of freedom of one system, is shown by an ion trap example . Here the states of one degree of freedom are the ground and first excited state of the ion in the harmonic well trap. The corresponding states of the other are the ground and first excited electronic state of the ion. This type of generalization will not be pursued here.
### A Representation of Numbers and Arithmetic Operations in $`^{phy}`$
The goal here is for states in $`^{phy}`$ to represent numbers. However, it is clear that, a priori, neither the product states $`|\underset{¯}{t}=_{aϵA}|\underset{¯}{t}(a),a`$ nor linear superpositions of these states represent numbers. For the $`|\underset{¯}{t}`$ the reason is that there is no association between the labels $`a`$ and powers of $`k`$; also there is no association between the range set $`B`$ of $`\underset{¯}{t}`$ and the numbers $`0,1,\mathrm{},k1`$.
This can be remedied by use of unitary maps from $`^{arith}`$ to $`^{phy}`$ that preserve the tensor product structure. One way of doing this is to let $`g`$ and $`d`$ be any bijections (one-one onto) maps from $`1,2,\mathrm{},L`$ to $`A`$ and from $`0,1,\mathrm{},k1`$ to $`B`$. For each pair $`g,d`$ and each $`j`$ there is a corresponding unitary operator $`w_{g,d,j}`$ that maps states $`|h,j`$ in $`_j`$ where $`0hk1`$ to states in $`_{g(j)}`$ according to $`w_{g,d,j}|h,j=|d(h),g(j)`$. This induces a unitary operator $`W_{g,d}=_{j=1}^Lw_{g,d,j}`$ from the product space $`^{arith}`$ to $`^{phy}`$ where
$`W_{g,d}|\underset{¯}{s}=`$ $`_{j=1}^Lw_{g,d,j}|\underset{¯}{s}(j),j`$ (10)
$`=`$ $`_{j=1}^L|d(\underset{¯}{s}(j)),g(j)=|\underset{¯}{s}_g^d.`$ (11)
Here $`|\underset{¯}{s}_g^d`$ is the physical parameter based state in $`^{phy}`$ that corresponds, under $`W_{g,d}`$ to the number state $`|\underset{¯}{s}`$ in $`^{arith}`$.
This process can be inverted, using the adjoint $`W_{g,d}^{}`$ to relate physical parameter states in $`^{phy}`$ to number states in $`^{arith}`$. One has
$`W_{g,d}^{}|\underset{¯}{t}=`$ $`_{aϵA}w_{g,d,g^1(a)}^{}|\underset{¯}{t}(a),a`$ (12)
$`=`$ $`_{aϵA}|d^1(\underset{¯}{t}(a)),g^1(a)=|\underset{¯}{t}_{g^1}^{d^1}.`$ (13)
Here $`|\underset{¯}{t}_{g^1}^{d^1}`$ is the number state in $`^{arith}`$ corresponding to the physical state $`|\underset{¯}{t}`$. Note that $`W_{g,d}^{}=W_{g^1,d^1}`$ where $`g^1,d^1`$ are the inverses of $`g`$ and $`d`$, and $`w_{g^1,d^1,a}=w_{g,d,g^1(a)}^{}`$.
The operators $`W_{g,d}`$ also induce representations of the $`V_j^{+1},+,`$ and $`\times `$ operators on the physical parameter states in $`^{phy}`$. For the $`V_j^{+1}`$ one defines $`V_{g,j}^{d,+1}`$ by
$$V_{g,j}^{d,+1}=W_{g,d}V_j^{+1}W_{g,d}^{}.$$
(14)
An equivalent definition can be given by direct reference to the maps $`g,d`$ and the operators $`w_{g,d,j}`$:
$`V_{g,j}^{d,+1}`$ $`=`$ $`{\displaystyle \underset{n=j}{\overset{L}{}}}u_{g(n)}^dP_{d(k1),g(n)}{\displaystyle \underset{\mathrm{}=j}{\overset{n1}{}}}u_{g(\mathrm{})}^dP_{d(k1),g(\mathrm{})}`$ (16)
$`+{\displaystyle \underset{\mathrm{}=j}{\overset{L}{}}}u_{g(\mathrm{})}^dP_{d(k1),g(\mathrm{})}.`$
Here $`P_{d(k1),g(\mathrm{})}=w_{g,d,\mathrm{}}P_{k1,\mathrm{}}w_{g,d,\mathrm{}}^{}`$ and $`u_{g(\mathrm{})}=w_{g,d,\mathrm{}}u_{\mathrm{}}w_{g,d,\mathrm{}}^{}`$.
In a similar fashion one can use the $`W_{g,d}`$ to define the operator $`+_{g,d}`$ acting on the physical parameter states in $`^{phy}^{phy}`$. The definition is based on that given for the operator $`+`$ acting on $`^{phy}^{phy}`$ Eq. 8. One has
$$+_{g,d}=(W_{g,d}W_{g,d})+(W_{g,d}^{}W_{g,d}^{}).$$
(17)
The operator $`\times _{g,d}`$ is defined similarly from $`\times `$ as defined in the Appendix.
It is clear from the above that there is no unique correspondence between states in the arithmetic and physical Hilbert spaces. There are $`L!`$ possible bijections $`g`$ and $`k!`$ possible bijections $`d`$. Thus some or many of the $`L!k!`$ unitary operators $`W_{g,d}`$ associate a different physical parameter state $`|\underset{¯}{s}_g^d`$ with the number state $`|\underset{¯}{s}`$. Conversely the $`g`$ and $`d`$ dependence of $`W_{g,d}^{}`$ shows that many different number states $`|\underset{¯}{t}_g^d`$ can be associated with the physical state $`|\underset{¯}{t}`$. The multiplicity of these correspondences depends on the states $`|\underset{¯}{s}`$ or $`|\underset{¯}{t}`$ and the choices of $`g`$ and $`d`$.
It follows from the unitarity of $`W_{g,d}`$ that if the operators $`V_j^{+1},+,\times `$ and the states $`|\underset{¯}{s}`$ in $`^{arith}`$ satisfy the axioms of modular arithmetic, then so do the operators $`V_{g,j}^{d,+1},+_{g.d},\times _{g,d}`$ and states $`|\underset{¯}{s}_g^d`$ in $`^{phy}`$. In this way all the states $`|\underset{¯}{s}_g^d`$ in $`^{phy}`$ and the operators $`V_{g,j}^{d,+1},+_{g,d},\times _{g.d}`$ are a model of the axioms of modular arithmetic. The fact that superposition of the states $`|\underset{¯}{s}_g^d`$ plays an important role in quantum computation does not affect this conclusion.
This argument also applies to any unitary map $`U`$ from $`^{arith}`$ to $`^{phy}`$ independent of whether $`U`$ is tensor product preserving or not. However most of these maps are not of interest because the operators $`UV_j^{+1}U^{}`$ are not physically implementable (Section V). Also the states $`U|\underset{¯}{s}`$ nay not be stable or even preparable.
### B Grover’s and Shor’s Algorithms
Since the spaces $`^{arith}`$ and $`^{phy}`$, and arithmetic models constructed on these spaces are unitarily equivalent, one might think that dynamically an algorithm is independent of the unitary map used. This is not true in general even if one restricts the maps to have the form of $`W_{g,d}`$: some algorithms are independent of these maps and others are not.
To see this one notes that dynamically any quantum algorithm carried out on a composite physical system must be sensitive to the values of the physical parameters for the system. This means that the physical dynamics of an algorithm must be described by some evolution operator acting on the states in $`^{phy}`$ or some other physical model of the system states. The physical dynamics is not described on $`^{arith}`$.
It follows that any algorithm that can be described in terms of states based on physical parameters is independent of the unitary maps $`W_{g,d}`$. The dynamics does not depend on these maps because what number a physical state represents is irrelevant to the algorithm. On the other hand, algorithms that compute numerical functions must be described on $`^{arith}`$ as number is of the essence for these. It follows that the dynamics of these algorithms depends on the maps $`W_{g,d}`$.
Grover’s Algorithm and Shor’s Algorithm are examples of the two types of algorithm. Grover’s Algorithm corresponds to a quantum search of a set of data where each element of the data base corresponds to a quantum state. The goal is to find the one unknown but unique state with some property different from the others. Here the quantum state representing each data element will be taken to be a tensor product of qubit states. This is not necessary, as Lloyd has shown. However, the price for this is the need for an exponential overhead of resources.
Here the relevant feature of Grover’s Algorithm is that it can be both defined and implemented on $`^{phy}`$ with no reference to numbers represented by states in $`^{arith}`$. To see this let $`k=2`$ and $`B=\{,\}`$ for spin up, spin down. The initial state can be written as $`\psi =(1/\sqrt{N})_{\underset{¯}{t}}|\underset{¯}{t}`$ where $`|\underset{¯}{t}=_{aϵA}|\underset{¯}{t}(a),a`$ and $`N=2^L`$.
Dynamically Grover’s Algorithm consists of iterations of the unitary operator $`WI_\underset{¯}{}WI_{\underset{¯}{t_u}}`$ on $`^{phy}`$. Here $`I_\underset{¯}{}=12|\underset{¯}{}\underset{¯}{}|`$ where $`|\underset{¯}{}`$ is the state with all $`L`$ systems in the $`|`$ state. $`I_{\underset{¯}{t_u}}=12|\underset{¯}{t_u}\underset{¯}{t_u}|`$ and $`W`$ is the Walsh Hadamard transformation. Here $`|\underset{¯}{t_u}`$ is the unknown product state that is to be amplified, and $`W=_{aϵA}(1/\sqrt{2})(\sigma _x+\sigma _z)_a`$ is a tensor product of single qubit operators. The $`\sigma _x,\sigma _z`$ are the Pauli spin operators and $`\psi =W|\underset{¯}{}`$.
Shor’s Algorithm for finding the two prime factors of a large number is quite different in that it is essential that the tensor product states represent numbers. This can be seen from the steps of the algorithm
$`{\displaystyle \frac{1}{\sqrt{N}}}{\displaystyle \underset{\underset{¯}{s}}{}}|\underset{¯}{s}|\underset{¯}{i}{\displaystyle \frac{1}{\sqrt{N}}}{\displaystyle \underset{\underset{¯}{s}}{}}|\underset{¯}{s}|\underset{¯}{f_m(s)}`$ (18)
$``$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{\underset{¯}{w}}{}}|\underset{¯}{w}{\displaystyle \underset{\underset{¯}{s}}{}}e^{2\pi iws/N}|\underset{¯}{f_m(s)}`$ (19)
Here $`|\underset{¯}{i}`$ is the initial product state, usually shown as a constant sequence of $`0s`$. $`f_m`$ is a numerical function defined by $`f_m(x)=m^xmodM`$ where m and $`M`$ are relatively prime. The number $`M`$, which is to be factored, and $`N`$ are related by $`M^22^N2M^2`$ .
Eq. 19 shows that the dynamics of Shor’s algorithm can be initially formulated as a unitary step operator $`U_{Sh}`$ acting on $`^{arith}`$. However, physically, the dynamics is represented by the operator $`W_{g,d}U_{Sh}W_{g,d}^{}`$ acting on $`^{phy}`$. This shows that physically the dynamical implementation of Shor’s Algorithm depends on the numberings $`g`$ and $`d`$ of the physical parameter sets $`A`$ and $`B`$.
More generally the requirement that the numerical function calculated by the algorithm be invariant under any unitary map from $`^{arith}`$ to $`^{phy}`$ means that the physical implementation of the algorithm depends on the unitary map. For example, let $`W_{g,d}`$ be a unitary map as defined by Eq. 11 and $`A`$ be a set of space locations of spin $`1/2`$ systems with spin up $`()`$ spin down $`()`$ representing (through $`d^1`$) $`0,1`$. Then the algorithm dynamics clearly depends on $`g`$ as $`g`$ determines which space location is associated with which power of $`2`$. A similar argument holds for the dynamics dependence on $`d`$. Also the correct interpretation of the measurement of the output depends on both $`g`$ and $`d`$.
## V Efficient Implementability of Arithmetic Operations
Probably the most important requirement is that of efficient implementablity of basic arithmetic operations. This means that, for states of a physical system to represent numbers, it must be possible to physically implement these operations and the implementation must be efficient. This includes at least the operations described by the axioms as efficient implementation of these is a necessary condition for states of a quantum system to represent numbers.
In the case of the $`V_{g,j}^{d,+1}`$ physical implementability means there must exist a physically realizable Hamiltonian $`H_{g,j}^d`$ such that for some time $`t_j`$, $`U_{g,j}^d(t_j)=e^{iH_{g,j}^dt_j}`$ corresponds to carrying out $`V_{g,j}^{d,+1}`$ on the states of the system. As $`V_{g,j}^{d,+1}`$ is unitary, one has $`e^{iH_{g,j}^dt_j}=V_{g,j}^{d,+1}`$. The presence of the indices $`d,g`$ shows the dependence of $`H_{g,j}^d`$ on the $`W_{g,d}`$.
Efficient implementation means that the time $`t_j`$ must be short. For microscopic systems this is equivalent to the condition that $`t_j`$ must be less than the decoherence time $`t_{dec}`$. If the Hamiltonian and system are such that $`V_{g,j}^{d,+1}`$ is carried out in a number $`n_j`$ of basic switching steps of duration $`\mathrm{\Delta }`$, then $`n_j=t_j/\mathrm{\Delta }<t_{dec}/\mathrm{\Delta }`$ must hold.
For macroscopic systems the efficiency requirement is different as $`t_{dec}<<\mathrm{\Delta }`$. In this case $`n_j`$ must be polynomial and not exponential in $`L`$. This means that $`n_j=O(L^c)`$ with $`c0`$ and $`c`$ not too large. $`O()`$ means ”of the order of”.
The efficiency requirement is much stricter for microscopic systems than for macroscopic ones. The reason is that for most systems $`t_{dec}`$ is small . This is one reason why quantum computers are so hard to implement compared to macroscopic computers. However, the requirement that $`n_j`$ be polynomial in $`L`$ would also apply to any microscopic system for which $`t_{dec}/\mathrm{\Delta }`$ is very large, (e.g. $`t_{dec}`$ is several hours or even longer).
The above is rather general in that it assumes that for each $`j`$ there is a distinct Hamiltonian $`H_{g,j}^d`$ to implement $`V_{g,j}^{d,+1}`$. However for many systems all the $`V_{g,j}^{d,+1}`$ may be implemented by just one Hamiltonian $`H_g^d`$ with the different values of $`j`$ expressed by different states of some ancillary systems.
The requirement of efficient implementation is the reason that the $`V_{g,j}^{d,+1}`$ are defined separately for each $`j`$ rather than defining them from $`V_{g,1}^{d,+1}`$ by $`V_{g,j}^{d,+1}=(V_{g,1}^{d,+1})^{k^{j1}}`$. Here $`V_{g,1}^{d,+1}`$ corresponds to the successor operation $`\mathrm{"}+1\mathrm{"}`$ in axiomatic arithmetic . The exponential dependence on $`j`$ shown by this equation shows that if efficient implementation were required just for $`V_{g,1}^{d,+1}`$, then carrying out of the $`V_{g,j}^{d,+1}`$ is not efficient as exponentially many repetitions of the procedure for $`V_{g,1}^{d,+1}`$ would be required.
For many physical systems, efficient implementation of the $`V_{g,j}^{d,+1}`$ can be carried out by shifting the procedure for implementation of $`V_{g,1}^{d,+1}`$ along path $`g`$ in $`A`$ until a component system in the state $`|g(j)`$ is encountered. At this point implementation of $`V_{g,1}^{d,+1}`$ is started.
Efficient implementability for the basic arithmetic operations also implies that there exist Hamiltonians $`H_{g,d}^+`$ and $`H_{g,d}^\times `$ that efficiently carry out $`+_{g,d}`$ and $`\times _{g,d}`$. Since the definitions of $`+`$ and $`\times `$ are given in terms of the $`V_j^{+1}`$,(Eq. 8 and the Appendix), it follows that if the $`V_{g,j}^{d,+1}`$ can be efficiently implemented, so can $`+_{g,d}`$ and $`\times _{g,d}`$. For microscopic systems the fact that the times $`t_+,t_\times `$ required for these implementations are greater than those for the $`V_{g,j}^{d,+1}`$ means that the values of $`L`$ for which $`t_+<t_{dec}`$ and $`t_\times <t_{dec}`$ may be less than those possible for just the $`V_{g,j}^{d,+1}`$.
Another aspect of the efficient implementability condition is that the thermodynamic resources required to implement $`V_{g,j}^{d,+1}`$ must be polynomial and not exponential in $`j`$. This takes account of the fact that all computations occur in a noisy environment and one must spend thermodynamic resources to protect the system from errors. This is especially the case for quantum computation for which entanglements of states that develop as the computation progresses must be protected from decoherence . Methods of protecting these states include the use of quantum error correction codes and possibly generation and use of EPR pairs . These considerations are another reason why it is important to minimize the time required to implement $`V_{g,j}^{d,+1}`$.
There are many physical systems where the resources needed to implement $`V_{g,j}^{d,+1}`$ (other than those involved in the shift) are either independent of $`j`$ or are at most polynomial in $`L`$. The needed resources do not depend exponentially on $`j`$ or $`L`$. These systems satisfy the requirement of efficient implementability. There are others that do not. Consider, for example, a 1-D lattice of systems where the intensity of environmental interference and noise grows exponentially with $`j`$. Here the thermodynamic resources needed to protect the system from decoherence, etc., would grow exponentially with $`j`$. Another simpler type of system that would be excluded would be a row of isolated harmonic oscillator potentials each containing a single spinless particle. The proposed two qubit states are the ground and first excited states in the well. However the spring constants of the wells depend exponentially on $`j`$. For example the spring constant $`p(j+1)`$ of the $`j+1st`$ well is related to that for the $`jth`$ well by $`p(j+1)=kp(j)`$.
For networks of quantum gates efficient implementation of the basic arithmetic operations has two components. The number of quantum gates (or steps) in the network must be polynomial in $`L`$, as in , and the resources needed to implement individual quantum gates must be polynomial in the locations of the individual systems addressed by each gate. In the physical models described above, this second requirement is not satisfied as resources needed to implement a quantum gate between the $`jth`$ and $`j^{}th`$ qubits depend exponentially on $`j`$ and $`j^{}`$. The fact that one would not build such models or could not build such models for large $`L`$ is not relevant here.
The condition of efficient implementability also places restrictions on the values of $`k`$ allowed for $`kary`$ representations. In general values of $`k`$ are used that are quite small (e.g. $`k=2,k=10`$, etc.). Except for special cases, $`k=1`$ (unary) representations are excluded as arithmetic operations are exponentially hard. Also the value of $`k`$ cannot be too large. One reason is that there are physical limitations on the amount of information that can be reliably stored and distinguished per unit space time volume . Also the requirement of efficient implementation enters in that for large $`k`$ (e.g. $`k=10^6`$), even a simple process such as adding two single digit numbers becomes quite lengthy.
## VI Is the Model $`^{arith}`$ Necessary?
The preceding was based on first constructing a purely mathematical Hilbert space model $`^{arith}`$for mdoular arithmetic and then using this to construct a physical model on a space $`^{phy}`$ that has the same tensor product structure as $`^{arith}`$. The question arises if the purely mathematical model based on $`^{arith}`$ is necessary. Can one go directly from the axioms of modular arithmetic to physical models without the use of the model based on $`^{arith}`$?
In general this is possible as any structure, physical or mathematical, that satisfies the axioms is acceptable. However, the intermediate mathematical models serve as a useful reference point for discussions. This is clear from the literature in which much use is made of such a model. For instance any reference to product qubit states $`|\underset{¯}{0},|0110110\mathrm{}`$, etc. and linear superpositions of these states is implicitly using a model based on $`^{arith}`$.
Another point, already noted, is that the axioms of arithmetic, modular or not, make no mention of efficient implementability. Models based on unary representations are just as valid as are any others. This is true even if additional axioms are added giving the properties of the $`V_j^{+1}`$ operators.
This raises the following questions: Suppose one starts with an arbitrary quantum system with states in a space $`^{phy}`$ whose tensor product structure (if any) is unknown. Can operators, indexed by values in a set of physical parameters for the system, be defined with properties such that they satisfy the axioms of modular arithmetic? As will be seen in the following, this seems possible. If one also requires that the operators and those for the basic arithmetic operations be efficiently implementable, does it follow that $`^{phy}`$ must have a tensor product structure based on the defined operators and their properties? At present, the answer is not known.
To be specific, the interest is in constructing a model of arithmetic $`modk^L`$ directly on the state space $`^{phy}`$ of a quantum system where $`^{phy}`$ has an arbitrary tensor product structure. A set $`A`$ of $`L`$ operators $`V_a`$ on $`^{phy}`$ indexed by the physical parameters $`aϵA`$ is required to have properties that are necessary conditions for $`^{phy}`$ to have the tensor product structure suitable for length $`L`$ $`kary`$ representations of numbers. These properties are,
1. Each $`V_a`$ is a cyclic shift.
2. The $`V_a`$ all commute with one another.
3. For each $`aϵA`$, if $`(V_a)^k1`$ there is a unique $`a^{}a`$ such that $`(V_a)^k=V_a^{}`$.
4. For each $`a^{}`$ , if there is an $`aa^{}`$ such that $`(V_a)^k=V_a^{}`$, then $`a`$ is unique.
5. There is just one $`a`$ for which $`(V_a)^k=1`$.
6. For just one $`a`$ there are no $`a^{}`$ such that $`(V_a^{})^k=V_a`$.
The properties reflect those possessed by the $`V_j^{+1}`$, note especially Eq. 4. Properties 3-6 can be used to establish a numbering of the label set $`A`$ with the maximum and minimum labels given by properties 5 and 6. The commutativity and cyclic shift properties give the existence of of a set $``$ of pairwise orthogonal subspaces of states such that for each $`a`$ and each subspace $`\beta `$ in $``$, $`V_a\beta `$ is in $``$ and is orthogonal to $`\beta `$. In the special case that the subspaces in $``$ are one dimensional, the subspaces $`\beta `$ in $``$ correspond to pairwise orthogonal states $`|\beta `$ such that for each $`|\beta `$ in $``$, $`V_a|\beta `$ and $`|\beta `$ are orthogonal.
One can use property 3 along with iterations $`(V_a)^h`$ for $`h=0,1,\mathrm{}k1`$ for each $`a`$ to generate a cyclic ordering or numbering of the states in $``$ and show that the set contains $`k^L`$ states. However none of this is sufficient to select a state as the zero state. This must be done by making an arbitrary choice.
There is a unique state $`|\beta _{\underset{¯}{0}}`$ in $``$ which is the zero state.
Based on this choice one can associate with each string of numbers, $`n_L,n_{L1},\mathrm{}n_{\mathrm{}},\mathrm{},n_2,n_1=\underset{¯}{n}`$ with $`0n_{\mathrm{}}k1`$ for each $`\mathrm{}`$ a unique state $`|\beta _{\underset{¯}{n}}`$. The association is given by
$$|\beta _{\underset{¯}{n}}=\underset{\mathrm{}=1}{\overset{L}{}}(V_a_{\mathrm{}})^n_{\mathrm{}}|\beta _{\underset{¯}{0}}.$$
where the properties of the $`V_a`$ show that the states $`|\beta _{\underset{¯}{n}}`$ for different number strings $`\underset{¯}{n}`$ are orthogonal.
The above can also be used to define addition as in Eq. 8 and show that $`|\beta _{\underset{¯}{0}}`$ is the additive identity. This and use of the discussion in Section III suggests that these operators and the associated states do satisfy the axioms of arithmetic $`modk^L`$. However examples can be constructed to show that it is very unlikely that the existence of operators with these properties are sufficient conditions for $`^{phy}`$ to have a tensor product structure suitable for $`kary`$ representations of length $`L`$. If one adds the additional requirement that these operators be efficiently implementable, then it is an open question if all these conditions are sufficient to require that $`^{phy}`$ has a tensor product structure suitable for $`kary`$ representations of length $`L`$.
## VII Discussion
Several points about the work done here should be noted. The state descriptions of composite quantum systems used in this paper have not taken account of whether or not the component systems are distinguishable by properties other than those explicitly shown in the states. This is based on the consideration that the only properties used by a quantum algorithm are those expressed explicitly in the states and operators representing the basic arithmetic operations. For indistinguishable systems, it is suspected that taking account of their bosonic or fermionic nature, as has been done elsewhere , will not change the results obtained. However, this must be investigated.
The condition of efficient implementation of the basic arithmetic operations is the main restrictive condition on states of quantum systems that represent numbers. As noted it excludes $`k=1`$ and large $`k`$. It also greatly restricts which unitary operators from $`^{arith}`$ to $`^{phy}`$ are allowed. To see this note that any unitary operator $`U`$, tensor product preserving or not, from $`^{arith}`$ to $`^{phy}`$ gives a model of the axioms of modular arithmetic on $`^{phy}`$. The numbers are represented by the states $`U|\underset{¯}{s}`$ and the basic operators by $`UV_j^{+1}U^{}`$ and $`(UU)+(U^{}U^{})`$ and similarly for $`\times `$. However most of these $`U`$ can be excluded because the corresponding basic operators on $`^{phy}`$ are not efficiently implementable. Also for most $`U`$ there is no way to physically prepare the states $`U|\underset{¯}{s}`$. This is the main reason for the restriction that $`U`$ be tensor product preserving with the form of $`W_{g,d}`$.
Unfortunately there is no way to define exactly which $`U`$ operators are allowed and which are not. The reason is that there is no way to precisely define the meaning of physical realizability. One needs an hypothesis for physical realizability equivalent to the Church-Turing Hypothesis for computable functions. Earlier attempts to characterize realizable physical procedures as collections of instructions , or state preparation and observation proceedures as instruction booklets or programs for robots have not been generally accepted. This problem also arises in describing exactly the class of tasks that a quantum robot can carry out.
Another aspect of the representation of numbers by quantum states is that the sets of numbers $`1,\mathrm{},L`$ and $`0,\mathrm{},k1`$ have been used to describe $`kary`$ representations of numbers of length $`L`$ by quantum states. For example numbers in either of these sets are used to describe the $`V_j^{+1}`$ operations. Also the definitions of $`+`$ and $`\times `$ were given in terms of numbers of iterations of $`V_j^{+1}`$ and $`+`$ respectively.
Two components of this should be noted. One is that the role of these numbers is limited to the dynamical implementation of the $`V_{g,j}^{d,+1},+_{g,d}`$, and $`\times _{g,d}`$. For example, any method based on a Hamiltonian $`H_g^d`$ that implements $`V_{g,j}^{d,+1}`$ as a translation of a procedure for implementing $`V_{g,1}^{d,+1}`$ by $`j`$ sites along $`g`$ requires motion along $`g`$ until the site $`g(j)`$ is reached. This can be done by repeated subtraction of $`1`$ from $`j`$, interleaved with motion of some system, such as a head or quantum robot , along $`g`$ until $`g(j)`$ is reached. Also the ”carry $`1`$” operation, which is part of $`V_{g,j}^{d,+1}`$ means that motion along the remaining $`Lj`$ elements of path $`g`$ must be built into $`H_g^d`$.
Similar arguments apply for the efficient carrying out of the $`+_{g,d}`$ operation as this requires up to $`k`$ iterations of $`V_{g,j}^{d,+1}`$ for each $`j`$. One method of implementation requires interleaving the implementation of a procedure for $`V_{g,j}^{+1}`$ with subtractions of $`1`$ from a state $`|\underset{¯}{s}_j`$, Eq. 8, until $`|\underset{¯}{0}_j`$ is obtained.
Implementation of these operations by quantum systems means that numbers up to $`L`$ and $`k`$ must also be represented by quantum states of systems. These systems can either be mobile and part of the head or fixed external systems. Thus the arguments and conditions already discussed apply to these representations too.
The other component is that the magnitudes of the numbers represented by the states of systems that are part of the dynamics are exponentially smaller than those represented by the system on which the dynamics is acting. States of a composite quantum system satisfying the conditions for $`kary`$ number representations of length $`L`$, represent the first $`k^L`$ numbers. Numbers appearing in the dynamics range up to $`k`$ and $`L=\mathrm{log}_kk^L`$. This exponential decrease is a consequence of the requirement of efficient implementability of arithmetic operations.
The conditions discussed in this paper, including the requirement of efficient physical implementability, also apply to the quantum states of ancillary systems that are used to implement the dynamics of an algorithm. This is evident in any algorithm which interleaves evaluation of some numerical function with carrying out an action until a specified function value is reached. For instance, implementation of the $`V_{g,j}^{d,+1}`$, e.g. by use of a head or quantum robot with an on board quantum computer , would require a quantum computer with at least $`O([\mathrm{log}_m(L)]+1)`$ qubytes for an $`mary`$ representation of numbers up to $`L`$. ($`[]`$ denotes the largest integer in.) Here the dynamics that carries out these operations is subject to all the requirements described so far. It is also part of the dynamics for implementing $`V_{g,j}^{d,+1}`$.
These considerations suggest that it may not be possible to describe the representation of numbers by states of a composite quantum system without the use of states of other systems already assumed to represent numbers. These states are part of the dynamics of the basic arithmetic operations.
Whether this is true or not is a question for the future. However, if this impossibility is the case, one is helped by the fact that the number of states needed to represent numbers in the dynamics is exponentially smaller than the number of states representing numbers of the composite system on which the dynamics acts.
Finally it should be noted that much of the discussion, including the efficient implementability condition, which has been applied to microscopic quantum systems, also applies to macroscopic quantum systems. In this case $`t_{dec}t_{sw}`$ so the limitation that the number of steps is $`<t_{dec}/t_{sw}`$ is not applicable. Instead efficient implementation means that there exists a dynamics such that the number of steps needed to carry out arithmetic operations is polynomial in $`L`$. Also the states of the system used to represent numbers are those that are stabilized by the interactions with the environment, the ”pointer states” . The fact that these conditions are much less onerous than the limitations on microscopic systems is shown by the widespread use of macroscopic computers and counting devices and timers.
In conclusion it is reemphasized that this work is one approach to making explicit the assumptions and conditions involved in the representation of natural numbers by states of quantum systems. It is based on separating the mathematical concept of numbers, as models of a set of axioms, from the physical concept of efficient implementabiliy of the basic arithmetic operations described by the axioms. Whether this approach will turn out to be a good one or not depends on future work.
## Acknowledgements
Discussions with Murray Peshkin on several points of this paper were much appreciated. This work is supported by the U.S. Department of Energy, Nuclear Physics Division, under contract W-31-109-ENG-38.
## Appendix: Definition of $`\times `$
The goal is to define a unitary times operator according to Eq. 9 based on efficient iteration of the $`+`$ operator. To this end define $`Q_j(2,3)`$ for $`j=1,\mathrm{},L`$ as operators on the second and third product states that convert $`|\underset{¯}{s},\underset{¯}{w},\underset{¯}{w0^{j1}},\underset{¯}{z}`$ to $`|\underset{¯}{s},\underset{¯}{w},\underset{¯}{w0^j},\underset{¯}{z}`$. It has the effect of multiplying $`|\underset{¯}{w0^j}`$ by $`k`$. An efficient reversible implementation of this, acting on the state $`|\underset{¯}{s},\underset{¯}{w},\underset{¯}{y},\underset{¯}{z}`$ is obtained by subtraction, $`modk`$, of the $`Lj+1st`$ component qubyte state of $`|\underset{¯}{w}`$ from the $`Lth`$ component state of $`|\underset{¯}{y}`$, shifting all the elements of $`|\underset{¯}{y}`$ by one site and putting the result of the subtraction at the newly opened first site. This works because, if $`|\underset{¯}{y}=|\underset{¯}{w0^{j1}}`$, then $`|\underset{¯}{y}_L=|\underset{¯}{w}_{Lj+1}`$. The result, $`|\underset{¯}{0}_L`$, of the subtraction is moved to the first site of $`|\underset{¯}{y}`$ after the shift. One has
$$Q_j(2,3)|\underset{¯}{s},\underset{¯}{w},\underset{¯}{y},\underset{¯}{z}=|\underset{¯}{s},\underset{¯}{w},\underset{¯}{y^{}},\underset{¯}{z}$$
(20)
where $`|\underset{¯}{y_{j+1}^{}}=|\underset{¯}{y}_j`$ for $`1jL1`$ and $`|\underset{¯}{y_1^{}}=|\underset{¯}{y}_L|\underset{¯}{w}_{Lj+1}`$. Here $``$ denotes subtraction $`modk`$. Note that $`Q_j(2,3)`$ is unitary.
The operator $`\times `$ is defined from the $`Q_j(2,3)`$ and $`+`$ by
$`\times |\underset{¯}{s},\underset{¯}{w},\underset{¯}{y},\underset{¯}{z}=Q_L(2,3)(+_{3,4})^{s_L}Q_{L1}(2,3)(+_{3,4})^{s_{L1}}`$
$`\mathrm{},(+_{3,4})^{s_2}Q_1(2,3)(+_{3,4})^{s_1}+_{2,3}|\underset{¯}{s},\underset{¯}{w},\underset{¯}{y},\underset{¯}{z}`$
Here $`+_{m,n}`$ carries out the action defined in Eq. 8 on the $`mth`$ and $`nth`$ product state. The $`mth`$ state remains unchanged in this action. $`s_h`$ is the number $`\underset{¯}{s}(h)`$ in the state component $`|\underset{¯}{s}(h),h`$ of $`|\underset{¯}{s}`$. Note that since each operator in the righthand product of the equation is unitary, so is $`\times `$.
To see that $`\times `$ as defined above does carry out the intended multiplication operation on initial states of the form $`|\underset{¯}{s},\underset{¯}{w},\underset{¯}{0},\underset{¯}{0}`$ one carries out the action of the $`2L+1`$ operators shown above. The steps give
$`|\underset{¯}{s},\underset{¯}{w},\underset{¯}{0},\underset{¯}{0}\begin{array}{c}+_{2,3}\hfill \\ \hfill \end{array}|\underset{¯}{s},\underset{¯}{w},\underset{¯}{w},\underset{¯}{0}\begin{array}{c}(+_{3,4})^{s_1}\\ \end{array}|\underset{¯}{s},\underset{¯}{w},\underset{¯}{w},\underset{¯}{s_1w}`$ (25)
$`\begin{array}{c}Q_1(2,3)\\ \end{array}|\underset{¯}{s},\underset{¯}{w},\underset{¯}{w0},\underset{¯}{s_1w}\begin{array}{c}(+_{3,4})^{s_2}\\ \end{array}|\underset{¯}{s},\underset{¯}{w},\underset{¯}{w0},\underset{¯}{s_1t+s_2t0}`$ (30)
$`\mathrm{}\begin{array}{c}Q_L(2,3)\\ \end{array}|\underset{¯}{s},\underset{¯}{w},\underset{¯}{0},\underset{¯}{s_1w+s_2t0+\mathrm{}+s_Lt0^{L1}}`$ (33)
Note that $`Q_L(2,3)`$ acting on $`|,\underset{¯}{w},\underset{¯}{w0^{L1}},`$ gives $`|,\underset{¯}{w},\underset{¯}{0},`$ in accordance with Eq. 9 as $`|\underset{¯}{w0^L}=|\underset{¯}{0}`$. Here $`|\underset{¯}{s_1w}`$ denotes $`s_1`$ iterations of adding $`|\underset{¯}{w}`$ to $`|\underset{¯}{0}`$; also $`s_jw0^{j1}`$ denotes the result of $`s_j`$ additions of $`|\underset{¯}{w0^{j1}}`$ to the $`4th`$ product state.
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# Sunyaev-Zel’dovich constraints from black hole-seeded proto-galaxies
## 1 Introduction
Most theories of hierarchical structure formation are based on the study of the evolution of density perturbations under their own gravity. A density fluctuation, which represents an over- (or under-) density with respect to the mean matter distribution, contains both baryonic and dark matter (DM). The baryonic component sinks into the gravitational potential of the DM halo. It collapses and cools, resulting in star formation. In these scenarios, after the gravitational collapse of the DM halo, stars are assumed to be the first objects to form. A structure will thus end up as an emitting object after virialisation has occurred.
An alternative picture involves the formation of a super-massive black hole (BH) that powers the central regions of galaxies Lynden-Bell (1969). Numerous studies have been performed that relate the quasar luminosity function to galaxy formation scenarios by assuming that the formation of quasars (i.e., BH) in the potential well of the DM halos constitutes one of the phases in the galaxy formation process Efstathiou & Rees (1988); Haehnelt & Rees (1993); Nusser & Silk (1993); Haiman & Loeb (1997). Recent observations even suggest that a super-massive BH may be present in the centres of all galaxies with spheroidal components Kormendy & Richstone (1995).
Several authors have looked at several consequences of the presence of massive BHs on galaxy formation and evolution Haiman & Loeb (1997); Natarajan et al. (1998); Silk & Rees (1998). In this paper, we investigate the cosmological implications of such an alternative scenario for the Cosmic Microwave Background (CMB) anisotropies and spectral distortions. More specifically, we study the effects of the outflows, driven by the BH activity, on the gas within the seeded proto-galaxy. In fact, the outflow expands and shock-heats the ambient medium (proto-galactic gas), and then interacts with the inter-galactic medium (IGM). Three regimes of interest may be considered: 1) the high density region of the proto-galaxy, 2) the low density IGM and 3) the thin compressed layer (four times denser than the IGM) induced by the front shock. The second and third regimes give results very similar to those computed in Aghanim et al. (1996). We thus focus on the first regime, i.e. the localised effects of the BH-driven shock on the gas within the seeded proto-galaxy. This shock-heated gas will Compton scatter the CMB photons and induce spectral distortions and temperature anisotropies through the so-called Sunyaev-Zel’dovich effect Sunyaev & Zel’dovich (1980). The thermal SZ effect depresses the CMB brightness in the Rayleigh-Jeans region and increases it above a frequency of about 219 GHz. Its amplitude represents the integral along the line of sight of the electron pressure. It is proportional to the electron density $`n_e`$ and is characterised by the Compton parameter $`y`$ defined by:
$$y=2_0^R\frac{kT}{m_ec^2}\sigma _Tn_e(l)𝑑l,$$
(1)
where $`T`$ is the temperature of the gas, $`R`$ the physical size of the structure, $`m_e`$ the electron mass, $`k`$ the Boltzmann constant, $`c`$ the speed of light and $`\sigma _T`$ the Thomson cross section. An additional secondary anisotropy arises due to the first-order Doppler effect of the CMB photons when they scatter on a structure moving with respect to the Hubble flow, with radial peculiar velocity $`v_r`$. This interaction is called the SZ kinetic effect. It generates an anisotropy, with no specific spectral signature, whose amplitude is given by:
$$\frac{\delta T}{T}=\frac{v_r}{c}\times \left(2_0^R\sigma _Tn_e(l)𝑑l\right).$$
(2)
Previous work on galaxy formation has evaluated the global distortion of a population of galaxies in the virialised regime. The global Compton parameter was found to be much smaller than the constraints set by the FIRAS instrument (Far InfraRed Absolute Spectro-photometer) on board COBE (COsmic Background Explorer) Fixsen et al. (1996) on the global SZ distortion of the universe. By contrast, our model focuses on a regime in which proto-galaxies undergo a BH formation phase that induces larger distortions. The paper is organised as follows: in §2, we model the shock in an individual structure and give its physical characteristics (size and temperature). In §3, we compute the predicted number density of primordial galaxies, using the Press & Schechter (1974) mass function. In §4, we generalise the description of the shock to the whole population of proto-galaxies and we simulate maps of the induced secondary anisotropies. We estimate this contribution to the CMB anisotropies. We also compare our predicted global $`y`$ parameter to the COBE-FIRAS value and derive constraints on the model. Conclusions are given in the last section.
## 2 Modelling the shock
In the galaxy formation canonical scenario, a galaxy forms in the gravitational potential well of a DM halo of mass $`M_{halo}`$. Following Silk & Rees (1998), we consider that each forming galaxy, which is fully described by the mass $`M_{halo}`$, hosts a super-massive BH. The fraction of seeded proto-galaxies is thus 100%. For the sake of simplicity, we assume that all proto-galaxies are spheroids. A more refined and accurate description would take into account morphological segregation Balland et al. (1998). However, these effects do not introduce significant differences. We assume that the BH radiates during its lifetime $`t_{BH}`$ a fraction $`ϵ_E=0.1\mathrm{\hspace{0.17em}0.2}`$ Natarajan & Sigurdsson (1999) of its Eddington luminosity ($`L_{BH}=ϵ_EL_{edd}`$). The latter is fixed by the BH mass, $`M_{BH}`$:
$$L_{edd}=\frac{4\pi M_{BH}m_pGc}{\sigma _T},$$
(3)
where $`m_p`$ is the proton mass and $`G`$ the gravitational constant. The mass of the BH should be directly related to the mass of the proto-galaxy it seeds, and thus to the fraction of baryonic matter locked up in the spheroid in which the central BH collapses. We assume a simple relation of proportionality ($`M_{BH}=ϵ_{BH}M_{sph}`$) between the mass of the BH and the mass of the spheroid $`M_{sph}`$. Haehnelt & Rees (1993) give $`ϵ_{BH}10^3`$ and they assume that $`ϵ_{BH}`$ declines with mass. However, observations of galactic nuclei Kormendy et al. (1997); Magorrian et al. (1998) indicate that $`ϵ_{BH}`$ is relatively constant. Here, we take the value $`ϵ_{BH}=2.10^3`$ as advocated by Magorrian et al. (1998) and Silk & Rees (1998). We assume that the medium is instantaneously ionised, which is very likely due to the intense radiation emitted by the central BH Voit (1996). We therefore study the propagation of a strong shock driven by a mechanical energy which represents some fraction of the BH luminosity. The first analytic solutions of spherically symmetric explosions were given by Taylor (1950) and Sedov (1959) and applied to several astrophysical problems such as stellar winds and supernovae explosions Ikeuchi et al. (1983); Bertschinger (1986); Koo & McKee (1992a, b); Voit (1996). Following Voit (1996), we study the expansion of the shock and its effects on the gas within the seeded proto-galaxy. We characterise the shock in an expanding universe by a set of physical properties (its radius, velocity and temperature). Adapting Voit’s solutions to our study, the shock at a redshift $`z<z_{}`$, where $`z_{}`$ is the redshift at which the BH switches on, has a radius $`R_s`$ which reads as:
$$R_s\left(\frac{\pi GE_0}{H_0^4\delta \mathrm{\Omega }_0^3}\right)^{1/5}\frac{(1+\mathrm{\Omega }_0z_{})^{1/5}}{(1+z_{})^{2/5}}\times $$
$$\left[1\left(\frac{1+\mathrm{\Omega }_0z}{1+\mathrm{\Omega }_0z_{}}\right)^{1/2}\right]^{2/5}(1+z)^1.$$
(4)
It propagates at a velocity
$$v_s=\frac{2}{5}\left[\frac{\pi E_0H_0G}{3\delta \mathrm{\Omega }_b}\frac{\mathrm{\Omega }_0^2(1+z_{})^3}{(1+\mathrm{\Omega }_0z_{})^{3/2}}\right]^{1/5}\times $$
$$\left[1\left(\frac{1+\mathrm{\Omega }_0z}{1+\mathrm{\Omega }_0z_{}}\right)^{1/2}\right]^{3/5}\left(\frac{1+z}{1+z_{}}\right).$$
(5)
In the two previous equations, $`E_0=ϵL_{BH}t_{BH}`$ represents the mechanical energy which drives the shock. Very little is known about the fraction $`ϵ`$ which can be on the order of 0.5 or even higher Natarajan & Sigurdsson (1999). For the lifetime of the BH in the bright phase, we choose the recent value derived by Haiman & Loeb (1997): $`t_{BH}=10^6`$ years (our results, however, are not very sensitive to the exact value of $`t_{BH}`$). $`\delta `$ is the mean over-density of the proto-galaxy (assumed to be spherical), $`H_0`$ is the Hubble constant, $`\mathrm{\Omega }_0`$ is the density parameter and $`\mathrm{\Omega }_b`$ is the baryon density in the universe. Zero subscripts denote present day values. In the following and throughout the paper, we use $`h=H_0/(100\text{km/Mpc/s})=0.5`$ and $`\mathrm{\Omega }_b=0.06`$ Walker et al. (1991) and give the results as a function of $`\mathrm{\Omega }_0`$, for $`\mathrm{\Omega }_0=1`$ and $`\mathrm{\Omega }_0=0.3`$.
Following Natarajan & Sigurdsson (1999), we assume that the BH radiates 10% of its Eddington luminosity, i.e. $`ϵ_E=0.1`$ and that half of the corresponding energy is mechanical. Furthermore, we consider that galaxies form at sufficiently high redshifts ($`>5`$) so that only hydrogen and helium are present in substantial amounts. We compute both the size and velocity of the shock and find that the host proto-galaxy is always embedded in the shocked region.
### 2.1 Cooling mechanisms
The temperature of the shocked medium can be directly derived from the equipartition of energy:
$$Tv_s^2\mu m_p/k.$$
(6)
Here, $`\mu =0.6`$ is the mean molecular weight of a plasma with primordial abundances. The galactic matter can be shock-heated up to very high temperatures. For the most massive galaxies, we find that the matter is heated up to a few $`10^8`$ K, a value comparable to the temperature of the intracluster medium in galaxy clusters. For these temperatures and redshifts ($`>5`$), the main cooling process at play is bremsstrahlung. Therefore, the shocked gas looses heat with a cooling rate given by a temperature-dependent cooling function $`\mathrm{\Lambda }(T)`$. We have compared three different cooling functions Bertschinger (1986); Koo & McKee (1992a); Voit (1996) given in the literature for temperatures between a few $`10^5`$ and $`10^8`$ K. We find that the condition required for efficient cooling is always satisfied in our redshift range. This means that the cooling time, given by $`t_{cool}=5kT/n_H\mathrm{\Lambda }(T)`$ with $`n_H`$ the hydrogen number density, is always smaller than the age of Universe. Cooling is thus very efficient in our picture. The comparison between the three cooling functions shows that they all give essentially the same final temperature after cooling. We follow Bertschinger (1986) and we find that the shocked matter cools down to $`T10^5`$ K at $`z5`$.
## 3 Number counts of primordial galaxies
In order to quantify the global effect of the formation of primordial galaxies on the CMB, we apply our formalism to a synthetic population of galaxies with masses $`10^9M_{\mathrm{}}M10^{12}M_{\mathrm{}}`$.
We first assume that the galaxy number density traces, within a linear bias, the abundance of collapsed DM halos, as predicted by the Press–Schechter (PS) mass function Press & Schechter (1974). We use an initial power-law spectrum with an effective spectral index $`n=2`$ on galaxy scales. We express the amplitude of primordial matter fluctuations in terms of the rms variance in spheres of 8$`h^1`$ Mpc, $`\sigma _8=0.6`$ (as cluster-normalised, e.g. Viana & Liddle (1996), which corresponds to a bias factor $`b1.67`$. For $`\mathrm{\Omega }=1`$, the set of parameters corresponds to the “standard” biased cold DM model, which does fit neither small- and large-scale velocities Vittorio et al. (1986) nor COBE normalisation. However, we take it as a study case for the computations, our second model is the low density cosmological model with $`\mathrm{\Omega }=0.3`$. In our picture (no stars are formed yet), the spheroid is gaseous. Its mass is related to the mass of the DM halo via $`M_{sph}=\frac{1}{3}M_{halo}`$. The mass and luminosity of the central BH and thus the predicted SZ distortions are therefore inferred from $`M_{sph}`$ ($`M_{BH}=ϵ_{BH}M_{sph}`$, and $`ϵ_{BH}=\mathrm{2.\hspace{0.17em}10}^3`$; cf §2).
To compute the kinetic SZ term of a population of proto-galaxies, we need an estimate of their peculiar velocities with respect to the reference frame. As suggested by numerical simulations Bahcall et al. (1994); Moscardini et al. (1996), we assume that velocities follow a Gaussian distribution. The peculiar velocity of each proto-galaxy is drawn from a Gaussian which is completely defined by its rms value $`\sigma _v`$. In the range of redshifts we have adopted, the structures are in the linear regime, so that $`\sigma _v(z)=\sigma _0f(z)`$, where the redshift dependence of the velocities is given by $`f(z)`$Peebles (1980, 1993) as a function of the cosmological parameters. In this equation, $`\sigma _0`$ is the present-day rms peculiar velocity. It is related to the mass variance on mass scale $`M`$, $`\sigma (M)=(1.19\mathrm{\Omega }_0)^{(n+3)/6}\sigma _8M^{(n+3)/6}`$ Mathiesen & Evrard (1998), where $`n`$ is the index of the power spectrum. The rms velocity can thus be computed for each mass scale.
## 4 Results and discussion
The shock-heated gas within proto-galaxies interacts with the CMB photons through the SZ effect (thermal and kinetic). These interactions generate secondary temperature anisotropies and spectral distortions. We simulate maps of the secondary anisotropies generated by a population of seeded proto-galaxies formed between redshift 5 and 10. The maps have a resolution of about $`0.2`$ arcseconds to resolve the galaxies and contain $`600\times 600`$ pixels. The number of sources of mass $`M`$ at redshift $`z`$ is derived from the PS mass function. Their positions are drawn at random in the map. The $`y`$ and $`\delta T/T`$ profiles for, respectively, the thermal and the kinetic effects are directly derived from the integration of the gas profile $`n_e(R)`$ along the line of sight (Eqs. 1 and 2) assuming spherical symmetry.
Similarly to the case of galaxy clusters, we assume that in the early stages of formation the gas settles into a hydrostatic equilibrium within the DM potential. A universal density profile is motivated by Navarro et al. (1996). However, the gas profile may be softer than that of the DM, and moreover the existence of a central cusp is “unobserved” Kravtsov et al. (1998).We thus conservatively adopt the following parametrised profile for the gas distribution:
$$n_e(R)=n_0\left[1+\left(\frac{R}{R_c}\right)^2\right]^\alpha ,$$
(7)
where $`n_0`$ is the central density. $`\alpha `$ is left as a free parameter describing the steepness of the profile, whereas $`R_c`$ is identified with a core radius as in galaxy clusters. On cluster scales, $`R_c`$ is typically 10 to 30 times smaller than the cluster virial radius $`R_{vir}`$. In our model we introduce the parameter $`p=R_{vir}/R_c`$ which we vary, similarly to clusters, between 10 and 30. The central density $`n_0`$ can be derived from the gas mass of the proto-galaxy using the following equation:
$$M_{sph}\left(\frac{\mathrm{\Omega }_b}{\mathrm{\Omega }_0}\right)=m_p\mu _0^{R_{vir}}n_e(R)\mathrm{\hspace{0.17em}4}\pi R^2𝑑R,$$
(8)
where the virial radius of the structure is given by:
$$R_{vir}=\frac{(GM)^{1/3}}{(3\pi H_0)^{2/3}}\frac{1}{1+z_{coll}},$$
(9)
for a critical universe. It is fixed solely by the mass and the collapse redshift $`z_{coll}`$.
We will give the results for a flat model with no cosmological constant ($`\mathrm{\Omega }_0=1`$) and an open model ($`\mathrm{\Omega }_0=0.3`$). Varying the cosmological parameters will vary the number of proto-galaxies along the line of sight as well as their peculiar velocities. It will also modify their physical properties, i.e. the size and velocity of the shock and thus the gas temperature. The two cosmological models represent the upper and lower bounds between which all other cosmological models involving a non-zero cosmological constant fall.
### 4.1 Compton distortion
The CMB photons, scattering off the electrons of the ionised hot gas, induce a spectral distortion whose amplitude is given by Eq. 1. The FIRAS experiment has measured the mean Compton parameter resulting from all the interactions undergone by the photons. The result is $`\overline{y}_{FIRAS}=\mathrm{1.5\hspace{0.17em}10}^5`$ Fixsen et al. (1996). This stringent observational limit incorporates the (negligible) contribution of the rather cold intergalactic medium and that of all other extragalactic signals. Among these signals, there is the contribution of the hot ionised gas in galaxy clusters. The global distortion induced by intra-cluster gas has been computed De Luca et al. (1995); Barbosa et al. (1996), and found to be of the order of a few $`10^6`$. In addition to galaxy clusters, one has to take into account the contribution of the proto-galaxy population in terms of the overall Compton distortion, $`\overline{y}_{PG}`$, induced by the scattering of CMB photons on the shock-heated gas.
Based on simulated maps, we predict $`\overline{y}_{PG}`$ and we compare it to the limit set by COBE-FIRAS. Among all the parameters of the model, there are four major quantities that substantially affect the predictions of the mean Compton parameter. Two of them, $`\alpha `$ and $`p`$, are related to the gas distribution (Eq. 7). The two others are the fraction $`f`$ of BH-seeded proto-galaxies and the BH-to-spheroid mass ratio $`ϵ_{BH}`$. We compare our predicted overall distortion to the COBE-FIRAS limit and look for the combinations of parameters for which our predictions fit the observations. This allows us to constrain the assumptions of our model.
For $`\alpha =1/2`$, $`f=100`$%, $`ϵ_{BH}=\mathrm{2.\hspace{0.17em}10}^3`$ and both cosmological models, we find $`\overline{y}_{PG}10^4`$ which exceeds the observational value. In order for our prediction to be reconciled with the COBE-FIRAS limit, $`f`$ must be only a few percent. This constraint on $`f`$ strongly violates the actual observations Magorrian et al. (1998); Richstone et al. (1998). $`\alpha =1/2`$ is thus excluded by the limit on the global distortion whatever value we choose for $`p=R_{vir}/R_c`$.
For $`\alpha =1`$ (i.e. an isothermal profile), an $`\mathrm{\Omega }_0=1`$ model, $`f=100`$% and $`ϵ_{BH}=\mathrm{2.\hspace{0.17em}10}^3`$, we find $`\overline{y}_{PG}>\overline{y}_{FIRAS}`$ whatever we adopt for $`p`$. The fraction $`f`$ must be smaller than 75% for the prediction to be compatible with the observational limit. Again, this fraction is significantly smaller than the 95% advocated by Magorrian et al. (1998). Such a constraint could rule out the isothermal profile. However, up to now, $`ϵ_{BH}`$ was assumed to be constant and equal to $`\mathrm{2.\hspace{0.17em}10}^3`$. If we now use the lower limit of Magorrian et al. (1998), that is $`ϵ_{BH}=10^3`$, together with $`f=95`$% or higher there is only a marginal agreement for $`p30`$ between the predicted and measured distortions. In the open model case ($`\mathrm{\Omega }_0=0.3`$), we find approximately the same results. For $`\overline{y}_{PG}`$ to be compatible with $`\overline{y}_{FIRAS}`$ if $`p=10`$, the fraction of BH seeded proto-galaxies should be smaller than 80% if $`ϵ_{BH}=\mathrm{2.\hspace{0.17em}10}^3`$ or $`f<95\%`$ if $`ϵ_{BH}=10^3`$. The predictions for $`p=30`$ agree with observations in all the cases.
For $`\alpha =3/2`$ (i.e. the gas profile approximates a King profile) and for both cosmological models, we find $`\overline{y}_{PG}`$ of about $`10^6`$ to a few $`10^6`$, a prediction compatible with $`\overline{y}_{FIRAS}`$. This result remains valid for all values of $`p`$ between 10 and 30 including the highest boundaries of $`f`$ and $`ϵ_{BH}`$, respectively, 100% and $`\mathrm{2.\hspace{0.17em}10}^3`$.
### 4.2 Predicting the angular power spectrum
We choose the set of parameters associated with the isothermal profile which agrees with the COBE-FIRAS limit: $`\alpha =1`$, $`p=30`$, $`f=95`$% and $`ϵ_{BH}=\mathrm{2.\hspace{0.17em}10}^3`$. Within this context, we predict the upper limit on the contribution to secondary temperature anisotropies induced by the SZ, thermal and kinetic effects, of the proto-galactic gas. We express this contribution in terms of an angular power spectrum plotted in figure 1 together with the main other well-known secondary anisotropies.
At very small scales ($`l`$ a few $`10^5`$) corresponding to galactic scales, the kinetic SZ contribution of the shock-heated gas (Fig. 1, thick solid line for $`\mathrm{\Omega }_0=1`$ and thick dashed line for $`\mathrm{\Omega }_0=0.3`$) is very large. It is interesting to note the good agreement between our results and those obtained by Peebles & Juszkiewicz (1998) for the scattering of the CMB photons by the cloudy proto-galactic plasma. The power spectrum of the kinetic SZ anisotropies for the $`\mathrm{\Omega }_0=0.3`$ model is significantly larger than the $`\mathrm{\Omega }_0=1`$ model. This is mainly due to the higher number of sources per unit comoving volume in open models. In all other flat cosmological models involving a non-zero cosmological constant, the power spectrum will lie between the two curves. The expected power spectrum due to the thermal SZ effect is not plotted in this figure. It is more than one order of magnitude smaller than the kinetic effect contribution. This is due to the efficiency of bremsstrahlung which lowers the temperature down to a few $`10^5`$ K.
We compare the contribution of the proto-galaxies due to their SZ kinetic effect to the major sources of secondary temperature anisotropies. In each case, we choose the most extreme cases for the comparison with our upper limit prediction. The power spectra displayed in figure 1 are taken from the literature. The dotted line represents the upper limit of the contribution of the inhomogeneous reionisation as computed by Aghanim et al. (1996) for a quasar lifetime of $`10^7`$ yrs. The dot-dashed line represents the Rees-Sciama effect Rees & Sciama (1968) taken from Seljak (1996) ($`\sigma _8=1`$, $`\mathrm{\Omega }_0h=0.25`$). The dashed line represents the galaxy cluster contribution due to kinetic SZ effect from Aghanim et al. (1998) with a cut-off mass of $`10^{14}M_{\mathrm{}}`$. The triple-dot-dashed line represents the Vishniac-Ostriker effect Ostriker & Vishniac (1986); Vishniac (1987) computed by Hu & White (1996) with a total reionisation occurring at $`z_i=10`$. Finally, the solid thin line represents the power spectrum of the primary CMB anisotropies in a standard CDM model computed using the CMBFAST code Seljak & Zaldarriaga (1996). The primary CMB anisotropies dominate at all scales larger than the damping around 5 arcminutes. At intermediate scales, several effects take place among which the inhomogeneous reionisation, the Ostriker-Vishniac and the SZ effect. In figure 1, we do not plot the power spectra of the thermal SZ effect of galaxy clusters. It is about one order of magnitude larger than the kinetic SZ effect. At very small scales, the anisotropies are totally dominated by the proto-galactic contribution.
## 5 Conclusions
Previous studies on galaxy formation have evaluated the global distortion of a population of galaxies in the virialised regime. In these studies, the global Compton parameter was found to be very small, and smaller than the COBE-FIRAS value. In contrast, our model focuses on a regime in which proto-galaxies undergo a BH formation phase. During this phase, the proto-galactic matter is shock-heated up to a few $`10^8`$ K and cools down to $`10^5`$ K. CMB photons undergo inverse Compton scattering on the heated gas. In addition, galaxy peculiar motions induce temperature anisotropies through the SZ kinetic effect. We have estimated the global Compton parameter due to a population of proto-galaxies and the expected power spectrum of the induced secondary anisotropies. We find that there are four main parameters that control our model: the fraction $`f`$ of BH-seeded proto-galaxies, the fraction $`ϵ_{BH}`$ of the spheroid mass in the BH, the steepness of the density profile $`\alpha `$ and the gas core radius $`p=R_{vir}/R_c`$. The comparison between our predictions and the COBE-FIRAS observation constrains these parameters. Given the observed fraction of seeded galaxies, $`f=95`$%, our results put rather strong constraints on the density profile and on $`ϵ_{BH}`$. Indeed, our predictions agree with the observations whatever $`p`$ if the density profile is an approximation to a King profile. On the contrary, if the density profile is isothermal, then the core radius must be at least 30 times smaller than the virial radius and the BH-to-spheroid mass ratio has to be small, of the order of $`10^3`$. The computations in the two extreme cosmological models show that the global Compton parameter due to proto-galaxies is not very sensitive to $`\mathrm{\Omega }_0`$.
We compare the power spectra of the different contributions to the temperature anisotropies. Our results show that the SZ effect of the very early shock-heated proto-galaxies could constitute the major source of CMB distortions on very small scales (arcsecond and sub-arcsecond scales). The anisotropies are likely to be detected and measured by future long baseline interferometers such as ALMA. The shock heating is likely to contribute to the re-heating of the proto-galactic gas, which plays a role in galaxy formation theory. Blanchard et al. (1992) used preheating to modify the galaxy luminosity function, suppressing and finally delaying dwarf galaxy formation. We do not take into account this effect in our model, therefore, our results should be taken as an upper limit to the proto-galaxy contribution in terms of secondary anisotropies.
###### Acknowledgements.
This work was initiated while C.B. was at the Center for Particle Astrophysics and while N.A. was at the Astronomy Department of the University of Berkeley (U.S.A.). N.A. acknowledges partial funding from the Centre National d’Études Spatiales. During this work, J.S. was partially supported by a Blaise Pascal professorship. The authors thank O. Forni, G. Lagache and J.-L. Puget for discussions and careful reading of the manuscript. They also wish to thank R. Juszkiewicz for his comments that helped to improve the final version of the paper.
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# The 2dF QSO Redshift Survey.
## 1. Introduction
The observational aim of the 2dF QSO Survey is to use the new AAT 2dF fibre-optic coupler to obtain redshifts for 25000 B$`<`$20.85, 0$`<`$z$`<`$3 QSO’s in two $`75\times 5`$deg<sup>2</sup> strips of sky in the RA ranges 21h50 - 03h15 at $`\delta `$ = -30 and 09h50 - 14h50 at $`\delta `$ = +00. The 2dF instrument allows spectra for up to 400 QSO candidates to be obtained simultaneously. The input catalogue is based on APM UBR magnitudes for $`7.5\times 10^6`$ stars to B=20.85 on 30 UKST fields. The final QSO catalogue will be an order of magnitude bigger than previous complete QSO surveys. The prime scientific aims of the 2dF QSO survey are:
1. To determine the QSO clustering power spectrum, P(k), in the range of spatial scales, $`0<r<1000h^1Mpc`$.
2. To measure $`\mathrm{\Omega }_\mathrm{\Lambda }`$ from geometric distortions in clustering.
3. To trace the evolution of QSO clustering in the range, 0$`<`$z$`<`$3, to obtain new limits on $`\mathrm{\Omega }_0`$ and QSO bias.
Other aims include determining the evolution of the QSO LF evolution to z=3, cross-correlating the QSO’s with 2dF Galaxy Survey groups to measure $`\mathrm{\Omega }_0`$ via gravitational lensing (cf. Croom & Shanks, 1999a) and constraining $`\mathrm{\Omega }_\mathrm{\Lambda }`$ by finding the sky density of close (6-20<sup>′′</sup>) lensed QSO pairs.
Previous QSO clustering results have generally been based on the following complete QSO surveys - the Durham/AAT survey of Boyle et al (1988) comprising 392 B$`<`$21 QSO’s, the CFHT survey of Crampton et al (1989) with 215 B$`<`$20.5 QSO’s, the ESO survey of Zitelli et al (1992) with 28 B$`<`$22 QSO’s, the LBQS survey of Hewett et al (1995) with 1053 B$`<`$18.8 QSO’s and the ESO survey of LaFranca et al (1998) with 300 B$`<`$20.5 QSO’s. These have produced results on QSO clustering at small (r$`<10h^1`$ Mpc) scales which are consistent with an r<sup>-1.8</sup> power-law with amplitude $`r_06h^1`$Mpc which shows no evolution in comoving coordinates (Shanks et al 1987, Andreani & Cristiani, 1992, Shanks & Boyle, 1994, Georgantopoulos & Shanks, 1994, Croom & Shanks, 1996). At larger scales (10$`<r<1000h^1`$Mpc), no clustering has previously been detected in $`\xi _{qq}`$.
## 2. 2dF QSO Redshift Survey Status and Current Results.
We now have redshifts for $``$6000 QSO’s where 5600 QSO’s have been observed using 2dF itself. A further 400 QSO’s selected from the input catalogue have been observed on different telescopes for associated projects. Observations have included bright 17$`<B<`$18.25 QSO’s using UK Schmidt Telescope FLAIR fibre coupler, radio-loud QSO’s identified in the NRAO VLA Sky Survey (NVSS) and observed at Keck and finally 30 QSO’s in close pairs ($`<20^{\prime \prime }`$) from the ANU 2.3-m telescope. This makes the 2dF survey already the biggest, complete QSO survey by a factor of $``$6.
From these observations, we know that 53.5% of candidates are QSO’s which means there will be 26000 QSO’s in final survey. The QSO number count, n(B), has been found to be in good agreement with previous surveys. The QSO redshift distribution, n(z), extends to z$``$3 because of our multi-colour UBR selection. In Fig. 1 we show the current Northern and Southern redshift cone plots from the 2dF survey. The rectangles shows the sky distribution of the fields that have already been observed.
We have detected 8 close ($`<20^{\prime \prime }`$) QSO pairs. Only one is a candidate gravitational lens; the separation is $`16^{\prime \prime }`$ (see Fig. 6 of Croom et al., 1999).
The new 2dF results for the QSO Luminosity Function continue to be consistent with Pure Luminosity Evolution models throughout the range 0.35$`<`$z$`<`$2.3 (Boyle et al.,1988,1991,1999a). The luminosity function based on the current $``$6000 QSOs of the 2dF survey combined with $``$1000 LBQS QSOs are shown in Fig. 2. The large sample size makes it possible to define the QSO Luminosity Function in much smaller redshift bins than used previously and the accuracy of pure luminosity evolution as a description of QSO evolution is clear. Indeed, this accuracy is so high that it has prompted a new investigation of pure-luminosity evolution as a physical, as well as a phenomenological, model for QSO evolution. (Done & Shanks, 1999, in prep.)
The most interesting individual QSO that has been discovered from the 2dF QSO survey is UN J1025-0040, a unique, post-starburst radio QSO at z=0.634, identified in the 2dF-NVSS catalogue and followed up spectroscopically at the Keck (Brotherton et al., 1999). As well as broad emission lines, the spectrum also shows strong Balmer absorption lines indicative of a post-starburst galaxy. The starburst component of the spectrum at M<sub>B</sub>=-24.7 dominates the AGN continuum spectrum by $``$ 2mag. This 2dF-NVSS collaboration has previously also uncovered a new class of radio-loud BAL QSO’s (Brotherton et al., 1998) and clearly has great potential for further exciting discoveries.
Finally, we present a preliminary 2dF QSO correlation function from our most complete subset of 4115 QSO’s. We have taken into account the current incompleteness of the 2dF survey as best we can; however, this process is complicated by the fact that many observed areas still have overlapping 2dF ‘tiles’ as yet unobserved. We show the preliminary correlation function in Fig. 3(a) as a log-log plot and in Fig. 3(b) as a log-linear plot. As can be seen the correlation function is consistent with being a -1.8 power-law with $`r_04h^1`$Mpc and in good agreement with previous results (eg Croom and Shanks, 1996). The QSO correlation function thus appears to be remarkably similar to the correlation function for local, optically selected galaxies and continues to show no evolution in comoving coordinates. This behaviour is consistent with evolution due to gravitational clustering either in a low $`\mathrm{\Omega }_0`$ model or in a biased, $`\mathrm{\Omega }_0=1`$ model. The correlation function is consistently positive out to about 20h<sup>-1</sup>Mpc at 3$`\sigma `$ and out to about 50h<sup>-1</sup>Mpc at 1$`\sigma `$, thus giving a hint of power extending to larger scales. At $`r>50h^1`$Mpc the errors in $`\xi `$ are now as low as $`\pm `$0.016. There is also a suggestion of anti-correlation at r$``$100h<sup>-1</sup>Mpc but still only at 1.5$`\sigma `$ significance. At all scales in the range 100$`<`$r$`<`$1000h<sup>-1</sup>Mpc the correlation function is within 1$`\sigma `$ of zero.
## 3. Future Results.
We have used mock catalogues drawn from simulations using the Zeldovich approximation to demonstrate the accuracy of P(k) estimates at large scales that will be obtainable from the full survey (Croom et al., in prep.). The results from this simulation show that our estimates of the power spectrum can recover the input spectrum up to scales of several hundreds of Mpc (k$``$0.01hMpc<sup>-1</sup>) where the power spectrum is expected to turn over to its primordial form (see Boyle et al., 1999b).
We are now using the Hubble Volume N-body simulation from the Virgo Consortium with of order 10<sup>9</sup> particles to produce mock QSO survey catalogues. We have ‘light cone’ output for a model with $`\mathrm{\Omega }_\mathrm{\Lambda }`$=0.7, $`\mathrm{\Omega }_0`$=0.3 and $`\sigma _8`$=1.0 to test the potential of the geometric test of Alcock & Paczynski (1979) for $`\mathrm{\Omega }_\mathrm{\Lambda }`$ directly in a biased $`\mathrm{\Lambda }`$CDM model (Hoyle et al., 1999 in prep.). We shall be using these simulations to see how robust the test is against different models for the bias. Constraints on these bias models come from the QSO-galaxy cross-corrrelation function (Ellingson et al., 1991, Smith et al., 1995, Croom & Shanks, 1999b) which suggest radio-quiet QSO’s inhabit similar environments to average galaxies out to z=1.5. We are now extending QSO environment tests out to z$``$2 using the INT Wide Field Camera at B$`<`$26 and the AAT Taurus Tunable Filter (Croom & Shanks, in prep.) and these results will further constrain possible models of QSO bias at high redshift.
## 4. Conclusions
1. The 2dF QSO Survey has obtained redshifts for 6000/26000 B$`<`$20.85, z$`<`$3 QSOs which means it is already the biggest complete QSO survey by a factor of $``$6.
2. Previous surveys detected QSO clustering at 4$`\sigma `$ level for r$`<`$10h<sup>-1</sup>Mpc where the QSO clustering appears stable when measured in comoving coordinates, suggesting either a low $`\mathrm{\Omega }_0`$ model or a biased, $`\mathrm{\Omega }_0`$=1 model. Previous surveys detected no significant clustering in $`\xi _{qq}`$ at larger scales.
3. The 2dF QSO survey confirms the accuracy of the PLE model of Boyle et al (1991) for the evolution QSO Luminosity Function for 0$`<`$z$`<`$2.2.
4. The 2dF QSO survey has already detected many individually interesting objects, including a post-starburst QSO and several close QSO pairs.
5. A preliminary 2dF QSO correlation function based on 4115 QSOs shows results consistent with previous surveys and consistent with the correlation function of local, optically selected galaxies but with a hint of extended power to $``$50h<sup>-1</sup>Mpc.
6. Mock catalogues from the Hubble Volume suggest the 2dF QSO survey will be able to determine QSO power spectrum out to $`>`$1000h<sup>-1</sup>Mpc and so detect the expected turnover to the primordial, n=1, slope.
7. The potential of 2dF QSO survey to constrain $`\mathrm{\Omega }_\mathrm{\Lambda }`$ from both geometric distortions and gravitational lensing is being tested in the Hubble Volume mock catalogues.
### Acknowledgments.
We thank Fiona Hoyle, Carlton Baugh and Adrian Jenkins for allowing us to use the results from their analyses of the mock catalogues from the Hubble Volume simulation. The Hubble Volume is provided courtesy of the Virgo Consortium.
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# Search for New Particles Decaying to 𝑡𝑡̄ in 𝑝𝑝̄ Collisions at √𝑠=1.8 TeV
## Abstract
We use 106 $`\mathrm{pb}^1`$ of data collected with the Collider Detector at Fermilab to search for narrow-width, vector particles decaying to a top and an anti-top quark. Model independent upper limits on the cross section for narrow, vector resonances decaying to $`t\overline{t}`$ are presented. At the 95% confidence level, we exclude the existence of a leptophobic $`Z^{}`$ boson in a model of topcolor-assisted technicolor with mass $`M_Z^{}`$ $`<`$ 480 $`\mathrm{GeV}/c^2`$ for natural width $`\mathrm{\Gamma }`$ = 0.012 $`M_Z^{}`$, and $`M_Z^{}`$ $`<`$ 780 $`\mathrm{GeV}/c^2`$ for $`\mathrm{\Gamma }`$ = 0.04 $`M_Z^{}`$.
PACS numbers: 14.65.Ha, 13.85.Ni, 13.85.Qk
In this letter, we present a model-independent search for narrow, vector resonances decaying to $`t\overline{t}`$. This search is sensitive to, for example, a $`Z^{}`$ predicted by topcolor-assisted technicolor. This model anticipates that the explanation of spontaneous electroweak symmetry breaking is related to the observed fermion masses, and that the large value of the top quark mass suggests the introduction of new strong dynamics into the standard model. It accounts for the large top quark mass by predicting the existence of a residual global symmetry SU(3)$`\times `$U(1) at energies below 1 TeV. The SU(3) results in the generation of topgluons which we have searched for previously in the $`b\overline{b}`$ channel. The U(1) gives the $`Z^{}`$ we search for here. In one model, the $`Z^{}`$ decays exclusively to quarks (leptophobic) resulting in a large cross section for $`t\overline{t}`$.
With the $`z`$-axis defined along the proton beam, the Collider Detector at Fermilab (CDF) coordinate system defines $`\varphi `$ as the azimuthal angle in the transverse plane, $`\theta `$ as the polar angle, and pseudorapidity $`\eta `$ as $`\mathrm{ln}(\mathrm{tan}\frac{\theta }{2})`$. Tracking chambers, immersed in a 1.4-Tesla solenoidal magnetic field, are used for the detection of charged particles and the measurement of their momenta. The precision track reconstruction of the silicon microstrip vertex detector (SVX), located immediately outside the beampipe, is used for the detection of displaced secondary vertices resulting from $`b`$-quark decays. Outside the SVX is the vertex time projection chamber (VTX) which provides further tracking information for $`|\eta |`$ $``$ $`3.25`$. Both the SVX and VTX are housed within the central tracking chamber (CTC), a wire drift chamber used to measure charged particle momenta. Electromagnetic and hadronic calorimeters, located beyond the CTC and superconducting solenoid, measure energy in segmented $`\eta `$-$`\varphi `$ towers out to $`|\eta |`$ $`<`$ 4.2. Drift chambers used for muon detection reside outside the calorimetry. A more detailed description of the CDF detector can be found elsewhere.
Standard model $`t\overline{t}`$ production in $`p\overline{p}`$ collisions at a center of mass energy of $`\sqrt{s}=1.8`$ TeV is dominated by $`q\overline{q}`$ annihilation, while $``$10% is attributable to gluon-gluon fusion. Once a $`t\overline{t}`$ pair is produced, each of the top quarks is expected to decay almost exclusively to $`Wb`$. The search presented here focuses on the $`t\overline{t}`$ event topology in which one $`W`$ boson decays hadronically while the other decays to an electron or muon and its corresponding neutrino. The fragmentation of the $`b`$-quarks, as well as the hadronic daughters of the $`W`$ boson, form jets. Accordingly, $`t\overline{t}`$ candidates in this “lepton + jets” channel are characterized by a single lepton, missing transverse energy, $`\mathrm{}\mathrm{E}_\mathrm{T}`$ , due to the undetected neutrino, and at least four jets. Furthermore, a jet resulting from a $`b`$-quark can be identified (or “tagged”) as such by the reconstruction of a secondary vertex from the $`b`$ hadron decay using the SVX, or by using the soft lepton tagging (SLT) algorithm to find an additional lepton from a semileptonic $`b`$ decay.
Like the top quark mass measurement, events included in our measurement of the $`t\overline{t}`$ invariant mass spectrum must first contain a lepton candidate in the central detector region ($`|\eta |`$ $`<`$ 1.0). This lepton is required to be either an isolated electron with transverse energy ($`\mathrm{E}_\mathrm{T}`$) in excess of 20 $`\mathrm{GeV}`$ or an isolated muon with transverse momentum ($`\mathrm{P}_\mathrm{T}`$) in excess of 20 $`\mathrm{GeV}/\mathrm{c}`$. Events must also include at least 20 $`\mathrm{GeV}`$ of $`\mathrm{}\mathrm{E}_\mathrm{T}`$, attributable to the presence of a neutrino, as well as at least four jets with $`|\eta |`$ $`<`$ 2.0 and raw $`\mathrm{E}_\mathrm{T}`$ $`>`$ 15 $`\mathrm{GeV}`$. Raw jet energies are the values which result from clustering signals in the calorimeter towers before any offline jet corrections are applied. To increase the acceptance rate for $`t\overline{t}`$ events, the requirements for the fourth jet are relaxed such that the raw $`\mathrm{E}_\mathrm{T}`$ must only be greater than 8 $`\mathrm{GeV}`$ with $`|\eta |`$ $`<`$ 2.4 in events where at least one of the leading three jets is tagged by the SVX or SLT algorithms. All jets in this analysis are formed as clusters of calorimeter towers within cones of fixed radius $`\mathrm{\Delta }R\sqrt{(\mathrm{\Delta }\eta )^2+(\mathrm{\Delta }\varphi )^2}=0.4`$. In 106 $`\mathrm{pb}^1`$ of data, we observe 83 events which satisfy these requirements.
This method builds upon the techniques developed for the top quark mass measurement by fitting each event to the hypothesis of $`t\overline{t}`$ production followed by decay in the lepton+jets channel:
The four-momenta of these 13 objects fully describe a $`t\overline{t}`$ event. The three-momenta of the charged lepton and four jets are measured directly. To compute the energies of these objects, the $`b`$ and $`\overline{b}`$ quark masses are taken to be 5 $`\mathrm{GeV}/c^2`$, the $`q`$ and $`\overline{q}^{}`$ masses are taken to be 0.5 $`\mathrm{GeV}/c^2`$, and the charged lepton mass is assigned according to its identification as either an electron or a muon. The components of transverse momentum for the recoiling system, $`\xi `$, are measured directly from extra jets in the event and unclustered energy deposits that are not included in lepton or jet energies. The transverse momentum components of the neutrino are computed by requiring that the total $`\mathrm{E}_\mathrm{T}`$ in the event sums to zero. While the neutrino is assumed to be massless, its longitudinal momentum is a free parameter in the kinematic fit in which the $`q\overline{q}^{}`$ and $`\mathrm{}\nu `$ invariant masses are constrained to equal the $`W`$ boson mass. We perform a kinematic fit to the production and decay of the $`t\overline{t}`$ pair as described by the decay chain shown above. This fitting procedure, which depends on the minimization of a $`\chi ^2`$ expression, allows the lepton energy, the jet energies and the unclustered energy to vary within their respective uncertainties. The fitted results for these values determine the $`t`$ and $`\overline{t}`$ four-momentum, from which the $`t\overline{t}`$ invariant mass ($`M_{t\overline{t}}`$) can be computed. To improve the $`M_{t\overline{t}}`$ resolution, we also constrain the two $`Wb`$ invariant masses to 175 GeV/c<sup>2</sup>, in agreement with the most recent measurement of the top quark mass. We use only the four highest $`\mathrm{E}_\mathrm{T}`$ jets, leading to 12 combinations for assigning jets to the $`b`$, $`\overline{b}`$, and hadronic $`W`$ daughters. However, because we measure only the transverse component of the total energy, thereby determining $`\mathrm{}\mathrm{E}_\mathrm{T}`$, a two-fold ambiguity in the longitudinal component of the neutrino momentum results in 24 combinations. We further require that jets that are SVX or SLT-tagged be assigned to $`b`$-quarks, thereby reducing the number of combinations.
Electron energies and muon momenta are measured with the calorimeter and tracking chambers, respectively. A set of generic jet corrections is applied to the energies of all the jets in an event to account for absolute energy scale calibration, contributions from the underlying event and multiple interactions, as well as energy losses in cracks between detector components and outside the clustering cone. These corrections are determined from a combination of Monte Carlo simulations and data. The four leading jets in a $`t\overline{t}`$ event undergo an additional energy correction that depends on the type of parton that they are assumed to be in the fit: a light quark, a hadronically decaying $`b`$ quark, or a $`b`$ quark that decayed semileptonically. These parton-specific corrections account for (a) the differences in the expected $`\mathrm{P}_\mathrm{T}`$ distributions of jets from $`t\overline{t}`$ and the shape which was assumed to derive the generic jet corrections mentioned above, and (b) the energy losses from semileptonic $`b`$ and $`c`$-hadron decays. These corrections were derived from a study of $`t\overline{t}`$ events generated with the herwig Monte Carlo program.
Using Monte Carlo simulations of signal and background events, we explored several event selection criteria in an attempt to optimize our discovery potential. Of the 24 possibilities for each event, we select the $`M_{t\overline{t}}`$ value which corresponds to the configuration with the lowest $`\chi ^2`$. To reduce the probability of selecting configurations with incorrect parton assignments which tend to yield artificially low values of $`M_{t\overline{t}}`$, we refit each event after releasing the constraint that the $`Wb`$ invariant mass be equal to 175 GeV/c<sup>2</sup> and demand that the fit for this particular configuration return a value for the top quark mass between 150 GeV/c<sup>2</sup> and 200 GeV/c<sup>2</sup>. To further reduce incorrect combinations and to increase discovery potential for a new particle decaying to $`t\overline{t}`$, we apply a $`\chi ^2`$ cut. For narrow width $`t\overline{t}`$ resonances, simulation predicts that the width of the $`M_{t\overline{t}}`$ spectrum is $`6`$% of the resonance mass for cases in which the correct jet configuration is selected. For resonances with a natural width $`\mathrm{\Gamma }`$ that is significantly less than 6% of the nominal mass, the CDF detector resolution will dominate and the resonances will all have approximately the same shape (shown in the inset of Fig. 1 for a mass of 500 $`\mathrm{GeV}/c^2`$). At low $`M_{t\overline{t}}`$, the presence of residual events with incorrect parton assignments is evident in this figure.
The selection criteria described above eliminate an additional 20 events from our data sample and the resulting $`M_{t\overline{t}}`$ spectrum is shown in Fig. 1, along with the expected standard model $`t\overline{t}`$ and QCD $`W`$+jets background shapes normalized to the data. While the non-$`t\overline{t}`$ background is dominated by $`W`$+jets events, it also includes contributions from multijet $`b\overline{b}`$ events with one jet misidentified as a lepton, $`Z`$+jets events, events with a boson pair, and single-top production. However, it has been shown that the vecbos $`W`$+jets shape alone is sufficient in modeling the entire non-$`t\overline{t}`$ background spectrum. For this analysis, the expected non-$`t\overline{t}`$ background prediction of 31.1 $`\pm `$ 8.5 events is calculated as in Ref. , but accounts for differences in selection criteria. We find that the $`M_{t\overline{t}}`$ distribution of 63 data events is consistent with the hypothesis that the spectrum is comprised of standard model $`t\overline{t}`$ production and the predicted rate of non-$`t\overline{t}`$ background events, as shown in Fig. 1.
Because we cannot present evidence for a narrow $`t\overline{t}`$ resonance, we establish upper limits on the production cross-section for a new vector particle, $`X`$, of mass $`M_\text{X}`$ decaying and to $`t\overline{t}`$. For natural widths $`\mathrm{\Gamma }=0.012M_\text{X}`$ and $`\mathrm{\Gamma }=0.04M_\text{X}`$, and for each $`M_\text{X}`$ between 400 $`\mathrm{GeV}/c^2`$ and 1 $`\mathrm{TeV}/c^2`$ in increments of 50 $`\mathrm{GeV}/c^2`$, we perform a binned-likelihood fit of the data. To determine the likelihood function for a given $`M_\text{X}`$ and $`\mathrm{\Gamma }`$, we fit the $`M_{t\overline{t}}`$ spectrum from the data to the expected Monte Carlo shapes for both the $`t\overline{t}`$ and QCD $`W`$+jets background sources as well as a resonance signal $`Xt\overline{t}`$ which we model using $`Z^{}t\overline{t}`$ in pythia.
Our analysis is subject to several sources of systematic uncertainty on the expected shape of background and signal $`M_{t\overline{t}}`$ spectra and/or the signal acceptance rate. Treating these two types of systematic effect separately, we establish the magnitude of each source through a Monte Carlo procedure which quantifies the effect of varying the source of uncertainty by one standard deviation. We determine the uncertainty contributions due to the jet $`\mathrm{E}_\mathrm{T}`$ scale, initial and final state gluon radiation, and the non-$`t\overline{t}`$ background spectrum using methods described in Ref. . The uncertainty in the measurements of the top quark mass and total integrated luminosity are included in our study of systematic effects, as well as the uncertainty due to the choice of parton distribution functions (PDF). The remaining sources of systematic uncertainty considered are all small and include trigger efficiency, lepton identification, tracking efficiency, $`z`$-vertex efficiency, and Monte Carlo statistics. The uncertainties resulting from jet $`\mathrm{E}_\mathrm{T}`$ scale and top quark mass are correlated and we conservatively take this correlation to be 100%.
The percent uncertainty in $`\sigma _\text{X}\text{BR}\{\text{X}t\overline{t}\}`$ is listed in Table I for each of the systematic sources at several different resonance masses. The systematic effect due to uncertainty in the top quark mass ($`M_{top}`$) is dominant at low $`M_\text{X}`$, whereas the effect due to the uncertainty in modeling final state radiation dominates at large $`M_\text{X}`$. To ensure that our estimates are conservative, the systematic uncertainty is taken to be a constant number of pb below the value of $`\sigma _\text{X}\text{BR}\{\text{X}t\overline{t}\}`$ corresponding to the 95% C.L. limit obtained with statistical uncertainties only. That constant is the estimate of the systematic uncertainty at the 95% C.L. limit. Above the same value of $`\sigma _\text{X}\text{BR}\{\text{X}t\overline{t}\}`$, we use a systematic uncertainty that rises with $`\sigma _\text{X}\text{BR}\{\text{X}t\overline{t}\}`$ at the fixed percent rate listed in Table I.
For each resonance mass and width, we convolute the statistical likelihood shape with the Gaussian total systematic uncertainty and extract the 95% C.L. upper limit on $`\sigma _\text{X}\text{BR}\{\text{X}t\overline{t}\}`$ which is listed in Table II and shown in Fig. 2. The systematic uncertainties increase the 95% C.L. upper limit by 27% for $`M_\text{X}`$ = 400 $`\mathrm{GeV}/c^2`$, but only 7% (6%) for $`M_\text{X}`$ = 600 (800) $`\mathrm{GeV}/c^2`$ because statistical uncertainties dominate the likelihood. Also shown in Fig. 2 are the theoretical predictions for cross-section times branching ratio for a leptophobic $`Z^{}`$ with natural width $`\mathrm{\Gamma }`$ = $`0.012M_Z^{}`$ and $`\mathrm{\Gamma }`$ = $`0.04M_Z^{}`$. At 95% confidence, we exclude the existence of a leptophobic topcolor $`Z^{}`$ with mass $`M_Z^{}`$ $`<`$ 480 $`\mathrm{GeV}/c^2`$ for natural width $`\mathrm{\Gamma }`$ = 0.012 $`M_Z^{}`$, and mass $`M_Z^{}`$ $`<`$ 780 $`\mathrm{GeV}/c^2`$ for $`\mathrm{\Gamma }`$ = 0.04 $`M_Z^{}`$. For larger widths, detector resolution will no longer be the dominant factor in determining the $`Z^{}`$ signal shape, so our limits are no longer applicable.
In conclusion, after investigating 106 $`\mathrm{pb}^1`$ of data collected at CDF, we find no evidence for a $`t\overline{t}`$ resonance and establish upper limits on cross-section times branching ratio for narrow resonances. We have used these limits to constrain a model of topcolor assisted technicolor. We thank the Fermilab staff and the technical staffs of the participating institutions for their vital contributions. This work is supported by the U.S. Department of Energy and the National Science Foundation; the Natural Sciences and Engineering Research Council of Canada; the Istituto Nazionale di Fisica Nucleare of Italy; the Ministry of Education, Science and Culture of Japan; the National Science Council of the Republic of China; and the A.P. Sloan Foundation.
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# Berry’s Phase in the Presence of a Dissipative Medium
## 1 Introduction
The existence of geometric phases in quantum systems, ever since its discovery by Berry , has attracted considerable interest both from the theoretical and experimental viewpoints. Several applications of this phenomenon in different areas of physics have also been studied. Comparatively much less work has been devoted to the question of the dynamical evolution of such systems in the presence of a weakly dissipative medium. Basically the available results can be summarized as follows: nonhermitian operators lead to a modification of Berry’s phase, stochastically evolving magnetic fields produce both energy shift and broadening, phenomenological weakly dissipative Liouvillians alter Berry’s phase by introducing an imaginary correction or causing damping and mixing of the density matrix elements. Ellinas et al. obtain their results studying the eigenmatrices of the complete Liouvillian superoperator in a time-independent basis, while Gamliel and Freed, find closed formal expressions for the density matrix representation in the instantaneous frame of the hamiltonian in the adiabatic limit, and in the weak dissipation approximation. However, even in this regime the results in the cases studied can only be extracted through approximations (in addition to the adiabatic approximation) or numerical computation.
In the present work, we consider the celebrated example of a spin 1/2 particle coupled to a slowly evolving magnetic field in the presence of a weakly dissipative medium as represented by a Lindblad form superoperator, incorporated in a convenient and physically motivated frame. In the absence of dissipation, a simple geometrical interpretation of the results emerges in terms of both the Bloch’s vector and the phase vector. The geometrical phase appears as a “delay” or an “advance” in the precession period of the Bloch’s vector with respect to the period dictated by the magnetic field’s frequency. The precession of this vector occurs at a fixed angle with respect to a fixed axis about which the external magnetic field precesses, as usual. We will introduce the dissipation via a semi-group type dynamics, such that dissipation does not alter the precession frequency. However, the $`x`$, $`y`$ and $`z`$ components of Bloch’s vector are altered in different ways by nonunitary effects. In particular, the modulus of Bloch’s vector will shrink. The relation between the reduction of Bloch’s vector modulus and the loss of coherence has been explored by Stodolsky and collaborators. In order to make it explicit, we introduce the linear entropy (or idempotency defect) as a measure of purity loss
$$\delta (t)=1Tr\rho ^2(t)$$
(1)
and write
$$\rho (t)=\frac{1}{2}\left(\begin{array}{cc}1+S_z(t)& S_x(t)iS_y(t)\\ S_x(t)+iS_y(t)& 1S_z(t)\end{array}\right)=\frac{1}{2}(1\mathrm{l}+\stackrel{}{S}(t)\stackrel{}{\sigma }),$$
(2)
where $`\stackrel{}{S}(t)`$ is the Bloch’s vector. We get
$$\delta (t)=\frac{1|\stackrel{}{S}(t)|^2}{2},$$
(3)
which confirms that the shrinking of the modulus of $`\stackrel{}{S}(t)`$ is a measure of coherence loss.
In order to discuss dissipative effects on interference patterns due to Berry’s geometrical phase the natural framework is to consider the evolution of the density matrix, which completely characterizes the interference effect due to the different geometric phases acquired by the eigenvectors of the hamiltonian of the system, even in the absence of dissipation. Of course, this time we add a nonunitary Liouville operator contribution to the dynamics.
In section 2 we present the two level quantum system of a spin 1/2 in the presence of an external magnetic field precessing with constant angular velocity around a fixed axis. The master equation of the model is written in what we call the diagonal frame where Lindblad superoperators are introduced to describe the nonunitary part of the Liouville operator. In subsection 2.1 we consider the adiabatic limit of the quantum system in thermal equilibrium with a reservoir of electromagnetic fields and in subsection 2.2 we consider the case of a dephasing process. In both cases, we treat the coupling of the quantum system with its environment in the weak regime. In section 3 we show how the geometric phases and the dissipation effects acts on the $`z`$ component of the magnetization of the spin 1/2 coupled to a reservoir of electromagnetic fields at thermal equilibrium. In section 4 we summarize our conclusions; finally, in appendix A, we present the non-unitary part of the liouvillians of the models under consideration in the instantaneous basis of the hamiltonian.
## 2 The adiabatic limit of the spin 1/2 model in the weak coupling regime
We consider a spin 1/2 variable (two level model) coupled to a time dependent magnetic field precessing around the z-axis. The unitary contribution for this evolution is given by the hamiltonian
$`H_s(t)`$ $`=`$ $`\mu \stackrel{}{\sigma }\stackrel{}{B}(t)`$ (6)
$`=`$ $`\mu B\left(\begin{array}{cc}\mathrm{cos}(\theta )\hfill & \hfill \mathrm{sin}(\theta )e^{i\omega t}\\ \mathrm{sin}(\theta )e^{i\omega t}\hfill & \hfill \mathrm{cos}(\theta )\end{array}\right),`$
written in the basis of the eigenstates of the $`z`$-component of the spin, where $`B`$ is the norm of the external magnetic field, $`\theta `$ its azimuthal angle, $`\omega `$ the precession frequency and the constant $`\mu =\frac{g\mu _B}{2}`$, being $`g`$ Landé’s factor and $`\mu _B`$ the Bohr magneton. We are using natural units ($`c=\mathrm{}=1`$). For the sake of later calculations, it is convenient to define two unitary transformations: the first one, $`R(\omega ,t)`$, takes us to the rotating frame where the hamiltonian is no longer time dependent; the second one, $`D(B,\theta ,\omega )`$, diagonalizes the effective hamiltonian (time independent) that drives the dynamics of the final matrix representation of the density operator. After the first transformation, $`R(\omega ,t)=e^{\frac{i\omega t}{2}\sigma _z}`$, the density matrix and the hamiltonian read
$$\rho _R(t)=e^{i\frac{\omega t}{2}\sigma _z}\rho (t)e^{i\frac{\omega t}{2}\sigma _z},$$
(7)
and
$$H_R=\mu B(\mathrm{sin}(\theta )\sigma _x+\mathrm{cos}(\theta )\sigma _z).$$
(8)
In analogous manner, after the second transformation we get, in the diagonal frame, the density matrix
$$\rho _D(t)=𝐃^T\rho _R(t)𝐃,$$
(9)
and the effective hamiltonian
$`H_D`$ $`=`$ $`𝐃^T(H_R{\displaystyle \frac{\omega }{2}}\sigma _z)𝐃`$ (12)
$`=`$ $`\left(\begin{array}{cc}\lambda _1\hfill & \hfill 0\\ 0\hfill & \hfill \lambda _1\end{array}\right),`$
where $`\lambda _1=\sqrt{\mu ^2B^2\mathrm{sin}^2(\theta )+(\mu B\mathrm{cos}(\theta )\frac{\omega }{2})^2}`$. The rotation matrix $`𝐃=𝐃^T`$ is equal to
$$𝐃=\sqrt{\frac{1}{2}\frac{1}{2\lambda _1}(\mu B\mathrm{cos}(\theta )\frac{\omega }{2})}\sigma _x+\sqrt{\frac{1}{2}+\frac{1}{2\lambda _1}(\mu B\mathrm{cos}(\theta )\frac{\omega }{2})}\sigma _z.$$
(13)
One possible way to add dissipative contributions to the above dynamics is to include a Lindblad type superoperator in the evolution equation in the diagonal frame
$$\frac{d}{dt}\rho _D\left(t\right)=i[\lambda _1\sigma _z,\rho _D\left(t\right)]+k_D\rho _D(t),$$
(14)
where $`k`$ is the dissipation constant. The weak coupling regime is characterized by the condition $`\frac{k}{\lambda _1}1`$.
Our aim is isolating, in the density matrix, the effects of dissipation on interference due to geometric phases. The representation of the density matrix in the diagonal frame is not very enlightening for this purpose. This is better realized in some basis of the instantaneous eigenvectors of hamiltonian (6), as was done by Gamliel and Freed . We define $`\rho _I(t)`$, the matrix density in a basis of the instantaneous eigenvectors of hamiltonian (6). The relation between $`\rho _I(t)`$ and $`\rho _D(t)`$ is
$$\rho _I(t)=𝐕^{}(t)𝐃\rho _D(t)𝐃𝐕(t),$$
(15)
where the matrix $`𝐕(t)`$ is equal to
$$𝐕(t)=\left(\begin{array}{cc}\mathrm{cos}(\frac{\theta }{2})e^{\frac{i\omega t}{2}}\hfill & \hfill \mathrm{sin}(\frac{\theta }{2})e^{\frac{i\omega t}{2}}\\ \mathrm{sin}(\frac{\theta }{2})e^{\frac{i\omega t}{2}}\hfill & \hfill \mathrm{cos}(\frac{\theta }{2})e^{\frac{i\omega t}{2}}\end{array}\right).$$
(16)
The time evolution of $`\rho _I(t)`$ is given by
$$\frac{d}{dt}\rho _I\left(t\right)=i[\left(\mu B+\frac{\omega }{2}\right)\sigma _z\frac{\omega }{2}\sigma _n(t),\rho _I\left(t\right)]+k_I\rho _I(t)$$
(17a)
where
$$\sigma _n(t)=\left(\begin{array}{cc}\mathrm{cos}(\theta )\hfill & \hfill \mathrm{sin}(\theta )e^{i\omega t}\\ \mathrm{sin}(\theta )e^{i\omega t}\hfill & \hfill \mathrm{cos}(\theta )\end{array}\right)$$
(17b)
and $`_I\rho _I(t)`$ is obtained from $`_D\rho _D(t)`$ through a similarity transformation equivalent to (15).
Before we specialize our discussion to any particular liouvillian, we study the adiabatic limit of eqs.(17) for the coupling constant $`k`$ in the weak regime. In this regime, the matrix $`_I\rho _I(t)`$ is written as a linear superposition of the elements $`\rho _{ij}^I(t)`$. We remind that the density matrix of a two level model must satisfy two conditions: i) $`Tr(\rho _I(t))=1`$ and ii) $`\rho _{21}^I(t)=(\rho _{12}^I(t))^{}`$. As a consequence of those conditions, the density matrix has only two independent elements. We take the elements $`\rho _{11}^I(t)`$ and $`\rho _{12}^I(t)`$ as our two independent entries. The general form for the time equations of those two elements in the weak coupling regime is
$$\frac{d}{dt}\rho _{11}^I(t)=a_{11}(\omega ,k;t)\rho _{11}^I(t)+a_{12}(\omega ,k;t)\rho _{12}^I(t)+a_{13}(\omega ,k;t)\rho _{21}^I(t)+b_1(\omega ,k;t)$$
(18a)
and
$`{\displaystyle \frac{d}{dt}}\rho _{12}^I(t)`$ $`=`$ $`a_{21}(\omega ,k;t)\rho _{11}^I(t)+(2i\mu B+a_{22}(\omega ,k;t))\rho _{12}^I(t)+a_{23}(\omega ,k;t)\rho _{21}^I(t)+`$ (18b)
$`+b_2(\omega ,k;t).`$
We make the change of variables:
$$\rho _{ij}^I(t)e^{i(E_iE_j)t}\stackrel{~}{\rho }_{ij}(t)$$
(19a)
with $`E_1=\mu B`$ and $`E_2=\mu B`$. We may disclose the time scale $`T`$ in the differential equations by introducing the following transformation upon the time parameter:
$$s\frac{t}{T},$$
(19b)
where $`T=\frac{2\pi }{\omega }`$.
With the new variables, eqs.(18) become
$`{\displaystyle \frac{d}{ds}}\stackrel{~}{\rho }_{11}(s)`$ $`=`$ $`Ta_{11}(\omega ,k;s)\stackrel{~}{\rho }_{11}(s)+Ta_{12}(\omega ,k;s)e^{2i\mu BTs}\stackrel{~}{\rho }_{12}(s)+`$ (20a)
$`+`$ $`Ta_{13}(\omega ,k;s)e^{2i\mu BTs}\stackrel{~}{\rho }_{21}(s)+Tb_1(\omega ,k;s),`$
$`{\displaystyle \frac{d}{ds}}\stackrel{~}{\rho }_{12}(s)`$ $`=`$ $`Ta_{21}(\omega ,k;s)e^{2i\mu BTs}\stackrel{~}{\rho }_{11}(s)+Ta_{22}(\omega ,k;s)\stackrel{~}{\rho }_{12}(s)+`$ (20b)
$`+`$ $`Ta_{23}(\omega ,k;s)e^{4i\mu BTs}\stackrel{~}{\rho }_{21}(s)+Tb_2(\omega ,k;s)e^{2i\mu BTs}.`$
We take the differential equation for $`\stackrel{~}{\rho }_{12}(t)`$ to exemplify the discussion of the adiabatic limit of eqs.(20). At this point we will follow closely the references . We point out that the adiabatic approximation is not recovered by an $`\omega `$ expansion of the terms on the r.h.s. of eqs.(20).
The Volterra equation obtained from eq.(20b) is
$`\stackrel{~}{\rho }_{12}(s)`$ $`=`$ $`\stackrel{~}{\rho }_{12}(0)+T{\displaystyle _0^s}a_{22}(\omega ,k;s^{})\stackrel{~}{\rho }_{12}(s^{})𝑑s^{}+T{\displaystyle _0^s}a_{21}(\omega ,k;s^{})e^{2i\mu BTs^{}}\stackrel{~}{\rho }_{11}(s^{})𝑑s^{}+`$ (21)
$`+`$ $`T{\displaystyle _0^s}a_{23}(\omega ,k;s^{})e^{4i\mu BTs^{}}\stackrel{~}{\rho }_{21}(s^{})𝑑s^{}+T{\displaystyle _0^s}b_2(\omega ,k;s^{})e^{2i\mu BTs^{}}𝑑s^{}.`$
In the limit $`T\mathrm{}`$, the Riemann-Lebesgue Theorem gives that
$$\underset{T\mathrm{}}{lim}_0^sF(s^{})e^{i\alpha \mu TBs^{}}𝑑s^{}=0,$$
(22)
if $`F(s^{})`$ is a piece-wise continuous function in the interval \[$`0,s`$\] and $`\alpha `$. As a consequence of this theorem the last three terms on the r.h.s. of eq.(21) vanish in the adiabatic limit ($`\frac{\omega }{\mu B}0T\mathrm{}`$).
Integrating by parts eq.(22) we get
$$_0^sF(s^{})e^{i\alpha \mu TBs^{}}ds^{}=\frac{1}{iT}\frac{1}{\alpha \mu B}[F(s)e^{i\alpha \mu TBs}F(0))]\frac{1}{iT}\frac{1}{\alpha \mu B}_0^s\frac{d}{ds^{}}(F(s^{}))e^{i\alpha \mu TBs^{}}ds^{}.$$
(23)
From eqs.(17) and (18b), the coefficients $`a_{2j}`$ and $`b_2`$ have the general dependence on $`\omega `$ and $`k`$:
$`a_{21}(\omega ,k;s)`$ $`=`$ $`\omega a_{21}^{(0)}(s)+k\left(\stackrel{~}{a}_{21}^{(0)}(s)+{\displaystyle \frac{\omega }{\mu B}}A_{21}(s)+𝒪\left(\left({\displaystyle \frac{\omega }{\mu B}}\right)^2\right)\right),`$ (24a)
$`a_{22}(\omega ,k;s)`$ $`=`$ $`i\omega (1\mathrm{cos}(\theta ))+k\left(\stackrel{~}{a}_{22}^{(0)}(s)+{\displaystyle \frac{\omega }{\mu B}}A_{22}(s)+𝒪\left(\left({\displaystyle \frac{\omega }{\mu B}}\right)^2\right)\right),`$ (24b)
$`a_{23}(\omega ,k;s)`$ $`=`$ $`k\left(\stackrel{~}{a}_{23}^{(0)}(s)+{\displaystyle \frac{\omega }{\mu B}}A_{23}(s)+𝒪\left(\left({\displaystyle \frac{\omega }{\mu B}}\right)^2\right)\right)`$ (24c)
and
$$b_2(\omega ,k;s)=\omega b_2^{(0)}(s)+k\left(\stackrel{~}{b}_2^{(0)}(s)+\frac{\omega }{\mu B}B_2(s)+𝒪\left(\left(\frac{\omega }{\mu B}\right)^2\right)\right).$$
(24d)
The integrands of all integrals on the r.h.s. of eq.(21) which contain the oscillatory function $`e^{i\alpha \mu BTs^{}}`$, with $`\alpha =2`$ or $`4`$, after integration by parts, each of those integrals has its order in T decreased by one unit, and acquires a multiplying constant of value $`\frac{k}{\mu B}`$ or $`\frac{\omega }{\mu B}`$. From the conditions satisfied by the two simultaneous regimes: i) adiabatic limit ($`\frac{\omega }{\mu B}1`$) and ii) weak coupling limit ($`\frac{k}{\mu B}1`$), we can neglect those terms in comparison to the first two terms on the r.h.s. of eq.(21). The two previous inequalities do not impose any constraint to the ratio $`\frac{k}{\omega }`$, though. The differential equation satisfied by $`\stackrel{~}{\rho }_{12}(s)`$ in the adiabatic limit and weak coupling regime is
$$\frac{d}{ds}\stackrel{~}{\rho }_{12}(s)=T\left[i\omega (1\mathrm{cos}(\theta ))+k\stackrel{~}{a}_{22}^{(0)}(s)\right]\stackrel{~}{\rho }_{12}(s).$$
(25)
The term proportional to $`A_{22}(s)`$ was dropped, since it is of higher order in $`\left(\frac{\omega }{\mu B}\right)`$. For $`k\omega `$, the off-diagonal elements of the density matrix vanish before the external magnetic field $`\stackrel{}{B}(t)`$ returns to its configuration at $`t=0`$. The other uninteresting situation from the point of view of the appearance of an imaginary correction to the geometric phase is the condition $`k\omega `$ when the dissipation effects can still be neglected after one period. In this work we discuss the case $`k\omega `$ when both terms on the r.h.s. of eq.(25) contribute to the time evolution of $`\stackrel{~}{\rho }_{12}(s)`$.
By a similar discussion we obtain the time equation of $`\stackrel{~}{\rho }_{11}(s)`$ in the adiabatic and weak coupling regimes,
$$\frac{d}{ds}\stackrel{~}{\rho }_{11}(s)=Tk\stackrel{~}{a}_{11}^{(0)}(s)\stackrel{~}{\rho }_{11}(s)+T\left[\omega b_1^{(0)}(s)+k\stackrel{~}{b}_1^{(0)}(s)\right].$$
(26a)
From eqs.(17) we can affirm that
$`a_{11}(\omega ,k;s)`$ $`=`$ $`k\left(\stackrel{~}{a}_{11}^{(0)}(s)+{\displaystyle \frac{\omega }{\mu B}}A_{11}(s)+𝒪\left(\left({\displaystyle \frac{\omega }{\mu B}}\right)^2\right)\right),`$ (26b)
$`b_1(\omega ,k;s)`$ $`=`$ $`k\left(\stackrel{~}{b}_1^{(0)}(s)+{\displaystyle \frac{\omega }{\mu B}}B_1(s)+𝒪\left(\left({\displaystyle \frac{\omega }{\mu B}}\right)^2\right)\right).`$ (26c)
In order to verify if we can get imaginary phases from eqs. (25) and (26a) due to the coupling of the quantum system to a dissipative medium, we write the density operator at $`t=0`$ in the instantaneous basis of hamiltonian (6)
$$\rho (0)=\underset{j,l=1}{\overset{2}{}}\alpha _{jl}(0)|\varphi _j^0\varphi _l^0|.$$
(27)
We have $`𝐇(0)|\varphi _j^0=E_j(0)|\varphi _j^0`$, $`j=1`$ and $`2`$. Being $`|\varphi _j^0(t)`$ the time evolution of the eigenvector $`|\varphi _j^0`$, we have
$$\rho (t)=\underset{j,l=1}{\overset{2}{}}\alpha _{jl}(t)|\varphi _j^0(t)\varphi _l^0(t)|,$$
(28)
where the time evolution of $`|\varphi _j^0(t)`$ is driven by $`𝐇(t)`$. Differently from Gamliel and Freed, we include the phases coming from the unitary evolution and the geometric phase in the dyadic product $`|\varphi _j^0(t)\varphi _l^0(t)|`$. In the adiabatic approximation, we get
$$|\varphi _j^0(t)=e^{i\gamma _j(t)}e^{iE_j(t)t}|\varphi _j;t,$$
(29)
where $`\gamma _j(t)`$ is the geometric phase acquired by the eigenvector $`|\varphi _j^0`$, $`|\varphi _j;t`$ is the instantaneous eigenvector of $`𝐇(t)`$ ( $`𝐇(t)|\varphi _j;t=E_j(t)|\varphi _j;t`$ ) and $`E_j(t)\frac{1}{t}_0^t𝑑t^{}E_j(t^{})`$.
The density operator at any time, in the adiabatic limit is
$$\rho (t)=\underset{j,l=1}{\overset{2}{}}\alpha _{jl}(t)e^{i(\gamma _j(t)\gamma _l(t))}e^{i(E_j(t)E_l(t))t}|\varphi _j;t\varphi _l;t|.$$
(30)
From eq.(30) we recognize that the phase $`e^{i(\gamma _j(t)\gamma _l(t))}`$ in the element $`\rho _{jl}^I(t)`$ is just the difference of the geometric phases of the instantaneous eigenstates $`|\varphi _j;t`$ and $`|\varphi _l;t`$ in the absence of dissipation.
The dynamics of the coefficients $`\alpha _{jl}(t)`$ is ruled by the nonunitary evolution of the quantum system and it is independent of the particular choice for the instantaneous eigenstates of the hamiltonian, up to a multiplicative constant. In the model under consideration (see hamiltonian (6)), the eigenvalues $`E_j(t)`$, $`j=1`$ and $`2`$, are time-independent. It is simple to get the time equations of $`\alpha _{11}(t)`$ and $`\alpha _{12}(t)`$ from eqs.(26a) and (25), respectively
$`{\displaystyle \frac{d}{dt}}\alpha _{11}(t)`$ $`=`$ $`k\left(\stackrel{~}{a}_{11}^{(0)}(t)\alpha _{11}(t)+\stackrel{~}{b}_1^{(0)}(t)\right),`$ (31a)
$`{\displaystyle \frac{d}{dt}}\alpha _{12}(t)`$ $`=`$ $`k\stackrel{~}{a}_{22}^{(0)}(t)\alpha _{12}(t).`$ (31b)
The constants $`\stackrel{~}{a}_{11}^{(0)}(t)`$, $`\stackrel{~}{a}_{22}^{(0)}(t)`$ and $`\stackrel{~}{b}_1^{(0)}(t)`$ depend on the particular master equation that describes the behaviour of the quantum system interacting with the dissipative medium. In the next sub-sections we consider two particular interactions of the two level model with a reservoir: i) two level model in thermal equilibrium with a reservoir of electromagnetic fields; ii) dephasing process in a two level model.
### 2.1 Adiabatic limit of a two level model in thermal equilibrium
As discussed before, to incorporate the dissipative effects in the two level model we introduce the Lindblad superoperator in the diagonal frame. The master equation of the spin 1/2 model coupled to a reservoir of electromagnetic fields in thermal equilibrium in the diagonal frame is
$`{\displaystyle \frac{d}{dt}}\rho _D\left(t\right)`$ $`=`$ $`i[\lambda _1\sigma _z,\rho _D\left(t\right)]+k\left(\overline{n}+1\right)\left(2\sigma _{}\rho _D\left(t\right)\sigma _+\rho _D\left(t\right)\sigma _+\sigma _{}\sigma _+\sigma _{}\rho _D\left(t\right)\right)+`$ (32)
$`+k\overline{n}\left(2\sigma _+\rho _D\left(t\right)\sigma _{}\rho _D\left(t\right)\sigma _{}\sigma _+\sigma _{}\sigma _+\rho _D\left(t\right)\right),`$
where $`k`$ is the dissipation constant at zero temperature and $`\overline{n}`$ is the average number of excitations of the weakly coupled thermal oscillators at inverse temperature $`\beta `$. An important requirement for the introduction of this Lindblad type superoperator is that it leads asymptotically to a thermal equilibrium. In appendix A we give the master equation of this physical process in the instantaneous basis of hamiltonian (6) for arbitrary value of $`\omega `$.
From the master equation in the instantaneous basis of the hamiltonian we obtain the equations for $`\alpha _{11}(t)`$ and $`\alpha _{12}(t)`$. These equations in the adiabatic and in the weak coupling limits become
$`{\displaystyle \frac{d}{dt}}\alpha _{11}(t)`$ $`=`$ $`2k(1+2\overline{n})\alpha _{11}(t)+2k\overline{n},`$ (33a)
$`{\displaystyle \frac{d}{dt}}\alpha _{12}(t)`$ $`=`$ $`k(1+2\overline{n})\alpha _{12}(t).`$ (33b)
The constant $`k`$ does not come up on the r.h.s. of eqs.(33) due to the time variation of any classical parameter that characterizes the reservoir. From its explicit definition it can not be written as: $`f(t)\dot{g}(t)`$, where $`f(t)`$ and $`g(t)`$ are two regular time dependent functions. The same is true for $`\overline{n}`$. Therefore the imaginary phase $`\chi (t)`$ defined as: $`\alpha _{12}(t)\alpha _{12}(0)e^{i\chi (t)}`$, with
$$\chi (t)=i_0^tk(1+2\overline{n})𝑑t^{}$$
(34)
is not geometric.
The solution of eq.(33a) is
$$\alpha _{11}(t)=\frac{\overline{n}}{1+2\overline{n}}+\left[\alpha _{11}(0)\frac{\overline{n}}{1+2\overline{n}}\right]e^{2k_0^t(1+2\overline{n})𝑑t^{}}.$$
(35)
The exponential decay on the r.h.s. of eq.(35) means that the population of the instantaneous eigenstates of hamiltonian (6) varies in time and consequently the Adiabatic Theorem is not valid for this dissipation mechanism.
Exactly soluble models are always important checks to approximation schemes. Eq.(32) is exactly solved and the solutions are
$`\rho _{11}^D(t)`$ $`=`$ $`{\displaystyle \frac{\overline{n}}{2\overline{n}+1}}\left[1e^{2k(2\overline{n}+1)t}\right]+\rho _{11}^D(0)e^{2k(2\overline{n}+1)t},`$ (36a)
$`\rho _{12}^D(t)`$ $`=`$ $`\rho _{12}^D(0)e^{(2i\lambda _1+k(2\overline{n}+1))t}.`$ (36b)
It is straightforward to obtain the adiabatic and weak coupling limit eqs.(33) from eqs.(36).
### 2.2 Dephasing process in two level system
Another interesting process well studied in the standard textbooks is the phase destroying process which might appear due to elastic collisions. In general, those effects are incorporated in the master equation of the two level model, besides the energy dissipation process studied in subsection 2.1. Since we are studying the coupling of the spin 1/2 to a dissipative medium in the weak coupling limit, the inclusion of the phase destroying process in eq.(32) gives corrections to the coefficients $`a_{ij}(\omega ,k;t)`$ in eqs.(33). Due to the linearity of the equations, the new imaginary phases coming from the dephasing process are added to the ones obtained previously. For the sake of simplicity we study the imaginary phases acquired by the variables $`\alpha _{11}(t)`$ and $`\alpha _{12}(t)`$ only due to the dephasing process. The master equation written in the diagonal frame is
$$\frac{d}{dt}\rho _D\left(t\right)=i[\lambda _1\sigma _z,\rho _D\left(t\right)]+\frac{k}{2}\left(\sigma _z\rho _D(t)\sigma _z\rho _D(t)\right).$$
(37)
In appendix A we present the master equation of the spin 1/2 with the dephasing effect included for arbitrary value of the angular velocity $`\omega `$ of the external magnetic field. In this subsection we study the time equations of coefficients $`\alpha _{11}(t)`$ and $`\alpha _{12}(t)`$ in the adiabatic and weak coupling regimes. Taking into account our discussion in section 2, eqs.(A.2) and (31) we obtain
$`{\displaystyle \frac{d}{dt}}\alpha _{11}(t)`$ $`=`$ $`0,`$ (38a)
$`{\displaystyle \frac{d}{dt}}\alpha _{12}(t)`$ $`=`$ $`k\alpha _{12}(t).`$ (38b)
that have the solutions:
$`\alpha _{11}(t)`$ $`=`$ $`\alpha _{11}(0),`$ (39a)
$`\alpha _{12}(t)`$ $`=`$ $`\alpha _{12}(0)e^{_0^tk𝑑t^{}}.`$ (39b)
From eq.(39a) we conclude that the Adiabatic Theorem is valid in this process, since the population at each quantum state does not vary along the adiabatic process. By analogous reasons to the ones discussed in subsection 2.1, the phase in eq.(39b) is not geometric but a time dependent imaginary phase that destroys the off-diagonal elements of the density matrix. Eq.(37) is exactly solved. The solutions (39) are easily recovered from the exact solutions when we calculate them in the adiabatic and weak coupling limits.
In order to understand why references - give an imaginary correction to the geometric phase and the models studied here do not, we compare eq.(2.3) of reference with our eqs.(34) and (39b). The coefficient that multiplies of variable $`C_\lambda (t)`$ in eq.(2.3) of reference has the form $`\theta (t)|\frac{d}{dt}|\psi (t)`$. Since the vector states depend on a periodic external parameter $`R(t)`$, this coefficient, in the adiabatic approximation, corresponds to a correction to Berry’s phase written as a closed curve in parameter space. The coefficient $`k(1+2\overline{n})`$ that multiplies $`\alpha _{12}(t)`$ in eq.(33b) does not arise from any variation of an external parameter. The same is true for the coefficient $`k`$ in eq.(38b). That is the reason that allows us to claim that the exponentials acquired from eqs.(34) and (39b) have no geometric origin, which means that the suppression terms are a function of time instead of some path parameter.
## 3 Contribution of the geometric phase to the <br>magnetization
In order to verify how the geometric phases and the dissipation affect the physical quantities, we return to the spin $`1/2`$ model in the presence of a reservoir of electromagnetic fields in thermal equilibrium. Under the initial condition $`\rho (0)=|\psi (0)\psi (0)|`$, with $`|\psi (0)=\mathrm{cos}(\alpha )|++\mathrm{sin}(\alpha )|`$, in the adiabatic limit and weak coupling limit we end up with
$`\rho _I^{11}(t)`$ $`=`$ $`{\displaystyle \frac{\overline{n}}{2\overline{n}+1}}+{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{1}{2\overline{n}+1}}+\mathrm{cos}(\theta 2\alpha )\right)e^{2k\left(2\overline{n}+1\right)t},`$ (40a)
$`\rho _I^{12}(t)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{sin}(2\alpha \theta )e^{2i\mu Bt}e^{k\left(2\overline{n}+1\right)t}e^{ia\omega (1\mathrm{cos}(\theta ))t}.`$ (40b)
We have introduced the Berry’s phase tracer $`a`$ that helps us identify the contribution of Berry’s phase to physical quantities. At the end, we take the tracer equal to one.
Thus we see that, in this model, dissipation does not affect Berry’s geometrical phase, but only makes it harder to observe their interference effect: such information is contained in the off-diagonal terms of the density matrix in the instantaneous frame, which vanish. By the way, we can observe that the system considered in this contribution is analogous to the case of the classical Foucault pendulum where dissipation diminishes the amplitude, but do not affect the rotation of the oscillation plane.
The definition of the geometrical phase can also be given in terms of a phase vector as in reference, where the discussion is confined to pure states. It is, however, a relatively simple matter to extend the definition of the phase vector for mixed states. In this case we find that the loss of coherence will shorten the phase vector in a manner which is completely analogous to what happens to Bloch’s vector. This norm reduction, as in the present case, reflects the asymptotic vanishing of off-diagonal density matrix elements.
Since we are interested on interference effects due to geometric phases, Bloch’s vector, defined in eq.(2), is suitable to provide a graphic visualization of the density matrix . From eq.(3), we have that these effects are clearly described by Bloch’s vector in the instantaneous frame. In this frame the projection of Bloch’s vector sweeps the $`xy`$ plane and makes an angle smaller than that of the magnetic field by an amount proportional to the solid angle $`\mathrm{\Omega }(\theta )`$, while its length decreases exponentially with a time rate of $`k(2\overline{n}+1)`$. We illustrate in figure 1 the shrinking of the projection of Bloch’s vector in the $`xy`$ plane, due to the presence of dissipation. Even though we have discussed the case $`k\omega `$, we choose $`k/\omega =10`$ in order to show more clearly the decreasing of this projection in an interval $`4\pi /\mu B`$. Shortening in $`z`$ is faster than in the $`xy`$ plane, causing a time dependent azimuthal angle. Figure 2 shows the plot of the time evolution of Bloch’s vector for $`k/\omega =0.2`$. For $`t\mathrm{}`$ Bloch’s vector has only non-zero $`z`$-component and $`S_z(t\mathrm{})=\frac{1}{2\overline{n}+1}`$.
For time $`t`$, such that $`tnT`$ ($`T=2\pi /\omega `$ and $`n`$ is an integer) the off-diagonal elements of the density matrix $`\rho _I(t)`$ depends on the chosen condition satisfied by $`\varphi _i;t|\frac{d}{dt}|\varphi _i;t`$ . For the sake of comparison with experiments it is necessary to calculate dissipation and geometrical phases effects on measurable quantities. In this model the natural candidates are the components of the magnetization vector $`\stackrel{}{m}(t)`$. Let us consider the $`z`$-component of magnetization whose average value has the expression $`\mathrm{Tr}(\rho (t)𝐦_𝐳)`$. Since the trace is independent of the particular basis applied to calculate it the result is independent of our particular choice of $`\varphi _i;t|\frac{d}{dt}|\varphi _i;t`$. In the adiabatic approximation and the weak coupling limit, using eqs.(40), we get the Fourier transform of $`m_z(t)`$,
$`\stackrel{}{m}_z(\omega ^{})`$ $`=`$ $`{\displaystyle \frac{\mu }{\sqrt{2\pi }}}[{\displaystyle \frac{\mathrm{cos}(\theta )}{2\overline{n}+1}}\pi \delta (\omega ^{})i\mathrm{cos}(\theta )({\displaystyle \frac{1}{2\overline{n}+1}}+\mathrm{cos}(\theta 2\alpha )){\displaystyle \frac{1}{\omega ^{}+2ki(2\overline{n}+1)}}+`$
$`i{\displaystyle \frac{1}{2}}\mathrm{sin}(\theta )\mathrm{sin}(\theta 2\alpha )({\displaystyle \frac{1}{\omega ^{}+\mathrm{\Gamma }+ki(2\overline{n}+1)}}+{\displaystyle \frac{1}{\omega ^{}\mathrm{\Gamma }+ki(2\overline{n}+1)}})],`$
where $`\alpha `$ is given by the initial condition, and the resonant frequency $`\mathrm{\Gamma }`$ is equal to
$$\mathrm{\Gamma }2\mu Ba\omega \mathrm{cos}(\theta ).$$
(41b)
The first term of eq.(LABEL:13) corresponds to the constant component of the magnetic field, and the second one shows the dissipation effects on this field component. The last term on the r.h.s. of the above equation displays a real frequency shift which contains the contribution of the geometrical phase, as can be seen from eq.(41b), due the presence of the tracer $`a`$, and a line broadening caused only by the dissipative evolution. These effects on the magnetization agree with the ones derived in reference where the path integral formalism was applied.
The expressions for the other components of magnetization are analogous, but somewhat lengthy.
## 4 Conclusions
In summary, we have presented an analytical solution of the adiabatic limit of a spin $`1/2`$ in a precessing magnetic field embedded in a weakly dissipative medium, introduced phenomenologically. We consider two distinct nonunitary contributions that were accounted for by a Lindblad type superoperator in the diagonal frame. We are able to derive analytical expressions for the geometric and imaginary phases in both cases in the presence of a weak dissipation in the adiabatic limit without further approximations.
From eqs.(34) and (39b) we get that the nature (path dependent or time dependent) of the imaginary phase acquired by $`\alpha _{12}(t)`$ depends on the mechanism that introduces the dissipation in the quantum system. In both cases that we have studied, the dissipation is present due to the two level system being in contact with a reservoir. The constant $`k`$ on the r.h.s. of eq.(34) is associated to the time rate of population and not due to the variation of any external parameter. Consequently the quantum geometric phase for $`k=0`$ is not modified by a complex value. An analogous argument is valid to explain why the imaginary phase acquired by the entry $`\rho _{12}^I(t)`$ in the dephasing process is not geometric, either. In this last model the Adiabatic Theorem continues to be true while it is not true anymore for the spin 1/2 coupled to the electromagnetic fields at thermal equilibrium. Differently from Ellinas et al. in reference we do not call this complex phase as Berry’s phase. We reserve the name of “Berry’s phase” only to phases (real or imaginary) that are path dependent.
Decoherence effects are present and their manifestation is the shortening of the three components of the Bloch’s vector. The fact that the dissipation effect causes the suppression of interference patterns due to the geometric phase is not a particular result for the chosen liouvillian. It is rather general, stemming from the fact that the dissipation mechanism is not related to the variation of any set of external periodic parameters.
## Acknowledgements
The authors are in debt with M.C. Nemes and J.G. Peixoto de Faria for useful discussions. The authors thank the referee for bringing to their attention the interesting model discussed in section 2.2. A.C.A.P. and M.T.T. thank E.V. Corrêa Silva for the careful reading of the manuscript. A.C. Aguiar Pinto thanks CNPq for financial support. M.T. Thomaz thanks CNPq for partial financial support.
## Appendix A: Master equations in the instantaneous basis of the hamiltonian
The master equation (32) in the instantaneous basis of the hamiltonian for arbitrary value of $`\omega `$ is
$$\frac{d}{dt}\rho _I\left(t\right)=i[\left(\mu B+\frac{\omega }{2}\right)\sigma _z\frac{\omega }{2}\sigma _n(t),\rho _I\left(t\right)]+k_I\rho _I(t)$$
(A.1a)
where $`\sigma _n(t)`$ is given by eq.(17b) and
$`_I\rho _I(t)`$ $`=`$ $`{\displaystyle \frac{2\overline{n}+1}{2}}\{2\rho _I(t)+(1\mathrm{\Lambda }^2)[\sigma _n(t)\rho _I(t)\sigma _n(t)e^{2i\omega t}\sigma _+(t)\rho _I(t)\sigma _+(t)`$ (A.1b)
$``$ $`e^{2i\omega t}\sigma _{}(t)\rho _I(t)\sigma _{}(t)]+(1+\mathrm{\Lambda }^2)[\sigma _+(t)\rho _I(t)\sigma _{}(t)+\sigma _{}(t)\rho _I(t)\sigma _+(t)]`$
$``$ $`\mathrm{\Lambda }\sqrt{1\mathrm{\Lambda }^2}[e^{i\omega t}(\sigma _n(t)\rho _I(t)\sigma _{}(t)+\sigma _{}(t)\rho _I(t)\sigma _n(t))+`$
$`+`$ $`e^{i\omega t}(\sigma _+(t)\rho _I(t)\sigma _n(t)+\sigma _n(t)\rho _I(t)\sigma _+(t))]\}`$
$``$ $`{\displaystyle \frac{1}{2}}\{\{\rho _I(t),\mathrm{\Lambda }\sigma _n(t)+\sqrt{1\mathrm{\Lambda }^2}(e^{i\omega t}\sigma _+(t)+e^{i\omega t}\sigma _{}(t)\}+2\mathrm{\Lambda }[\sigma _+(t)\rho _I(t)\sigma _{}(t)`$
$``$ $`\sigma _{}(t)\rho _I(t)\sigma _+(t)]+\sqrt{1\mathrm{\Lambda }^2}[e^{i\omega t}(\sigma _n(t)\rho _I(t)\sigma _+(t)\sigma _+(t)\rho _I(t)\sigma _n(t))+`$
$`+`$ $`e^{i\omega t}(\sigma _{}(t)\rho _I(t)\sigma _n(t)\sigma _n(t)\rho _I(t)\sigma _{}(t))]\}.`$
In eq.(A.1b), we define: $`\mathrm{\Lambda }\frac{1}{\lambda _1}\left(\mu B\mathrm{cos}(\theta )\frac{\omega }{2}\right)`$ and
$$\sigma _+(t)=e^{i\omega t}\left[\frac{\mathrm{sin}(\theta )}{2}\sigma _z+(\mathrm{cos}(\frac{\theta }{2}))^2e^{i\omega t}\sigma _+(\mathrm{sin}(\frac{\theta }{2}))^2e^{i\omega t}\sigma _{}\right]$$
(A.1c)
and $`\sigma _{}(t)\left(\sigma _+(t)\right)^{}`$.
The master equation (37) in the instantaneous basis of the hamiltonian for arbitrary value of $`\omega `$ is
$$\frac{d}{dt}\rho _I\left(t\right)=i[\left(\mu B+\frac{\omega }{2}\right)\sigma _z\frac{\omega }{2}\sigma _n(t),\rho _I\left(t\right)]+\frac{k}{2}_I\rho _I(t)$$
(A.2a)
where
$`_I\rho _I(t)`$ $`=`$ $`\mathrm{\Lambda }^2\sigma _n(t)\rho _I(t)\sigma _n(t)+\mathrm{\Lambda }\sqrt{1\mathrm{\Lambda }^2}[e^{i\omega t}(\sigma _n(t)\rho _I(t)\sigma _+(t)+\sigma _+(t)\rho _I(t)\sigma _n(t))+`$ (A.2b)
$`+`$ $`e^{i\omega t}(\sigma _n(t)\rho _I(t)\sigma _{}(t)+\sigma _{}(t)\rho _I(t)\sigma _n(t))]+(1\mathrm{\Lambda }^2)[e^{2i\omega t}\sigma _+(t)\rho _I(t)\sigma _+(t)+`$
$`+`$ $`e^{2i\omega t}\sigma _{}(t)\rho _I(t)\sigma _{}(t)+\sigma _+(t)\rho _I(t)\sigma _{}(t)+\sigma _{}(t)\rho _I(t)\sigma _+(t)]\rho _I(t).`$
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# Boundaries of zero scalar curvature in the AdS/CFT Correspondence
## 1. Introduction
In , Witten and Yau consider the AdS/CFT correspondence in the context of a Riemannian Einstein manifold $`M^{n+1}`$ of negative Ricci curvature which admits a conformal compactification (in the sense of Penrose ) with conformal boundary $`N^n`$. As discussed in , a conformal field theory on $`N`$ relevant to this correspondence is stable (with respect to the brane action on $`M`$) if the conformal class of the boundary contains a metric of positive scalar curvature and is unstable if it contains a metric of negative scalar curvature. In the borderline case of zero scalar curvature, the theory may be stable or unstable. Witten and Yau go on to prove that if the conformal class of $`N`$ contains a metric of positive scalar curvature, then $`M`$ and $`N`$ have several desirable properties: (1) $`N`$ is connected, which avoids the difficulty of coupling seemingly independent conformal field theories, (2) the $`n`$th homology of $`\overline{M}`$ vanishes, in particular $`M`$ has no wormholes and (3) at the fundamental group level, the topology of $`M`$ is “bounded by” the topology of $`N`$.
The aim of the present paper is to show that all of these results extend to the case where the conformal class of the boundary contains a metric of nonnegative scalar curvature. By a well known result of Kazdan and Warner , if $`N`$ has a metric of nonnegative scalar curvature, and if the scalar curvature is positive at some point, then N has a conformally related metric of positive scalar curvature. Hence, the essential case handled here is the case in which the conformal class of the boundary contains a metric of zero scalar curvature. The proof method used in this paper is quite different from, and in some sense dual to, that used in . The method of involves minimizing the co-dimension one brane action on $`M`$, and uses the machinery of geometric measure theory, while the arguments presented here use only geodesic geometry. We proceed to a precise statement of our results.
Let $`M^{n+1}`$ be a complete Riemannian manifold, with metric $`g`$, and suppose $`M`$ admits a conformal compactification, with conformal boundary (or conformal infinity) $`N`$. Thus it is assumed that $`M`$ is the interior of a compact manifold-with-boundary $`\overline{M}^{n+1}`$ and that there exists a smooth function $`r`$ on $`\overline{M}`$ such that (1) $`r>0`$ on $`M`$, (2) $`r=0`$ and $`dr0`$ along $`N=\overline{M}`$, and (3) $`r^2g`$ extends smoothly to a Riemannian metric $`\overline{g}`$ on $`\overline{M}`$. The induced metric $`h=\overline{g}|_{TN}`$ on $`N`$ changes by a conformal factor with a change in the *defining function* $`r`$, and so $`N`$ has a well defined conformal structure. If the conformal class of metrics on $`N`$ contains a metric of positive (resp., nonnegative, zero, etc.) scalar curvature we say that $`N`$ has positive (resp., nonnegative, zero, etc.) scalar curvature.
In , Witten and Yau consider conformally compactified orientable Einstein manifolds $`M^{n+1}`$ which satisfy $`\mathrm{Ric}=ng`$. More generally, their results allow $`\mathrm{Ric}ng`$ (i.e., $`\mathrm{Ric}(X,X)ng(X,X)=n`$ for all unit vectors $`X`$) provided $`\mathrm{Ric}ng`$ sufficiently fast as one approaches the conformal boundary $`N`$. (Regarding this fall-off, it is sufficient to require $`\mathrm{Ric}=ng+o(r^2)`$, where $`r`$ is a suitably chosen defining function for the conformal boundary. This is discussed in more detail below.) In this setting they prove that if $`N`$ has a component of positive scalar curvature then (1) $`N`$ is connected, (2) $`H_n(\overline{M},)=0`$, and (3) the map $`i_{}:\mathrm{\Pi }_1(N)\mathrm{\Pi }_1(\overline{M})`$ induced by inclusion $`i:N\overline{M}`$ is onto. The last result says that any loop in $`M`$ can be deformed to a loop in $`N`$. Thus, at the fundamental group level, the topology of $`M`$ can be no more complicated than the topology of $`N`$. In particular, if $`N`$ is simply connected, so is $`M`$. The following theorem generalizes these results by weakening the scalar curvature condition on $`N`$.
###### Theorem 1.
Let $`M^{n+1}`$ be a complete Riemannian manifold which admits a conformal compactification, with conformal boundary $`N^n`$, and suppose the Ricci tensor of $`M`$ satisfies, $`\mathrm{Ric}ng`$, such that $`\mathrm{Ric}ng`$ sufficiently fast on approach to conformal infinity. If $`N`$ has a component of nonnegative scalar curvature then the following properties hold.
1. $`N`$ is connected.
2. If $`M`$ is orientable, the $`n`$th homology of $`\overline{M}`$ vanishes, $`H_n(\overline{M},)=0`$.
3. The map $`i_{}:\mathrm{\Pi }_1(N)\mathrm{\Pi }_1(\overline{M})`$ ($`i`$ = inclusion) is onto.
The essential step in the proof is to establish part (a). Part (c) then follows from part (a) by a covering space argument, as noted in . In turn, as will be shown here, part (b) follows from part (c) via basic homology theory (essentially Poincaré duality; cf. , where similar arguments have been used). In order to give a flavor of the sorts of techniques that will be used to prove part (a), we will first give a proof of the connectedness of $`N`$ in the setting of Witten and Yau ; i.e., under the assumption that $`N`$ has a component of positive scalar curvature.
The result of Witten and Yau on the connectedness of the boundary is easily derived from the following proposition.
###### Proposition 2.
Suppose $`M^{n+1}`$ is a complete Riemannian manifold-with-boundary having Ricci curvature greater than or equal to $`n`$. If the boundary $`M`$ is compact and has mean curvature $`H>n`$ then $`M`$ is compact.
By our conventions, $`H=\mathrm{div}_MX`$, where $`X`$ is the outward pointing unit normal along $`M`$. Results similar to Proposition 2 are obtained in by minimizing the brane action and making use of the machinery of geometric measure theory. Here we give a proof of Proposition 2 using basic techniques in geodesic geometry. These arguments are reminiscent of the kinds of arguments used in the proof of the classical Hawking-Penrose singularity theorems.
Proof of Proposition 2: Suppose $`M`$ is noncompact. Then we can find a point $`qM`$ such that the distance from $`q`$ to $`M`$ is greater than $`\mathrm{coth}^1(1+\delta )`$. Here $`\delta >0`$ is chosen so that the mean curvature $`H`$ of $`M`$ satisfies $`Hn(1+\delta )`$. Let $`p`$ be a point on $`M`$ closest to $`q`$, and let $`\sigma :[0,\mathrm{}]M`$ be a unit speed minimal geodesic from $`p=\sigma (0)`$ to $`q=\sigma (\mathrm{})`$. Let $`\rho :M`$ be the distance function to the boundary,
$`\rho (x)=d(x,M)=\underset{yM}{inf}d(x,y).`$
In general, $`\rho `$ is continuous on $`M`$ and smooth outside the focal cut locus of $`M`$. In particular, since $`\sigma `$ realizes the distance to $`M`$, $`\rho `$ is smooth on an open set $`U`$ containing $`\sigma \{q\}`$.
For $`0s<\mathrm{}`$, let $`H(s)=\mathrm{}\rho (\sigma (s))=\mathrm{div}(\rho )(\sigma (s))`$. Geometrically, $`H(s)`$ is the mean curvature of the slice $`\rho =s`$, with respect to the unit normal $`\rho `$, at the point $`\sigma (s)`$. $`H=H(s)`$ obeys the well known traced Riccati equation ,
(1) $`H^{}=\mathrm{Ric}(\sigma ^{},\sigma ^{})+|B|^2,`$
where $`{}_{}{}^{}=d/ds`$ and $`B(s)=\mathrm{Hess}(\rho )(\sigma (s))=(\rho )(\sigma (s))`$. Since $`H(s)`$ is the trace of $`B(s)`$, the Schwarz inequality implies, $`|B|^2H^2/n`$. Equation 1, taken together with this inequality and the inequalities $`\mathrm{Ric}(\sigma ^{},\sigma ^{})n`$ and $`H(0)=H_M(p)n(1+\delta )`$, implies that $`(s):=H(s)/n`$ satisfies,
$`^{}^21,(0)1+\delta .`$
By comparison with the unique solution to: $`h^{}=h^21`$, $`h(0)=1+\delta `$, we obtain,
$`(s)\mathrm{coth}(as),`$
where $`a=\mathrm{coth}^1(1+\delta )<\mathrm{}=d(q,M)`$. This inequality implies that $`=(s)`$ is unbounded on $`[0,a)`$, which contradicts the fact that it is smooth on $`[0,\mathrm{})`$. Thus, $`M`$ must be compact.
We remark that the rigid version of Proposition 2, in which one assumes the weak inequality, $`Hn`$, holds, has previously been treated in the literature, cf., .
Return to the setting of Theorem 1, but in the case considered by Witten and Yau in which some component $`N_0`$, say, of $`N`$ has positive scalar curvature. We indicate how the connectedness of $`N`$ in this case follows from Proposition 2. Let $`r`$ be a defining function for the conformal boundary $`N`$, and let $`U`$ be a neighborhood of $`N_0`$ which does not meet any other components of $`N`$. Let $`M_t=M\{xU:r(x)<t\}`$. For $`t`$ sufficiently small, $`M_t`$ is a manifold-with-boundary, with boundary $`M_t`$ diffeomorphic to $`N_0`$, which satisfies the hypotheses of Proposition 2. (That the mean curvature of $`M_t`$ satisfies $`H>n`$ uses the assumption of positive scalar curvature on $`N`$. It also uses the fall-off condition on the Ricci curvature; without that, there are simple counter-examples to the Witten-Yau result.) We conclude that $`M_t`$ is compact, from which it follows that $`N`$ has no other components; i.e., $`N=N_0`$, and hence is connected.
We now describe briefly our approach to the proof of part (a) of Theorem 1. The idea, roughly, is as follows: First we show that if there is a sequence of compact hypersurfaces $`\mathrm{\Sigma }_k`$ in $`M`$ going to infinity in some end such that the mean curvature $`H_k`$ of $`\mathrm{\Sigma }_k`$ approaches $`n`$ “fast enough” then $`M`$ has only one end. Then we show that this rate is realized if the end admits a conformal compactification such that the conformal class of the boundary has a metric of zero scalar curvature. The first, and main, step is to establish a suitable refinement of Proposition 2.
###### Theorem 3.
Let $`M^{n+1}`$ be a complete Riemannian manifold having Ricci curvature greater than or equal to $`n`$. Fix a base point $`oM`$. Suppose there exists a sequence of compact hypersurfaces $`\{\mathrm{\Sigma }_k\}`$ satisfying the following conditions.
1. Each $`\mathrm{\Sigma }_k`$ separates $`M`$. We call the component of $`M\mathrm{\Sigma }_k`$ containing $`o`$ the inside of $`\mathrm{\Sigma }_k`$ and the other component the outside.
2. $`d(o,\mathrm{\Sigma }_k)\mathrm{}`$ as $`k\mathrm{}`$.
3. Denote by $`H_k`$ the mean curvature of $`\mathrm{\Sigma }_k`$ with respect to the outward normal, and let $`h_k`$ be the smaller of $`\mathrm{min}\{H_k(x):x\mathrm{\Sigma }_k\}`$ and n. Assume that
(2) $`\underset{k\mathrm{}}{lim}(nh_k)e^{2d(o,\mathrm{\Sigma }_k)}=0.`$
Then $`M`$ has one or two ends. If $`M`$ has two ends then $`M`$ is isometric to $`\times \mathrm{\Sigma }`$, with warped product metric $`dr^2+e^{2r}g_0`$, where $`\mathrm{\Sigma }`$ is compact and $`g_0`$ is a metric of nonnegative Ricci curvature on $`\mathrm{\Sigma }`$.
An *end* of $`M`$ is, roughly speaking, an unbounded component of the complement of a sufficiently large compact subset of $`M`$, cf. e.g., for a precise definition. If $`M`$ admits a conformal compactification then each of its ends is diffeomorphic to $`\times \mathrm{\Sigma }`$, where $`\mathrm{\Sigma }`$ is a component of the conformal boundary. In the special case of two ends in the theorem, $`M`$ has a *cusp* at one end, and hence does not admit a conformal compactification. (Regardless of curvature conditions, the ends of a conformally compactified manifold must have positive mean curvature near infinity.)
We recall an example considered in which satisfies all the hypotheses of Theorem 3 except for the mean curvature decay condition (2). Let $`(\mathrm{\Sigma },g_0)`$ be any compact negatively curved ($`\mathrm{Ric}=(n1)g`$) Einstein manifold of dimension $`n`$. Then $`M=\times \mathrm{\Sigma }`$, with warped product metric $`g=dr^2+\mathrm{cosh}^2(r)g_0`$ is an Einstein manifold satisfying, $`\mathrm{Ric}=ng`$. The slices $`\mathrm{\Sigma }_r=\{r\}\times \mathrm{\Sigma }`$, $`r>0`$, have mean curvature $`H(r)=n\mathrm{tanh}r`$. Hence $`H(r)n`$ as $`r\mathrm{}`$, but $`lim_r\mathrm{}e^{2r}(nH(r))=2n`$. This shows that the mean curvature condition in Theorem 3 is in some sense optimal.
In the next section we present the proof of Theorem 3, and in the final section we present the proof of Theorem 1.
## 2. Proof of Theorem 3
The proof of Theorem 3 is similar in spirit to the proof of the Cheeger-Gromoll splitting theorem , and makes use of (generalized) Busemann functions and the method of support functions . The method of support functions provides an elementary way to work with the Laplacian of certain geometrically defined functions, such as Busemann functions, which are in general only $`C^0`$.
Let $`M`$ be a Riemannian manifold and let $`fC^0(M)`$ be a continuous function on $`M`$. A *lower support* function for $`f`$ at $`pM`$ is a function $`\varphi `$ defined and continuous on a neighborhood $`U`$ of $`p`$ such that $`\varphi f`$ on $`U`$ and $`\varphi (p)=f(p)`$. We say that $`fC^0(M)`$ satisfies $`\mathrm{}fa`$ ($`a`$) in the *support sense* provided for each $`pM`$ and every $`ϵ>0`$ there exists a $`C^2`$ lower support function $`\varphi _{p,ϵ}`$ for $`f`$ at $`p`$ such that $`\mathrm{}\varphi _{p,ϵ}(p)aϵ`$. The Hopf-Calabi maximum principle asserts that if $`M`$ is connected and $`fC^0(M)`$ satisfies $`\mathrm{}f0`$ in the support sense then $`f`$ cannot attain a maximum unless it is constant. The proof of the Hopf-Calabi maximum principle is completely elementary; a short elegant proof is given in (cf., also ). We will need to make use of a slightly more general version of the Hopf-Calabi maximum principle.
###### Definition 1.
A function $`fC^0(M)`$ satisfies $`\mathrm{}fa`$, $`a`$, in the *generalized support sense* provided for each $`pM`$, there is a neighborhood $`U`$ of $`p`$ such that the following conditions hold.
1. There exists a sequence $`\{f_k\}`$, $`f_kC^0(U)`$, such that $`f_kf`$ uniformly on $`U`$.
2. $`\mathrm{}f_ka_k`$ on $`U`$ in the support sense, and $`a_ka`$.
###### Lemma 4.
*(Generalized maximum principle).* Suppose $`M`$ is a connected Riemannian manifold, and $`fC^0(M)`$ satisfies $`\mathrm{}f0`$ in the generalized support sense. Then, if $`f`$ attains a maximum, it is constant.
###### Proof.
The proof is a simple modification of the proof of the Hopf-Calabi maximum principle given in . We omit the details. ∎
One defines $`\mathrm{}fa`$ in the generalized support sense in a similar way, using $`C^2`$ *upper* support functions. By definition, $`\mathrm{}f=a`$ in the generalized support sense provided $`\mathrm{}fa`$ and $`\mathrm{}fa`$ in the generalized support sense. If $`\mathrm{}f=a`$ in the generalized support sense then $`fC^{\mathrm{}}(M)`$ and $`\mathrm{}f=a`$ in the usual sense. Indeed, for any small geodesic ball $`B`$, basic elliptic theory guarantees that the Dirichlet problem: $`\mathrm{}h=a`$, $`h|_B=f|_B`$, has a solution $`hC^{\mathrm{}}(B)C^0(\overline{B})`$. Then $`\mathrm{}(fh)=0`$ on $`B`$ in the generalized support sense, and the generalized maximum principle applied to $`\pm (fh)`$ implies that $`f|_B=h`$.
*Proof of Theorem 3:* Suppose that $`M`$ has more than one end. Then there is a compact set $`K`$ such that $`MK`$ has at least two unbounded components $`E_1`$ and $`E_2`$, say. Since $`MK`$ has at most finitely many components we may assume without loss of generality that $`\mathrm{\Sigma }_kE_1`$ for all $`k`$. We now construct a *line* in $`M`$. (Recall, a line is a complete unit speed geodesic, each segment of which realizes the distance between its endpoints.) Let $`\{q_k\}`$ be a sequence in $`E_2`$ going to infinity, $`d(o,q_k)\mathrm{}`$. Let $`p_k`$ be a point on $`\mathrm{\Sigma }_k`$ closest to $`q_k`$, and let $`\sigma _k:[a_k,b_k]M`$ be a unit speed minimal geodesic from $`p_k`$ to $`q_k`$. Since $`\sigma _k`$ meets $`K`$, we may parameterize $`\sigma _k`$ so that $`\sigma _k(0)K`$. By passing to a subsequence if necessary we have $`\sigma _k(0)\overline{o}M`$ and $`\sigma _k^{}(0)XT_{\overline{o}}M`$. Let $`\sigma :M`$ be the geodesic satisfying $`\sigma (0)=\overline{o}`$ and $`\sigma ^{}(0)=X`$. As $`\sigma `$ is the limit of minimal segments, it is a line in $`M`$.
We consider two Busemann functions on $`M`$, a Busemann function associated with the sets $`\mathrm{\Sigma }_k`$ (in the sense decribed in ) and the standard Busemann function associated with the ray (half-line) $`\sigma |_{[0,\mathrm{})}`$. For each $`k`$, let $`\beta _k:M`$ be the function defined by, $`\beta _k(x)=d(\overline{o},\mathrm{\Sigma }_k)d(x,\mathrm{\Sigma }_k)`$. The triangle inequality implies that $`\beta _k`$ is Lipschitz continuous, with Lipschitz constant one, and satisfies $`|\beta _k(x)|d(x,\overline{o})`$. Hence, the family of functions $`\{\beta _k\}`$ is equicontinuous and uniformly bounded on compact subsets. Thus, by Ascoli’s theorem, and passing to a subsequence if necessary, we have that $`\beta _k`$ converges on compact subsets to a continuous function $`\beta :M`$, called the Busemann function associated with $`\{\mathrm{\Sigma }_k\}`$. We will ultimately show that $`\beta C^{\mathrm{}}(M)`$ and satisfies $`\mathrm{}\beta =n`$, from which the special form of $`(M,g)`$ in the statement of Theorem 3 will readily follow. The first step is to establish the following.
Claim. *$`\mathrm{}\beta n`$ in the generalized support sense*.
Let $`p`$ be any point in $`M`$ and let $`B=B(p,r)`$ be a small geodesic ball centered at $`p`$ of radius $`r`$. To prove the claim we show that $`\mathrm{}\beta _kn_k`$ on $`B`$ in the support sense, where $`n_kn`$. Given $`qB`$ and $`ϵ>0`$, we construct a support function $`\beta _k^{q,ϵ}`$ for $`\beta _k`$ at $`q`$ as follows. Let $`z`$ be a point on $`\mathrm{\Sigma }_k`$ closest to $`q`$, and let $`\gamma :[0,\mathrm{}]:M`$ be a unit speed minimal geodesic from $`z=\gamma (0)`$ to $`q=\gamma (\mathrm{})`$. Let $`V`$ be a small neighborhood of $`z`$ in $`\mathrm{\Sigma }_k`$. By bending $`V`$ slightly toward the outside of $`\mathrm{\Sigma }_k`$ we obtain a smooth hypersurface $`V^{}`$ with the following properties: (1) $`zV^{}`$ is the unique closest point in $`V^{}`$ to $`q`$, (2) the second fundamental form of $`V^{}`$ at $`z`$ (with respect to the outward normal) is strictly less than that of $`V`$, and (3) the mean curvature of $`V^{}`$ at $`z`$ satisfies $`H_V^{}(z)H_V(z)ϵh_kϵ`$. By construction, $`\gamma `$ minimizes the distance from $`q`$ to $`V^{}`$, and there are no focal cut points to $`V^{}`$ on $`\gamma `$ (in particular, $`q`$ is not a focal cut point). Hence the function $`\beta _k^{q,ϵ}(x)=d(\overline{o},\mathrm{\Sigma }_k)d(x,V^{})`$ is a lower support function for $`\beta _k`$ at $`q`$ which is smooth on a neighborhood of $`\gamma `$.
For $`0s\mathrm{}`$, let $`H(s)=\mathrm{}\beta _k^{q,ϵ}(\gamma (s))`$. Arguing as in Proposition 2, $`(s)=H(s)/n`$ satisfies, $`^{}^21`$, $`(0)(h_kϵ)/n`$. Since $`(h_kϵ)/n<1`$, by comparing with the unique solution to $`h^{}=h^21`$, $`h(0)=(h_kϵ)/n`$, we obtain, $`(s)\mathrm{tanh}(as)`$, where $`a=\mathrm{tanh}^1((h_kϵ)/n)=\frac{1}{2}\mathrm{ln}(\frac{n+h_kϵ}{nh_k+ϵ})`$. Setting $`s=\mathrm{}`$ in this inequality, we obtain,
$`\mathrm{}\beta _k^{q,ϵ}(q)=H(\mathrm{})`$ $``$ $`n\mathrm{tanh}(a\mathrm{})=n{\displaystyle \frac{e^{2a}e^2\mathrm{}}{e^{2a}+e^2\mathrm{}}}`$
$`=`$ $`n{\displaystyle \frac{(n+h_kϵ)(nh_k+ϵ)e^2\mathrm{}}{(n+h_kϵ)+(nh_k+ϵ)e^2\mathrm{}}}.`$
Now, by the triangle inequality, $`\mathrm{}=d(q,\mathrm{\Sigma }_k)r+d(o,p)+d(o,\mathrm{\Sigma }_k)`$, and hence $`e^2\mathrm{}Ce^{2d(o,\mathrm{\Sigma }_k)}`$, where $`C=e^{2(r+d(o,p))}`$. Making use of this latter inequality in (2) we conclude that $`\mathrm{}\beta _kn_k`$ on $`B`$ in the support sense, where
(4) $`n_k=n{\displaystyle \frac{(n+h_k)C(nh_k)e^{2d(o,\mathrm{\Sigma }_k)}}{(n+h_k)+C(nh_k)e^{2d(o,\mathrm{\Sigma }_k)}}}.`$
Invoking the mean curvature condition (2), we see that $`n_kn`$. This yields the claim.
We now consider the standard Busemann function associated to the ray $`\sigma |_{[0,\mathrm{})}`$. For each $`s>0`$, define the function $`b_s:M`$ by, $`b_s(x)=d(\overline{o},\sigma (s))d(x,\sigma (s))=sd(x,\sigma (s))`$. For each $`xM`$, $`b_s(x)`$ is increasing in $`s`$ and bounded by $`d(\overline{o},x)`$. The Busemann function $`b:M`$ of $`\sigma |_{[0,\mathrm{})}`$ is defined to be the limit function, $`b(x)=lim_s\mathrm{}b_s(x)`$. Because the family $`\{b_s\}`$ is equicontinuous, $`b`$ is continuous.
In the present situation in which the Ricci curvature is greater than or equal to $`n`$, it is known that $`\mathrm{}bn`$ in the generalized support sense (in fact, in the support sense, cf. ). As the arguments involved to show this are similar to (but simpler than) the arguments used in the proof of the claim, we make only a few brief comments. The relevant support functions $`b_s^{q,ϵ}`$ for $`b_s`$ are defined as follows. Let $`\gamma :[0,\mathrm{}]M`$ be a unit speed minimal geodesic from $`q`$ to $`\sigma (s)`$. The function $`b_s^{q,ϵ}`$ defined by, $`b_s^{q,ϵ}(x)=s(ϵ+d(x,\gamma (\mathrm{}ϵ))`$ is a lower support function for $`b_s`$ at $`q`$ which is smooth near $`q`$. By standard comparison techniques like those used above, $`b_s^{q,ϵ}`$ satisfies $`\mathrm{}b_s^{q,ϵ}(q)n\mathrm{coth}(\mathrm{}ϵ)`$. From this it easily follows that $`\mathrm{}bn`$ in the generalized support sense.
To summarize, we have shown that the Busemann functions $`\beta `$ and $`b`$ satisfy $`\mathrm{}\beta n`$ and $`\mathrm{}bn`$ in the generalized support sense. Hence the sum $`f=\beta +b`$ satisfies $`\mathrm{}f0`$ in the generalized support sense. Moreover, $`f`$ satisfies, $`f0`$ on $`M`$. Indeed, we have,
$`\beta (x)+b_s(x)`$ $`=`$ $`\underset{k\mathrm{}}{lim}[d(\sigma (0),\mathrm{\Sigma }_k)d(x,\mathrm{\Sigma }_k)]+sd(x,\sigma (s))`$
$`=`$ $`\underset{k\mathrm{}}{lim}[d(\sigma _k(0),\mathrm{\Sigma }_k)d(x,\mathrm{\Sigma }_k)]+\underset{k\mathrm{}}{lim}[sd(x,\sigma _k(s))]`$
$`=`$ $`\underset{k\mathrm{}}{lim}[d(\sigma _k(0),\mathrm{\Sigma }_k)+s(d(\sigma _k(s),x)+d(x,\mathrm{\Sigma }_k))].`$
By the triangle inequality,
(6) $`d(\sigma _k(s),x)+d(x,\mathrm{\Sigma }_k)d(\sigma _k(s),\mathrm{\Sigma }_k)=d(\sigma _k(0),\mathrm{\Sigma }_k)+s.`$
The inequalities (2) and (6) imply $`\beta (x)+b_s(x)0`$. Letting $`s\mathrm{}`$, we obtain $`f=\beta +b0`$. But note that $`f(\overline{0})=\beta (\sigma (0))+b(\sigma (0))=0+0=0`$. Thus, by the generalized maximum principle $`f0`$. Hence, $`\beta =b`$ and so satisfies $`\mathrm{}\beta n`$ in the generalized support sense. Since $`\mathrm{}\beta n`$ in the generalized support sense, as well, we conclude from the discussion after Lemma 4 that $`\beta `$ is smooth and satisfies $`\mathrm{}\beta =n`$ in the usual sense.
It is known that Busemann functions, where differentiable, have unit gradient, and hence $`|\beta |=1`$ everywhere. (Briefly, this follows from the fact that $`\beta `$ satisfies, $`|\beta (q)\beta (p)|d(p,q)`$, with equality holding when $`p`$ and $`q`$ are on an *asymptotic ray*, cf. Lemma 6 in .) This has as well known consequences the fact that $`\beta `$ is a geodesic vector field, i.e. its integral curves are unit speed geodesics, and that $`\beta `$ satisfies (),
(7) $`\beta (\mathrm{}\beta )=\mathrm{Ric}(\beta ,\beta )+|\mathrm{Hess}\beta |^2.`$
(Compare with Equation 1.) Since the left hand side vanishes, we have $`|\mathrm{Hess}\beta |^2=\mathrm{Ric}(\beta ,\beta )n`$. But from the Schwarz inequality, $`|\mathrm{Hess}\beta |^2|\mathrm{}\beta |^2/n=n`$. Hence, equality holds, which implies,
(8) $`\mathrm{Hess}\beta |_\beta ^{}=g|_\beta ^{}\text{ and }\mathrm{Ric}(\beta ,\beta )=n.`$
Exponentiating out from the slice $`\mathrm{\Sigma }=\beta ^1(0)`$ along its normal geodesics ($`=`$ integral curves of $`\beta `$) establishes a global diffeomorphism $`M\times \mathrm{\Sigma }`$, with respect to which $`g`$ takes the form
(9) $`g=dr^2+g_{ij}(r,x)dx^idx^j,`$
where $`/r=\beta `$ and $`g_r=g_{ij}(r,x)dx^idx^j`$ is the induced metric on the slice $`\mathrm{\Sigma }^r=\beta ^1(r)\{r\}\times \mathrm{\Sigma }`$. Along $`\mathrm{\Sigma }^r`$, $`\mathrm{Hess}\beta (_i,_j)=\frac{1}{2}_rg_{ij}=g_{ij}`$ (by the first equation in (8)), which gives,
(10) $`g=dr^2+e^{2r}g_{ij}(0,x)dx^idx^j,`$
as required. The second equation in (8) and a calculation show that $`g_0=g_{ij}(0,x)dx^idx^j`$ is a metric of nonnegative Ricci curvature. Finally, $`\mathrm{\Sigma }`$ is compact, otherwise $`M\times \mathrm{\Sigma }`$ has only one end, contrary to assumption. This concludes the proof of Theorem 3.
## 3. Proof of Theorem 1
Let $`N_0`$ be the component in the statement of the theorem which admits in its conformal class a metric $`h`$ of nonnegative scalar curvature. As discussed in the introduction, we may assume, in fact, that $`h`$ is a metric of zero scalar curvature. Then there is a defining function $`r`$ such that near $`N_0`$, $`M`$ has the form, $`M=[0,r_0)\times N_0`$, with metric $`g`$ of the form,
$`g={\displaystyle \frac{1}{r^2}}\overline{g}={\displaystyle \frac{1}{r^2}}({\displaystyle \frac{1}{\overline{g}(\overline{}r,\overline{}r)}}dr^2+g_r),`$
where $`g_r`$ is the metric induced on $`N_r=\{r\}\times N_0`$ from $`\overline{g}`$, such that $`g_0=h`$. Assume that $`(M,g)`$ is Einstein, with $`\mathrm{Ric}=ng`$, or more generally that $`(M,g)`$ satisfies, $`\mathrm{Ric}ng`$ such that, as part of our fall-off assumption, the scalar curvature $`S`$ of $`(M,g)`$ satisfies,
$`Sn(n+1)\text{ as }r0.`$
As a computation shows, this implies that $`\overline{g}(\overline{}r,\overline{}r)=1`$ along $`N_0`$. Then as described in (cf., Lemma 2.1) the defining function $`r`$ can be chosen uniquely in a neighborhood of $`N_0`$ so that $`\overline{g}(\overline{}r,\overline{}r)=1`$ in this neighborhood. Thus, we may assume that in $`[0,r_0)\times N_0`$, $`g`$ has the form,
(11) $`g={\displaystyle \frac{1}{r^2}}(dr^2+g_r).`$
With respect to this distinguished defining function we impose the following fall-off requirement,
(12) $`r^2(Ric+ng)0\text{ uniformly as }r0.`$
This condition is compatible with the fall-off condition considered in . Let us also emphasize that these fall-off conditions are automatically satisfied when $`M`$ is Einstein with $`\mathrm{Ric}=ng`$.
The Gauss equation in $`(M,g)`$ implies,
(13) $`H^2=\widehat{S}S+\mathrm{Ric}(X,X)+|B|^2,`$
where, for each $`r`$, $`\widehat{S}`$, $`B`$ and $`H`$ are, respectively, the scalar curvature, second fundamental form and mean curvature of $`N_r`$, and $`X=r/r`$ is the outward unit normal to $`N_r`$. For each $`r`$, let $`\overline{S}`$ denote the scalar curvature of $`N_r`$ in the metric $`g_r`$; $`\overline{S}`$ and $`\widehat{S}`$ are related by $`\widehat{S}=r^2\overline{S}`$. Using this and the inequality $`|B|^2H^2/n`$, (13) implies the inequality,
(14) $`n^2H^2{\displaystyle \frac{n}{n1}}(r^2\overline{S}+\kappa ),`$
where $`\kappa =2\mathrm{Ric}(X,X)Sn(n1)`$. It follows from (12) that $`r^2\kappa 0`$ uniformly as $`r0`$. Since $`H>0`$ for $`r`$ sufficiently small, we have $`n^2H^2n(nH)`$ at points where $`Hn`$. This, together with (14) implies,
(15) $`r^2(nH){\displaystyle \frac{1}{n1}}(\overline{S}+r^2\kappa )\text{ where }Hn.`$
Pick a sequence $`r_k0`$, and set $`\widehat{N}_k=N_{r_k}`$. Given $`o\widehat{N}_{k_0}`$, we have,
(16) $`e^{d(o,\widehat{N}_k)}=e^{_{r_k}^{r_{k_0}}\frac{1}{r}𝑑r}=r_{k_0}r_k^1.`$
Let $`(r_k,x_k)\widehat{N}_k`$ be a point where the mean curvature of $`\widehat{N}_k`$ achieves a minimum. We may assume this minimum mean curvature $`h_k=H(r_k,x_k)`$ is less than or equal to $`n`$, otherwise by the Witten-Yau result we are done. By passing to a subsequence we may further assume $`(r_k,x_k)(0,x_0)N_0`$. Then, setting $`(r,x)=(r_k,x_k)`$ in (15) we obtain,
(17) $`e^{2d(o,\widehat{N}_k)}(nh_k){\displaystyle \frac{r_{k_0}^2}{n1}}(\overline{S}(r_k,x_k)+r_k^2\kappa (r_k,x_k)).`$
Since, as $`k\mathrm{}`$, $`\overline{S}(r_k,x_k)\overline{S}(0,x_0)=0`$ and $`r_k^2\kappa (r_k,x_k)0`$, we have that condition (2) in Theorem 3 is satisfied. Then by Theorem 3 and remarks in the paragraph following its statement, $`M`$ has only one end and hence $`N=N_0`$, i.e., $`N`$ is connected.
This concludes the proof of part (a). Part (c) follows from part (a), just as in the proof of Theorem 3.3 in . One passes to the covering space $`\overline{M}^{}`$ of $`\overline{M}`$ associated with the subgroup $`i_{}(\mathrm{\Pi }_1(N))`$ of $`\mathrm{\Pi }_1(\overline{M})`$. The boundary $`\overline{M}^{}`$ contains a copy of $`N`$ and has more than one component if $`i_{}:\mathrm{\Pi }_1(N)\mathrm{\Pi }_1(\overline{M})`$ is not onto, contradicting part (a) applied to $`\overline{M}^{}`$. Part (b) follows from part (c) by some basic homology theory, as we now describe. Similar arguments have been used in where related results in the spacetime setting have been obtained.
To prove part (c) consider the relative homology sequence for the pair $`\overline{M}N`$ (all homology is over $``$),
(18) $`\mathrm{}H_1(N)\stackrel{\alpha }{}H_1(\overline{M})\stackrel{\beta }{}H_1(\overline{M},N)\stackrel{}{}\stackrel{~}{H}_0(N)=0.`$
(Here $`\stackrel{~}{H}_0(N)`$ is the reduced zeroth dimensional homology group.) To make use of part (c), we use the fact that the first integral homology of a space is isomorphic to the fundamental group modded out by its commutator subgroup. Hence, modding out by the commutator subgroups of $`\mathrm{\Pi }_1(N)`$ and $`\mathrm{\Pi }_1(\overline{M})`$, we obtain a surjective linear map from $`H_1(N)`$ to $`H_1(\overline{M})`$, i.e., $`\alpha `$ in (18) is onto. Since $`\alpha `$ is onto, $`\mathrm{ker}\beta =\mathrm{im}\alpha =H_1(\overline{M})`$ which implies $`\beta 0`$. Hence $`\mathrm{ker}=\mathrm{im}\beta =0`$, and thus $``$ is injective. This implies that $`H_1(\overline{M},N)=0`$. But by Poincaré duality for manifolds-with-boundary, $`H_1(\overline{M},N)H^n(\overline{M})H_n(\overline{M})`$, where for the second isomorphism we have used the fact that $`H_n(\overline{M})`$ is free (cf., ). We conclude that $`H_n(\overline{M},)=0`$.
If $`M`$ is nonorientable (which, by part (a) and a covering space argument, can happen if and only if $`N`$ is nonorientable), essentially the same argument shows $`H_n(\overline{M},/2)`$ vanishes.
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# Excluded volume effects on the structure of a linear polymer under shear flow
## I Introduction
Excluded volume (EV) and hydrodynamics interactions (HI) have been recognized as fundamental physical phenomena to explain the peculiarities of single chain behavior in dilute polymer solutions. For a solution at rest, their effects on the single chain structure and dynamics have been carefully investigated over the years both theoretically and experimentally. When a dilute solution is subjected to a flow field, the situation is much more complex. In addition to finite extensibility (FE) effects to be considered at high strain, the role of EV and HI effects is still largely debated. This situation certainly arises from the difficulties to solve the kinetic equations for the dynamics of the simple coarse-grained models when EV and HI are considered. Some progresses have been made concerning the HI effects in absence of EV interactions by the use of a self consistent preaveraging approximations , a non equilibrium version of the more famous equilibrium Zimm theory. Renormalization Group Theory techniques have also been applied to solve suitable dynamic equations for the chain . The main focus of those investigations was the rheological response of the system although some results on the chain structure where also provided. An attempt to introduce EV effects in the Rouse model has also been performed by perturbation theory expansion.
On the experimental side the situation is similar in that measurements of single chain properties in flow are considerably more difficult than at rest. However, several scattering experiments have been performed, mainly Light Scattering (LS), with the aim to investigate the global properties of the chain such as the orientation and the deformation of the coil by the flow . The effect of EV interaction over global chain properties has been discussed in a series of recent LS studies where some systematic trends were reported. For dilute solutions in shear flow at a fixed reduced shear rate, a systematic weakening of the alignment of the main axis of the chain in the flow direction has been observed for increasing solvent quality. At a given reduced shear rate, the measured expansion ratio indicates that good solvent chains are less deformed than theta chains . Inspection of the flow effects on the internal structure of the coil can only be obtained by Small Angle Neutron Scattering (SANS) and indeed some data have been already published. However experimental limitations are here even more severe than for LS, and a quantitative study for various solvent qualities is still missing.
New experimental insight in the single chain behavior under flows and stretching forces has been provided recently by fluorescence techniques on biological molecules like DNA . Two papers on the single chain dynamics under shear flow have been published recently by S. Chu and coworkers and by D. Wirtz and coworkers in both cases employing DNA chains. In the former paper,the mean and time fluctuations of the fractional extension of DNA under shear flow in the flow direction was reported for reduced shear rates larger than one and culminating at $`\beta =76`$. In the second paper, the authors have shown that the dynamics of a DNA chain under weak shear flow ($`\beta 0.1`$) presents an unexpected wide range of relaxation times related to the particular initial configuration prior to the flow inset.
An interesting alternative to the experimental tools in the context of polymer physics are the numerical techniques such as Monte Carlo (MC), Molecular Dynamics (MD) and Brownian Dynamics (BD) simulations. They can be used to solve simple models for the system of interest and therefore test or suggest theoretical developments or stimulate new experiments. On the other hand simulations can help interpreting experimental data in a more complete and coherent fashion.
Many numerical studies of the behavior of a single chain under shear flow have already appeared by various techniques such as BD , MD and even non-equilibrium MC . To mention a recent simulation work , BD studies for a chain without HI and EV found a reasonable agreement with the infinite chain Rouse model predictions but came to the disturbing conclusion that, while the relative elongation of DNA chain in the flow direction is correctly explained by the chain model, the same model fails to explain the global expansion ratio of a polystyrene chain which is experimentally much lower than predicted.
To our knowledge, the role of the excluded volume effect on the structure of the chain under SF has not been systematically studied by simulations. In the present paper we try to elucidate this effect by comparing the analytical results of the Rouse model (in the continuum limit) with BD data for a suitable microscopic model of flexible chain with excluded volume interactions and unavoidable FE features.
The structure of a single chain in shear flow is often discussed by exploiting an analogy with the well known structure of a single chain subjected to a pair of equal and opposite stretching forces applied to its ends. In this case the ”macroscopic” behavior of long chains, as measured by the elasticity law, is characterized by a linear hookean regime for weak forces followed by a non linear power law regime for strong forces and ultimately by specific finite extensibility effects. The crossover between the first two regimes is linked to the appearance of the so called “tensile” blobs in the chain structure.
A very clear signature of the blobs can be detected in the chain structure factor as a crossover from ideal behavior at small q’s to excluded volume statistics at large q’s. This has been shown long time ago by Neutron Scattering experiments for “thermal” blobs in the semi-diluted solutions at rest , and recently by MC simulation for “tensile” blobs in the stretched chain problem. In the latter case, the crossover is detected for q orthogonal to the external force and can be ascribed to the fact that scattering at small q’s in these directions comes from pairs of monomers which are far apart in the elongation direction and therefore do not interact directly through the EV potential.
For the chain in shear flow Onuki has formulated a blob model as a simple extension of the Pincus-deGennes theory for the stretched chain case. The blob size is related to the bare shear rate and the chain at high shear rates is imagined as a string of shear blobs aligned most of the time in the flow direction. Such a structure which corresponds to an aligned fluctuating rod, should yield a scattering signal for $`q`$ in the flow direction (flow and elongation directions provide essentially the same scattering signal at high enough shear rate) similar to the stretched chain signal for $`q`$ aligned in the stretching direction . For the stretched chain such behavior is described quantitatively by the Rouse model with a suitable choice of the chain elongation and the longitudinal and transverse fluctuations . In the present paper we test these ideas on the form factor of a single EV chain in flow for shear rates which are sufficiently high to reach the plausible shear blob regime but also sufficiently small to remain in the scaling regime, i.e. avoiding finite chain effects. Our strategy is to analyze to which extend the low q behavior resembles the single chain structure factor of ideal chains in the presence of shear flow, which can be derived exactly, and whether the high q regime reproduces the structure of unperturbed EV chains. In order to observe such phenomena taking place on different length scales, we need sufficiently long chains and therefore we have performed BD simulations of EV chains in absence and presence of shear flow.
The paper is organized as follows. In section II we set up the essential phenomenology concerning a chain in shear flow and define the relevant structural characteristics of the chain. Section III describes the procedure to solve the Rouse model in flow and provides the results for the specific SF case. This is not a totally new analytical development as we are aware of at least two previous works on it . The novelty here is however the calculation of the structure factor in the whole range of $`q`$ vectors and the use of such behavior to pinpoint the signature of shear blobs in the EV chain structure factor. In the next section IV we describe the model for EV chains and the BD technique used. In section V we discuss the results of BD simulations. Finally section VI is devoted to our conclusion and perspectives.
## II Phenomenological framework
When a dilute polymer solution is subjected to a simple shear flow, characterized by the velocity field $`𝐮=\dot{\gamma }y\widehat{𝐱}`$ where$`\dot{\gamma }`$ is the shear rate, the chains in the solution are oriented and deformed according to the flow. Those phenomena can be measured through the anisotropy arising in any tensorial quantity associated to the chain, such as for instance the gyration tensor $`𝐆`$ and the order parameter tensor $`𝐎`$ defined as
$`𝐆`$ $`=`$ $`{\displaystyle \frac{1}{2N^2}}{\displaystyle \underset{i,j=1}{\overset{N}{}}}<(𝐑_i𝐑_j)(𝐑_i𝐑_j)>`$ (1)
$`𝐎`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N1}{}}}{\displaystyle \frac{𝐮_i𝐮_i}{|𝐮|^2}}{\displaystyle \frac{\mathrm{𝟏}}{3}}`$ (2)
where $`𝐑_i`$ are the coordinates of the i-th monomer in the chain, $`𝐮_i=𝐑_{i+1}𝐑_i`$ and $`\mathrm{𝟏}`$ is the unit tensor. At equilibrium the system is isotropic and any tensorial quantity, say $`𝐀`$ reduces to a scalar quantity. Under shear flow, symmetry imposes $`A_{xz}=A_{zx}=A_{yz}=A_{zy}=0`$. Therefore there are at most four independent components of any tensor to be monitored, namely the three diagonal elements and the off-diagonal $`A_{xy}=A_{yx}`$. The orientational angle $`\chi _A`$, defined through the relation
$$cot(2\chi _A)=\frac{A_{xx}A_{yy}}{2A_{xy}}$$
(3)
measure the rotation around $`\widehat{𝐳}`$ of the principal axes (I,II,III) of the tensor $`𝐀`$ with respect to the flow axes (x,y,z). In shear flow, $`A_{xy}`$ starts linearly with $`\dot{\gamma }`$ while the first contribution to $`A_{xx}A_{yy}`$ is of order $`\dot{\gamma }^2`$. Therefore the linear (Newtonian) regime is characterized by $`\chi _A=\pi /4`$. Outside the linear regime $`\chi _A`$ decreases to zero for increasing shear rate.
The deformation ratios, defined by
$$\delta A_{\alpha \alpha }=\frac{A_{\alpha \alpha }(\beta )}{A_{\alpha \alpha }(0)}1$$
(4)
measures the amount of deformation of the chain in flow, either in the flow reference frame ($`\alpha =x,y,z`$), or in the molecular reference frame ($`\alpha =I,II,III`$). When defined in the flow reference frame, the first contribution to those quantities is of order $`\dot{\gamma }^2`$, so that the Newtonian regime is characterized by the absence of deformation. Outside the linear regime, the chain is elongated in the flow direction ($`\delta A_{xx}>0`$) and compressed in the in-plane gradient direction ($`\delta A_{yy}<0`$), while the neutral (out of plane) direction is only slightly decreased ($`\delta A_{zz}0`$).
## III Structure factor of the Rouse chain under shear flow
To compute the structure factor of an ideal gaussian chain (Rouse chain) we adopt the continuous chain model introduced by Edwards in which the only energy term is a quadratic potential acting between nearest neighbors along the chain. For such a model the equilibrium probability distribution in configurational space is
$$\mathrm{\Psi }[𝐑]exp\left(\frac{\chi }{2}_0^N𝑑n\left(\frac{𝐑_n}{n}\right)^2\right)$$
(5)
where $`\chi =3k_BT/d^2`$, $`𝐑_n`$ represents the positions of the chain monomers, $`0nN`$ is the continuous monomer index and $`d`$ is the single bond variance of the model. The steady state distribution in presence of a generic homogeneous solvent flow field can be obtained in terms of the normal modes of the model $`(𝐗_p,(p=0,\mathrm{}))`$. For homogeneous flows to which we limit our discussion, the normal modes are defined by the same transformation as at equilibrium. In those new coordinates the distribution is factorized over the normal modes
$$\widehat{\mathrm{\Psi }}(\{𝐗\},t)=\underset{p=0}{\overset{\mathrm{}}{}}\widehat{\psi }_p(𝐗_p,t)$$
(6)
In SF, the steady state distribution of the p-th mode is
$$\widehat{\psi }_p(𝐗_p)=\left(\frac{\chi _1}{2\pi k_BT}\right)^{3/2}\frac{p^3}{\sqrt{1+\dot{\gamma }^2\left(\frac{\tau _1}{2p^2}\right)^2}}e^{\frac{\chi _1}{2k_BT}p^2𝐗_p\beta _p^1𝐗_p}$$
(7)
where the matrix $`\beta _p`$ is defined as
$$\beta _p=\mathrm{𝟏}+\frac{\tau _1}{2p^2}(𝐤+𝐤^T)+2\left(\frac{\tau _1}{2p^2}\right)^2𝐤𝐤^T$$
(8)
and $`𝐤=𝐮`$ is the velocity gradient tensor. The derivation of eqs. (7) and (8) is given in appendix A.
The chain structure factor is defined as
$$S(𝐪,N,\dot{\gamma })=\frac{1}{N}_0^N_0^N𝑑n𝑑m<e^{i𝐪𝐑_{nm}}>_{\dot{\gamma }}$$
(9)
where $`𝐑_{nm}=𝐑_n𝐑_m`$ and $`<\mathrm{}>_{\dot{\gamma }}`$ indicates a statistical averages over the non equilibrium distribution eq.(6). In terms of normal modes we have
$$𝐑_{nm}=2\underset{p=1}{\overset{\mathrm{}}{}}\left(cos\frac{p\pi n}{N}cos\frac{p\pi m}{N}\right)𝐗_p$$
(10)
and using the model solution eq.(7) we obtain, after a gaussian integration,
$$<e^{i𝐪𝐑_{nm}}>_{\dot{\gamma }}=exp\left\{\frac{2k_BT}{\chi _1}𝐪\left[\underset{p=1}{\overset{\mathrm{}}{}}\frac{\beta _p}{p^2}\left(cos\frac{p\pi n}{N}cos\frac{p\pi m}{N}\right)^2\right]𝐪\right\}$$
(11)
Substituting the expression of $`\beta _p`$ from eq.(8) we get
$$<e^{i𝐪𝐑_{nm}}>_{\dot{\gamma }}=exp\left\{\frac{2k_BT}{\chi _1}\left[S_0(n,m)q^2+\frac{\tau _1}{2}S_1(n,m)𝐪(𝐤+𝐤^T)𝐪+\frac{\tau _1^2}{2}S_2(n,m)𝐪(𝐤𝐤^T)𝐪\right]\right\}$$
(12)
with
$$S_k(n,m)=\underset{p=1}{\overset{\mathrm{}}{}}\frac{1}{p^{2+2k}}\left[cos\frac{p\pi n}{N}cos\frac{p\pi m}{N}\right]^2k=0,1,2$$
(13)
In eq.(12) the first term in the exponent on the r.h.s. is the equilibrium contribution leading to the well known Debye function. The other two terms arise from the presence of the flow. The three series above can be summed analytically as explained in appendix B and the result is
$`S_0(a,b)`$ $`=`$ $`{\displaystyle \frac{\pi ^2}{2}}|ab|`$ (14)
$`S_1(a,b)`$ $`=`$ $`{\displaystyle \frac{\pi ^4}{4}}(ab)^2\left[(a+b)\left(1{\displaystyle \frac{a+b}{2}}\right){\displaystyle \frac{|ab|}{3}}\right]`$ (15)
$`S_2(a,b)`$ $`=`$ $`{\displaystyle \frac{\pi ^6}{240}}(ab)^2\left\{|ab|^3+5(a+b)\left[(a+b)(2+a^2+b^2)3a^22ab3b^2\right]\right\}`$ (16)
where we have defined reduced indices $`a=n/N,b=m/N`$ and $`0(a,b)1`$.
Let us now introduce the reduced shear rate $`\beta `$ as the product of the bare shear rate and the global relaxation time of the chain at equilibrium. To compare the theory with the experiments we define $`\beta `$ as
$$\beta =\frac{M\eta _s[\eta ]\dot{\gamma }}{N_Ak_BT}$$
(17)
where $`M`$ is the molecular weight of the polymer, $`\eta _s`$ the solvent shear viscosity, $`[\eta ]`$ the intrinsic viscosity of the solution and $`N_A`$ the Avogadro number. In the Rouse model the intrinsic viscosity is $`[\eta ]=N_A/M\eta _s(N^2d^2\xi /36)`$ so that
$$\beta =\frac{\dot{\gamma }}{k_BT}\frac{N^2d^2\xi }{36}$$
(18)
which gives $`\tau _1\dot{\gamma }=12\beta /\pi `$. The chain structure factor in terms of the reduced shear rate is finally obtained as
$$\frac{S(𝐪,N,\dot{\gamma })}{N}=_0^1𝑑a_0^1𝑑be^{\frac{Nd^2}{6}\mathrm{𝐪𝐪}^T:𝐟(a,b,\beta )}$$
(19)
where $`𝐟(a,b,\beta )`$ is the following tensor
$$𝐟(a,b,\beta )=|ab|\mathrm{𝟏}+\frac{12\beta }{\pi ^4}\left[S_1(a,b)(\mathrm{\Lambda }+\mathrm{\Lambda }^T)+\frac{12\beta }{\pi ^2}S_2(a,b)(\mathrm{\Lambda }\mathrm{\Lambda }^T)\right]$$
(20)
and $`\mathrm{\Lambda }=\dot{\gamma }^1𝐤`$. Expressions (19) with (20) are our main interest here. The double integral over monomers cannot be performed analytically and we resort to numerical method. Before presenting specific results we want to make some further considerations. Firstly we note that the tensor $`𝐟`$ entering in the structure factor depends on $`N`$ and $`\dot{\gamma }`$ through $`\beta `$ only. This implies the following universality for the structure factor
$$\frac{S(𝐪,N,\dot{\gamma })}{N}=F(R_g𝐪,\beta )$$
(21)
where $`F`$ is a universal scalar function depending on the geometry of the applied flow and $`R_g=\sqrt{Nd^2/6}`$ is the equilibrium radius of gyration for the gaussian chain. Secondly we want to obtain the well known expression for the gyration tensor of the ideal chain under shear flow from which we can derive the orientation and deformation of the chain to be used below. The small $`q`$ expansion to order $`q^2`$ of the structure factor is directly related to the gyration tensor by
$$\frac{S(𝐪,N,\dot{\gamma })}{N}=1\mathrm{𝐪𝐪}^T:𝐆+O(q^4)$$
(22)
When applied to eq.(19) it provides
$$𝐆(N,\beta )=\frac{Nd^2}{6}_0^1𝑑a_0^1𝑑b𝐟(a,b,\beta )=\frac{Nd^2}{6}\left[\frac{1}{3}\mathrm{𝟏}+\frac{2\beta }{15}(\mathrm{\Lambda }+\mathrm{\Lambda }^T)+\frac{16\beta ^2}{105}(\mathrm{\Lambda }\mathrm{\Lambda }^T)\right]$$
(23)
Note that at fixed $`\beta `$ the gyration tensor has the same N dependence as at equilibrium (the scaling exponent is $`\nu =0.5`$). The orientation of the chain, as given by eq.(3), when defined through the gyration tensor is
$$cotg(2\chi _g)=\frac{G_{xx}G_{yy}}{2G_{xy}}=\frac{60\beta }{105}=\frac{\beta }{1.75}$$
(24)
implying that the orientational resistance $`m_g=1.75`$ and $`\beta `$-independent. This is the well known classical result.
Other used quantities in characterizing chains under SF are the deformation ratios as defined in section II. From eq.(23) we get $`\delta G_{xx}=16/35\beta ^2=0.457\beta ^2;\delta G_{yy}=\delta G_{zz}=0`$ as is well known. Note that the above values of the deformation ratios in the laboratory reference frame imply that in the molecular reference frame both elongation and compression directions present deformations as follows
$`\delta G_I`$ $`=`$ $`{\displaystyle \frac{2}{35}}\left[4\beta ^2+7\beta \sqrt{1+\left({\displaystyle \frac{4}{7}}\beta \right)^2}\right]`$ (25)
$`\delta G_{II}`$ $`=`$ $`{\displaystyle \frac{2}{35}}\left[4\beta ^27\beta \sqrt{1+\left({\displaystyle \frac{4}{7}}\beta \right)^2}\right]`$ (26)
with the following limiting behaviours
$`\delta G_I`$ $``$ $`{\displaystyle \frac{16}{35}}\beta ^2\beta >>{\displaystyle \frac{7}{4}}`$ (27)
$`\delta G_{II}`$ $``$ $`{\displaystyle \frac{49}{140}}=0.35\beta >>{\displaystyle \frac{7}{4}}`$ (28)
$`\delta G_I`$ $`=`$ $`\delta G_{II}{\displaystyle \frac{14}{35}}\beta \beta <<{\displaystyle \frac{7}{4}}`$ (29)
Now we go back to our original task, that is the calculation of the chain structure factor under shear flow. As already stated, the double integral over monomers in eq.(19) with eq.(20) cannot be performed analytically because of the exponential character of the integrand (to be compared with the double integral of the exponent which can be solved and provides the gyration tensor as in eq.(23)). We applied standard Romberg’s method to numerically evaluate those integrals in a wide range of $`q`$’s and for given $`\beta `$. In figure 1 we report the chain structure factor in its universal form ($`S(𝐪)(qR_g)^2/3N`$ vs $`qR_g/\sqrt{3}`$), for various values of $`\beta `$ and for $`\widehat{𝐪}`$ along the elongation and the compression directions in the flow plane (molecular reference frame). We note as in the elongation direction the crossover to the high $`q`$’s equilibrium power law behavior strongly depends on the shear rate. The Kratky plot in the compression direction presents a $`\beta `$ dependent overshoot due to the compression of the corresponding gyration tensor component. This effects however saturates very soon (the curve for $`\beta =3.2`$ is already very close to the behavior for $`\beta =10`$ which is indistinguishable from the one at $`\beta =100`$) and the crossover to the high $`q`$’s power law behavior is in practice $`\beta `$ independent.
In comparison with the stretched ideal chain we note the absence of oscillations in the elongation direction of the sheared chain signal. Such oscillations are characteristics of an aligned rod with uniform density of scattering centers. Our present results show definitively that, even at very high shear rates, the sheared chain cannot be thought in terms of such simplified model at variance with Onuki scheme .
Another important characteristics of the structure factor of the stretched chain, related to the presence of tensile blobs, is the Lorentzian shape of the scattering function at intermediate $`q`$’s as discussed by Benoit et al.. In this case one can easily show that, for large $`q`$, the scattering function for $`\widehat{𝐪}`$ along the stretching direction follows the equilibrium Ornstein-Zernike behavior $`S(q)=1/(1+q^2R_g^2/2)`$ if plotted in terms of the effective wave vector $`\stackrel{~}{q}^2=q^2+4/\xi _T^2`$, where $`\xi _T`$ is the tensile blob size. This led Pincus to formulate his scaling for stretched EV chains just changing the scaling exponent from the ideal to the EV value. Under shear flow the chain structure factor at intermediate $`q`$’s in the elongation direction has not a Lorentzian shape, but rather presents a behavior $`S(x)/N=1/(1+x^2/2+a(\beta )x^{\alpha (\beta )})`$ for $`x=qR_g/\sqrt{(}3)1`$, in analogy with the scattering signal from near critical fluids under shear. The fitting function $`a(\beta )`$ goes linearly from $`0`$ to $`73`$ in the interval $`0\beta 100`$, while $`\alpha (\beta )`$ goes almost linearly in $`log(\beta )`$ from $`0.62`$ to $`0.88`$ in the interval $`3.2\beta 100`$. Whether this form for $`S(q_I)`$ has some theoretical foundation remains unclear to us.
## IV Excluded volume chain model and Brownian Dynamics algorithm
The ideal chain model of previous section is only the starting point for a more realistic description of polymeric chains. When the aim is the study of dilute solutions at least two additional effects need to be taken into account namely the excluded volume and the hydrodynamic interactions between monomers. In this section we describe the model we have used to represent the excluded volume effect. In our previous work on stretched chains we have used a model of linear chains with rigid bonds and no angle hindrance. The excluded volume interaction where modeled by hard sphere interactions. The diameter $`\sigma `$ of each monomer was chosen as $`\sigma =0.65`$ in units of the bond length. This model exhibits good solvent scaling and is very suitable for Monte Carlo calculations. Here we deal with non equilibrium systems for which the stationary distribution under shear is unknown and therefore we cannot apply the Monte Carlo sampling procedure. The only way to perform the calculation is by a dynamical technique such as Brownian Dynamics(BD). For this technique, the model of the previous work is not very appropriate because of the strong discontinuity in the interactions. We thus replaced the hard sphere by a Lennard-Jones potential truncated at the minimum and shifted in such a way to have a continuous potential
$`v_{EV}(r)`$ $`=`$ $`v_{LJ}(r)v_{LJ}(r_c)rr_c=2^{1/6}\sigma `$ (30)
$`=`$ $`0r>r_c`$ (31)
where $`v_{LJ}(r)=4ϵ[(\sigma /r)^{12}(\sigma /r)^6]`$. In our calculation $`ϵ`$ was the energy unit and $`\sigma `$ was chosen to be $`0.65`$ as in the previous work. Moreover the imposition of the bond constraints in BD requires an iterative technique such as SHAKE that for long chains is quite demanding in terms of computer time. We therefore prefer to use flexible rather then rigid bonds and represent the bond interaction by an harmonic potential with force constant $`\chi `$ and minimum position $`d`$: $`v_{bond}(r)=\chi (|𝐫|d)^2/2`$. Lennard-Jones interactions between nearest neighbors along the chain were not considered.
The equation of motion of the i-th monomer is
$$\xi \dot{𝐑}_i=𝐅_i+𝐟_i+\xi 𝐤𝐑_i$$
(32)
where $`𝐅_i`$ is the total force on monomer $`i`$ deriving from the potential energy, and the properties of the noise $`𝐟_i`$ are
$`<𝐟_i(t)>`$ $`=`$ $`0`$ (33)
$`<𝐟_n(t)𝐟_j(t^{})>`$ $`=`$ $`2k_BT\delta (tt^{})\delta _{ij}\mathrm{𝟏}`$ (34)
The above equations can be integrated numerically by the simple finite differences scheme
$$𝐑_i(t+h)=[\mathrm{𝟏}+h𝐤]𝐑_i(t)+\frac{h}{\xi }𝐅_i(t)+𝐰_i(t)+O(h^{3/2})$$
(35)
where the moments of the white noise $`𝐰_i`$ are
$`<𝐰_i(t)>`$ $`=`$ $`0`$ (36)
$`<𝐰_i(t)𝐰_j(t^{})>`$ $`=`$ $`2hD_0\delta (tt^{})\delta _{ij}\mathrm{𝟏}`$ (37)
and $`D_0=k_BT/\xi `$ is the single monomer diffusion coefficient. Higher order schemes has been proposed in the literature which allow to take a larger time step but require more operations per time step. A further advantage of the simple first order scheme over higher order ones, is the possibility to use a uniformly rather than normally distributed noise, provided that the zero and second moment are correctly chosen. Indeed, it is easy to show that the error introduced by such a substitution is of order $`h^{3/2}`$ or higher, that is beyond the precision of the scheme itself. Use of uniformly distributed noise rather than a gaussian noise saves about $`50\%`$ of the computer time. In absence of a systematic comparison on the overall efficiency of integration schemes we have preferred to use the simple one above.
In all our calculations we have fixed $`k_BT=1`$ and $`\xi =1`$ which lead to $`D_0=1`$. The bond force constant has been chosen to be $`\chi =40`$ and in order to obtain a unitary average bond length we had to fix $`d=0.945`$.
The presence of stiff springs, ensuring the rigidity of the bonds, and strongly repulsive core interactions, ensuring the self avoidance and avoiding the self crossing of the chain, impose a quite small time step for the stability of the numerical scheme. On the other hand the long relaxation time of polymers and its rapid increase with the number of monomers ($`\tau _NN^{2.2}`$ for excluded volume chains without HI) would require a time step as large as possible. This is the main limitation of any dynamical scheme when dealing with polymers. In the present work we have chosen a time step $`h=0.00025`$ in reduced units. We have tested in some cases that halving the time step does not change the results for the chain structure. Even with such small time step, numerical instabilities occasionally show up. These occur when two particles end a time step at a reciprocal distance considerably smaller then $`\sigma `$. The following time step is then driven by very large repulsive forces which stretch the bonds related to the pair of particle at a distance considerably larger then the minimum of the bond potential and so forth. To cure this occasional pathology we checked at the end of any time step the minimum distance between any pair of particles and if it is found to be smaller then $`0.43\sigma `$ we reject the time step and we proceed the time integration for a time interval $`h`$ in many smaller time steps. When the time interval $`h`$ is reached the original time step is restored. With the chosen value for the original time step we have found that such pathological events happen on average any $`2x10^5`$ time steps and that dividing the time step by ten is always enough to escape from the instability. This is a craft made version of more sophisticated adaptive time step integration schemes which are however more demanding in terms of computer time.
To further minimize the computer time per step we have used a Verlet neighbors list which is very effective with such short range interactions.
## V Simulation results
### A Equilibrium scaling of the model
Before proceeding to study the chain model under shear flow we need to characterize its equilibrium behavior. In particular we need to know at which length the chains start following the known static and dynamic scalings and we need to determine the prefactor in the scaling laws in order to fix the values of the bare shear rate for a given chain length to work at fixed $`\beta `$. In figure 2 we report the radius of gyration $`R_g`$ and the end-to-end distance $`R`$ versus the number of links in the chain $`(N1)`$. The observed scalings are $`R_g=0.40448(N1)^{0.6}`$, and $`R=1.03(N1)^{0.6}`$ if we exclude $`N=9,20`$. Equilibrium structure factors for the various chain lengths are shown in figure 3a. In figure 3b we report the universal plot $`(qR_g)^2S(q)/N`$ versus $`(qR_g)`$. We also add the Debye curve for the ideal chain. We observe that, for $`N30`$, the EV chains follow indeed a universal $`q^{1/3}`$ power law as predicted on the basis of the static scaling and dimensional arguments .
By construction, the chain centre of mass diffusion coefficient is $`D_{cm}=D_0/NN^\nu ^{}`$ where $`\nu ^{}`$ is the dynamic exponent, which provides a characteristic time of the chain $`\tau _N=R_g^2/(6D_{cm})=(0.40448)^2(N1)^{1.2}N/6D_0N^{2\nu +1}`$. Here $`\nu =3/5`$ is the Flory scaling exponent for EV chains, and $`\nu ^{}=1`$ as HI are absent. In order to check the consistency between global and internal dynamics we extracted the characteristic time of the latter by inverting the chain dynamic structure factor $`I(q,t)=S(q,t)/S(q)`$ as we have previously done for MD chains . Indeed scaling and dimensional arguments predict that $`I(t,q,N)`$ be a universal function of $`(tq^x)`$ with $`x=2+\nu ^{}/\nu =11/3`$ in the present case. Figure 4 shows the behavior of $`(tq^{11/3})`$ vs $`q`$ for two values of $`I`$ and for various chain lengths. We observe that for $`N=20`$ the chain does not exhibit the correct scaling of internal times while from $`N=30`$ on, both static and dynamics are in the scaling laws regime. The data we show are results of quite long runs as reported in table I.
### B Chain structure at fixed $`\beta `$
In this subsection we describe the results for the chain structure under shear flow for various chain lengths at the same value of the reduced shear rate $`\beta `$ that we defined as $`\beta =\dot{\gamma }R_g^2/(6D_{cm})`$. We have limited our study to two values of $`\beta `$, namely $`\beta =3.2`$ as in our previous MD study and in several experiments, and $`\beta =10`$. In tables II and III, we report various technical details of our simulations and the results for the orientation and the deformation. We studied chains up to $`N=300`$ which require very long runs (up to $`310^9`$ time steps for the EV case, see table III).
As in our previous MD study of short chains (up to $`N=50`$), we analyze the results in the molecular reference frame on account of the observation that the orientational angle depends on $`N`$ and $`\dot{\gamma }`$ through $`\beta `$ only. This is also suggested by the Rouse model results of section III and by its extension when HI at the level of equilibrium preaveraging are considered . Our new results for EV chains at $`\beta =3.2`$ confirm this expectation within error bars, although short chains seems to be slightly less oriented as shown in figure 5. On the same figure we report the extinction angle related to the birefringence which also is found to be the same within error bars for all $`N`$ at fixed $`\beta `$.
Before discussing the results for our chain model with EV interactions it is instructive to consider the same coarse grained model with the EV interactions switched off ($`ϵ=0`$ in the LJ potential). Comparison of the chain structure factor of such a model with the continuous ideal model result eq.(19) (not shown here) validates the theory itself and the numerical procedure applied to compute the double integral. Moreover data for the orientation and the deformation of this model can be useful in understanding qualitatively the role of the finite extensibility effects due to the stiff harmonic potential which keeps the bond length nearly constant. Data for the orientation and deformations of ideal chains with $`N=100`$ and $`N=300`$ at $`\beta =3.2`$ and $`\beta =10`$ are reported in table II and compared to the Rouse model predictions. At $`\beta =3.2`$, the orientation for all chain lengths is in agreement with the long chain limit while the deformations for $`N=100`$ deviate considerably from the ones for $`N=300`$ (except in the direction II) which are instead in agreement with the Rouse prediction. Interestingly, the main eigenvalue $`G_I`$ of the gyration tensor follows an apparent power law $`N^{1.5}`$ qualitatively similar to the behavior observed for MD data on shorter chains with EV and HI. At $`\beta =10`$ both chains exhibit finite extensibility effects even on the orientation.
Now we turn to the EV results. In table III we collect the data for the orientation and the deformation of our chains. We also report data for the longest chains studied in ref. by molecular dynamics. We note at $`\beta =3.2`$, that the 300-bead chain appears to be slightly more oriented and deformed than the shorter chains, although the noise is quite large. At $`\beta =10`$ instead, the 300-bead chain is considerably more oriented and deformed than the 100-bead chain, showing an important influence of the finite extensibility. Comparison with previous MD data for N=50 shows a general agreement for all quantities.
From these results and those of table II, it appears reasonable to assume that the N=300 results at $`\beta =3.2`$ are representative of the global properties of long chains, namely that the remaining small systematic finite chain effects are masked by statistical errors of a few percents (due to finite statistics). Comparing EV chains to $`\theta `$ chains at the same reduced shear rate, we observe for EV chains a less marked alignment with respect to the flow (larger $`\chi `$ value) and a smaller global deformation, in qualitative agreement with recent LS results .
An important question concerns the universal character of the structure factor under shear flow. It is well known that at equilibrium the quantity $`S(q,N)/N`$ is a universal function of $`qR_g`$ only, for any coarse-grained model and for long enough chains. In the case of the chain stretched at both ends, we have shown that different universalities in terms of $`qR_g`$ hold separately in the longitudinal and the transverse directions. Alternatively one could investigate self-similarities in terms of $`q\sqrt{G_\alpha }`$ where $`G_\alpha `$ are the diagonal components of the gyration tensor. For the chain stretched by its ends, it is well known that fluctuations of the end-to-end distance (and therefore $`G_\alpha `$) at fixed reduced force scale with $`N`$ as the equilibrium radius of gyration and therefore the above two alternatives are equivalent. In our previous study of linear chains under shear flow we observed apparent scaling exponents $`\nu _\alpha `$ for the eigenvalues of the gyration tensor with $`(N1)`$ which were different from the equilibrium one in the flow plane directions while in the neutral (out of plane) direction the equilibrium scaling was preserved within error bars. In the present study on longer chains we find that at $`\beta =3.2`$ the equilibrium scaling $`G_\alpha N^{2\nu }`$ is preserved in the neutral direction ($`\alpha =III`$) and in the in-plane compression direction ($`\alpha =II`$) (see figure 6). In the elongation direction an apparent scaling exponent higher than the equilibrium value is still observed ($`\nu _I0.7`$) although the noise is considerably higher than in the other two principal directions. The results for the chain without EV interactions mentioned above demonstrate that such apparent scaling is related to finite extensibility effects rather than to EV interactions. For this reason we investigate universality in terms of $`qR_g`$.
In figure 7a,b,c we plot, at $`\beta =3.2`$ and for $`30N300`$, $`(qR_g/\sqrt{3})^{5/3}S(q)/N`$ versus $`qR_g/\sqrt{3}`$ for $`\alpha =III,II,I`$ respectively and we compare with the MD data for N=50 at the same $`\beta `$ from ref.. For the latter we have used the observed equilibrium exponent $`\nu _{MD}=0.57`$ rather then the Flory classical value $`\nu _F=0.6`$. Note that in that model the solvent was explicitly considered so that deviations from the mean field value are not contradictory. We observe that, with the exception of $`N=30`$ in the neutral direction, universality is indeed observed within error bars in all directions. A general feature of those curves is the presence of a quite slow crossover back to the equilibrium behavior at large $`q`$’s which would be represented by an horizontal line. That $`N=300`$ is still too short to cover all length scales of the problem is clearly shown by the fact that even for this length the equilibrium behavior is not completely recovered before the chain finite extensibility shows up, a sign that finite chain length effects are still present. While in the neutral direction the effect of the flow is tiny, in the in-plane directions the flow strongly changes the scattering function. We note that the noise level is considerable higher in the elongation direction than in the other directions. Moreover statistical convergence of the results at low $`\beta `$ appears to be considerably slower than at high $`\beta `$ and at equilibrium (see also table III). This is compatible with the possible appearance of new characteristic relaxation times in weak flow as indicated by recent experiments .
Finally we note that the structure factor obtained in our previous MD simulation on a different chain model embedded in an explicit atomic solvent under shear flow , follows very closely in all directions the universal behavior predicted by the present BD simulations. As the former study includes automatically the hydrodynamics interactions while the present calculation do neglect them from the start, it appears that HI play a secondary role on the chain structure under shear flow, once one works at a fixed value of the suitably defined reduced shear rate. Because of the limited length of the largest chain studied by molecular dynamics (N=50), we could not detect on the chain structure factors the slow trend back to the equilibrium behavior at large q’s, which explains our previous analysis in terms of anisotropic scaling laws.
In order to test the “shear” blob hypothesis we report in figures 8,9,10, $`(qR_g)^2S(q)/3N`$ versus $`qR_g/\sqrt{3}`$ for $`N=100`$ and $`N=300`$ at $`\beta =3.2`$ and for $`N=100`$ at equilibrium, and we compare with the ideal chain behavior. In all directions Rouse behavior describes the excluded volume chain up to $`qR_g/\sqrt{3}3`$. Similarly to the stretched chain case , we clearly see in the compression and the neutral directions a crossover from ideal statistics at low $`q`$’s to EV statistics at high $`q`$’s. The inset of EV statistics is at $`qR_g/\sqrt{3}9`$ in the two directions. Note that for $`N=100`$ the EV regime is completely missing and even with $`N=300`$ it is very narrow in the compression direction. In the elongation direction the noise level is quite high in the Rouse (small $`q`$) regime. As stated above, in this direction flow effects extends to higher $`q`$’s than in the other directions, above $`qR_g/\sqrt{3}=9`$. On the other hand, the finite extensibility takes place around $`qR_g/\sqrt{3}=20`$ for $`N=300`$ and therefore the EV regime cannot be observed in this direction for such chain length.
Onuki’s phenomenological model is based on the assumption that below a characteristic length scale $`\xi `$ the effects of the flow are negligible. The length scale separation is dictated by the typical relaxation time of chain density fluctuations at that scale. At a given shear rate $`\dot{\gamma }`$, the crossover occurs at a scale $`\xi =An_c^\nu `$ such that the longest relaxation time of the blob is $`\tau _c=\xi ^2/(6D_c)=\dot{\gamma }^1`$ where $`D_c=D_0/n_c^\nu ^{}`$ is the diffusion of the blob as a whole and $`\nu ^{}`$ is the dynamic scaling exponent. Working at fixed reduced shear rate $`\beta =\tau _N\dot{\gamma }`$, the number of monomers per blob results in $`n_c=N/\beta ^{1/(2\nu +\nu ^{})}`$ so that the number of blobs per chain is $`N_b=N/n_c=\beta ^{1/(2\nu +\nu ^{})}`$ and $`\xi =AN^\nu /\beta ^{\nu /(2\nu +\nu ^{})}=R_g/\beta ^{\nu /(2\nu +\nu ^{})}`$. For our model ($`A=0.40448,\nu =3/5,\nu ^{}=1`$) at $`\beta =3.2`$ we have $`N_b=1.7,n_c=177`$ and $`\xi =9.02`$ for $`N=300`$. The crossover from ideal to EV statistics in figure 8 is observed around $`qR_g/\sqrt{3}=4.0`$ which for $`N=300`$ provides $`q_c=0.565/\xi 2\pi /\xi `$. Therefore the shear blob hypothesis is compatible with our results although a more firm test on the basis of longer chains studied at various values of $`\beta `$ is still missing. To finally prove that our chains are too short for this purpose we show in figure 11 the Kratky plots in the three principal directions for $`N=300`$ and $`\beta =3.2`$ and $`\beta =10.0`$. It is evident that at $`\beta =10.0`$ the finite extensibility effect interferes with the EV regime giving rise to anomalous behaviours.
## VI Conclusions
In order to study the structure of a single chain subjected to a homogeneous and steady shear flow and to elucidate the effects of excluded volume interactions we have computed the scattering function of the continuous Rouse chain model in flow and we have performed BD simulations of a chain model of Fraenkel springs with fully developed EV interactions.
Concerning global properties, our simulations confirm that at a given reduced shear rate, EV chains are less deformed and less oriented along the flow lines than chains at the $`\mathrm{\Theta }`$ point .
We have shown that, even in the ideal case, the chain structure in shear flow cannot be explained by the known behavior of a chain stretched at its ends. In particular we have shown that the structure factor of the Rouse model under SF for the scattering vector lying along the in plane elongation direction is very different from the one of the Rouse chain stretched at its ends with $`q`$ along the stretching direction. It does not present the oscillations typical of permanent extension, but it is a monotonous curve qualitatively similar to the equilibrium Debye curve (see figure 10). In the elongation direction, the onset of the well known $`q^2`$ behavior signaling the static scaling of ideal chains is shifted to higher $`q`$’s than at equilibrium and is strongly dependent on the shear rate. In the compression direction it is at the same location as at equilibrium and it is almost insensible to shear rate.
Adding the EV interactions, our results suggest that the universality of $`S(𝐪,N,\dot{\gamma })`$ in terms of $`x=𝐪R_g`$ and $`\beta `$ only ($`R_g`$ is the equilibrium radius of gyration and $`\beta `$ is computed on the basis of the equilibrium characteristic relaxation time), exhibited by the Rouse model at any flow intensity, is preserved by the EV interactions provided the chains are long enough. Studying chains between 30 and 300 monomers we indeed observe that the eigenvalues of the gyration tensor in the compression (II) and in the neutral (III) directions follow scaling laws with N with the equilibrium Flory exponent. In the elongation direction a larger apparent scaling exponent is instead measured (around 0.7), in agreement with previous MD data which led us to propose an anisotropic scaling picture for intermediate reduced shear rates. From the results for the same chain model without EV interactions we infer that this anomalous scaling arises from finite extensibility effects rather than EV interactions so as the equilibrium scaling should be recovered even in the elongation direction for longer chains. Unfortunately the dynamical character of the BD algorithm and the unfavorable scaling with N of the chain relaxation time does not allow us to study chains longer than $`N=300`$.
The general agreement found between our present data (where HI are absent from the model) and the previous non-equilibrium MD data for short chains (where HI are automatically included) strongly suggests that HI effects play a minor role on the chain structure under shear flow, provided we compare results at the same reduced shear rate, which incorporates implicitly the HI effect through the longest relaxation time of the chain at equilibrium. The EV and the finite extensibility effects are by far the most relevant effects to be taken into account.
For our longest chain at $`\beta =3.2`$, the structure factor in the neutral (III) and in the compression (II) directions exhibit a clear crossover from the ideal behavior at small $`q`$’s, well predicted by the Rouse model, to the isotropic equilibrium EV behavior at high $`q`$’s, characterized by a power law $`S(q)q^{1/\nu }`$ with $`\nu =\nu _{eq}=3/5`$. This behavior is usually related to the presence of the so called “blobs”, in the present case“shear blobs”. Blobs were detected in the chain structure factor either experimentally for semi-dilute solutions or by Monte Carlo simulations for a single EV chain stretched at its ends . The presence of blobs in a chain under shear flow has been predicted theoretically but had never been observed so far, neither by SANS nor by simulation. The results of our BD simulations are compatible with the existence of blobs (see section V B) but, because of finite chain size effects, a quantitative test of this hypothesis would require chain lengths larger than N=300. It is interesting to note that the largest chain considered here has a size which corresponds to a polystyrene chain of the order of 300 Kuhn segments of about $`18\AA `$, which yields a molecular mass of 300.000 Daltons .
The present analysis should be tested by scattering experiments on dilute solutions for chains in good solvent subjected to steady shear flow. It would be interesting to study at low $`qR_g`$, either by Light or by Neutron scattering, whether the Rouse model indeed describes the global orientation and deformation of chains for large deformations, i.e. at $`\beta >>1`$. Remaining discrepancies could be interpreted as a direct measure of FE and HI influence on these structural properties far from the equilibrium structure. Ultimately, the most spectacular experimental result we seek for remains a SANS study of a suitable “polymer-good solvent” pair which would allow us to follow the form factor of a polymer of a few million Daltons under shear flow at large $`\beta `$ over a broad $`q`$ regime. This would allow a direct observation of the crossover from ideal chain to EV chain statistics across a crossover diffusion vector $`q_c`$ characterizing the shear blob dimension in a particular direction.
Concerning future simulation work, a more quantitative test of HI effects on the chain structure is in order and could be performed by BD on short chains.
## VII Acknowledgments
We thank M. Baus, B. Dünweg, P. Lindner and R.Winkler for useful discussions and suggestions.
## A Ideal chain under steady shear flow
In this appendix we derive the steady state distribution function of the continuous gaussian chain model under an homogenous shear flow, eqs. (6), (7) and (8).
The dynamics of the model under a generic flow field, specified by the velocity gradient tensor $`𝐤`$, is defined by the Langevin equation
$$\xi \frac{𝐑_n(t)}{t}=\chi \frac{^2𝐑_n(t)}{n^2}+𝐟_n(t)+\xi 𝐤(t)𝐑_n(t)$$
(A1)
with the boundary conditions
$$\left[\frac{𝐑_n(t)}{n}\right]_{n=0}=\left[\frac{𝐑_n(t)}{n}\right]_{n=N}=0$$
(A2)
$`𝐟_n(t)`$ is the gaussian random force acting on the chain and is defined by
$`<𝐟_n(t)>`$ $`=`$ $`0`$ (A3)
$`<𝐟_n(t)𝐟_m(t^{})>`$ $`=`$ $`2k_BT\delta (tt^{})\delta (nm)\mathrm{𝟏}`$ (A4)
where $`<\mathrm{}>`$ indicates an average over the distribution of the random noise and $`\mathrm{𝟏}`$ is the unit tensor. As it is well known the above Langevin dynamics corresponds to the following diffusion equation (Smoluchowsky) for the distribution function in the chain configurational space
$$\frac{\mathrm{\Psi }}{t}=_0^N𝑑n_{𝐑_n}\left[\dot{𝐑}_n\mathrm{\Psi }\right]$$
(A5)
with the dynamics of monomers given by
$$\dot{𝐑}_n=\frac{\chi }{\xi }\frac{^2𝐑_n}{n^2}\frac{k_BT}{\xi }_{𝐑_n}[ln\mathrm{\Psi }]+𝐤𝐑_n$$
(A6)
Let us introduce now the normal modes defined as
$$𝐗_p(t)=\frac{1}{N}_0^N𝑑ncos\left(\frac{p\pi n}{N}\right)𝐑_n(t)p=0,1,2,\mathrm{}.$$
(A7)
in terms of which the chain positions are given as
$$𝐑_n(t)=𝐗_0(t)+2\underset{p=1}{\overset{\mathrm{}}{}}cos\left(\frac{p\pi n}{N}\right)𝐗_p(t)0nN$$
(A8)
The completeness relations for the normal modes basis are
$`\delta _{pq}={\displaystyle \frac{2}{N}}{\displaystyle _0^N}𝑑ncos\left({\displaystyle \frac{p\pi n}{N}}\right)cos\left({\displaystyle \frac{q\pi n}{N}}\right)`$ (A9)
$`N\delta (nm)=1+2{\displaystyle \underset{p=1}{\overset{\mathrm{}}{}}}cos\left({\displaystyle \frac{p\pi n}{N}}\right)cos\left({\displaystyle \frac{p\pi m}{N}}\right)`$ (A10)
Note that, for homogeneous flows which we consider here, the normal modes of the dynamics are the same as at equilibrium.
The Langevin equations for the normal modes are
$$\xi _p_t𝐗_p(t)=\chi _p𝐗_p(t)+\stackrel{~}{𝐟}_p(t)+\xi _p𝐤𝐗_p(t)$$
(A11)
with the following definition of the symbols
$`\xi _0`$ $`=`$ $`N\xi ,\xi _p=2N\xi p=1,2,3\mathrm{}`$ (A12)
$`\chi _p`$ $`=`$ $`{\displaystyle \frac{6\pi ^2k_BT}{Nd^2}}p^2={\displaystyle \frac{2\pi ^2\chi }{N}}p^2=\chi _1p^2`$ (A13)
$`<\stackrel{~}{𝐟}_p(t)>`$ $`=`$ $`0,<\stackrel{~}{𝐟}_p(t)\stackrel{~}{𝐟}_q(t^{})>=2\xi _pk_BT\delta _{pq}\delta (tt^{})\mathrm{𝟏}`$ (A14)
The distribution function of normal modes is factorized
$$\widehat{\mathrm{\Psi }}(\{𝐗\},t)=\underset{p=0}{\overset{\mathrm{}}{}}\widehat{\psi }_p(𝐗_p,t)$$
(A15)
and one can easily derive the following diffusion equation for the $`pth`$ mode ($`p>0`$)
$$_t\widehat{\psi }_p(𝐗_p,t)=_{𝐗_p}\left[(𝐤𝐗_p\frac{\chi _p}{\xi _p}𝐗_p)\widehat{\psi }_p\frac{k_BT}{2\xi }_{𝐗_p}\widehat{\psi }_p\right]$$
(A16)
The center of mass motion is instead described by the diffusion equation for the zero mode
$$_t\widehat{\psi }_0(𝐗_0,t)=_{𝐗_0}\left[𝐤𝐗_0\widehat{\psi }_0\frac{k_BT}{2\xi }_{𝐗_0}\widehat{\psi }_0\right]$$
(A17)
The solution of eqs.(A16) is a multivariate guassian
$$\widehat{\psi }_p(𝐗_p,t)=\left(\frac{1}{(2\pi )^3det[\alpha _p(t)]}\right)^{1/2}e^{\frac{1}{2}𝐗_p\alpha _p^1(t)𝐗_p}$$
(A18)
where the time dependent covariance matrix $`\alpha _p(t)`$ is given by (see eq.(7.161) in ref.)
$$\alpha _p(t)=_{\mathrm{}}^t𝑑t^{}\frac{2k_BT}{\xi _p}𝐄(t,t^{})𝐄^T(t,t^{})e^{\frac{2p^2}{\tau _1}(tt^{})}$$
(A19)
Here $`\tau _1=\xi _1/\chi _1`$ is the relaxation time of the first normal mode and the tensor $`𝐄(t,t^{})`$ is the deformation tensor of the external field defined by
$$𝐄(t,t^{})=\mathrm{𝟏}+_t^{}^t𝑑t^{\prime \prime }𝐤(t^{\prime \prime })$$
(A20)
For the steady shear flow to which we are interested in this work we have
$$𝐤=\left(\begin{array}{ccc}0& \dot{\gamma }& 0\\ 0& 0& 0\\ 0& 0& 0\end{array}\right)$$
(A21)
$$𝐄(t,t^{})=\left(\begin{array}{ccc}1& \dot{\gamma }(tt^{})& 0\\ 0& 1& 0\\ 0& 0& 1\end{array}\right)$$
(A22)
so that the covariance matrices $`\alpha _p`$ results time independent and takes the form
$$\alpha _p=\frac{k_BT}{\chi _p}\left[\mathrm{𝟏}+\frac{\tau _1}{2p^2}(𝐤+𝐤^T)+2\left(\frac{\tau _1}{2p^2}\right)^2𝐤𝐤^T\right]=\frac{k_BT}{\chi _1}\left(\frac{\beta _p}{p^2}\right)$$
(A23)
where we have defined the matrix $`\beta _p`$ through the second equivalence (see eq.(8)). Noting that
$$det(\alpha _p)=\left(\frac{k_BT}{\chi _1p^2}\right)^3\left[1+\dot{\gamma }^2\left(\frac{\tau _1}{2p^2}\right)^2\right]$$
(A24)
the distribution function of the $`pth`$ normal mode under steady shear flow results in eq.(7). It is easy to check that eq.(7) reduces to the standard equilibrium distribution for $`\dot{\gamma }=0`$.
## B Derivation of the sums
In this appendix we outline the method to compute the sums in eqs. (13)
$$S_k(a,b)=\underset{p=1}{\overset{\mathrm{}}{}}\frac{1}{p^{2+2k}}\left[cos(p\pi a)cos(p\pi b)\right]^2k=0,1,2$$
(B1)
where $`0(a,b)1`$. It is evident that $`S_k(a,b)=S_k(b,a)0`$ and the equality holds for $`a=b`$ only. Developing the square inside the sum in eq.(B1) and using standard relations among trigonometric functions we obtain
$$S_k(a,b)=\sigma _k(0)+\frac{1}{2}\left[\sigma _k(a)+\sigma _k(b)\right]\sigma _k\left(\frac{a+b}{2}\right)\sigma _k\left(\frac{ab}{2}\right)$$
(B2)
with
$$\sigma _k(x)=\underset{p=1}{\overset{\mathrm{}}{}}\frac{cos(2\pi px)}{p^{2+2k}}$$
(B3)
so that the problem reduces to compute $`\sigma _k(x)`$ for $`k=0,1,2`$. Note that $`\sigma _k(x)`$ is an even and periodic function of $`x`$ with unitary period. Eq.(B3) expresses the Fourier series a function $`f_k(x)`$ with these properties, after removing the $`p=0`$ term. Therefore we have to seek for a function $`f_k(x)`$ that satisfies the following conditions
$`2p^{2+2k}{\displaystyle _0^1}𝑑xf_k(x)cos(2\pi px)`$ $`=`$ $`1p1`$ (B4)
$`{\displaystyle _0^1}𝑑xf_k(x)sin(2\pi px)`$ $`=`$ $`0`$ (B5)
Applying repeatedly integration by parts we easily obtain the following identities
$`{\displaystyle _0^1}𝑑xf_k(x)cos(2\pi px)`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi p)^2}}{\displaystyle _0^1}𝑑xf_k^{}(x)_xcos(2\pi px)`$ (B6)
$`=`$ $`{\displaystyle \frac{1}{(2\pi p)^4}}{\displaystyle _0^1}𝑑xf_k^{\prime \prime \prime }(x)_xcos(2\pi px)`$ (B7)
$`=`$ $`{\displaystyle \frac{1}{(2\pi p)^6}}{\displaystyle _0^1}𝑑xf_k^𝚟(x)_xcos(2\pi px)`$ (B8)
having imposed the boundary conditions ($`f_k(0)=f_k(1)`$) for the function and for its derivatives up to order five.
Let us now consider the three $`k`$ values separately. For $`k=0`$ we have
$`1`$ $`=`$ $`{\displaystyle \frac{1}{2\pi ^2}}{\displaystyle _0^1}𝑑xf_0^{}(x)_xcos(2\pi px)`$ (B9)
$`0`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi p)^2}}{\displaystyle _0^1}𝑑xf_0^{}(x)_xsin(2\pi px)p1`$ (B10)
$`f_0(0)`$ $`=`$ $`f_0(1)`$ (B11)
where we have used the first of the equalities in eq.(B6) and a similar one for the sinus function. These relations can be satisfied by choosing $`f_0(x)=c_0+c_1x+c_2x^2`$ with $`c_0=0,c_1=\pi ^2`$ and $`c_2=\pi ^2`$ which provides $`f_0(x)=\pi ^2(x^2x)`$. The function $`\sigma _0(x)`$ is then
$$\sigma _0(x)=f_0(x)_0^1𝑑xf_0(x)=\pi ^2(x^2x\frac{1}{6})$$
(B12)
and, after eq.(B2)
$$S_0(a,b)=\frac{\pi ^2}{2}|ab|$$
(B13)
which is the wanted result for $`k=0`$.
For $`k=1`$ we must solve the problem
$`1`$ $`=`$ $`{\displaystyle \frac{1}{8\pi ^4}}{\displaystyle _0^1}𝑑xf_1^{\prime \prime \prime }(x)_xcos(2\pi px)`$ (B14)
$`0`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi p)^4}}{\displaystyle _0^1}𝑑xf_1^{\prime \prime \prime }(x)_xsin(2\pi px)p1`$ (B15)
$`f_1(0)`$ $`=`$ $`f_1(1)`$ (B16)
where the second equalities in eq.(B6) has been used. These relations can be satisfied choosing $`f_1^{\prime \prime \prime }(x)=c_3+c_4x`$ with $`c_4=8\pi ^4`$ and $`c_3=4\pi ^4`$. Integrating three times over $`x`$ and imposing the boundary conditions on $`f_1^{\prime \prime },f_1^{}`$ and $`f_1`$ we obtain $`f_1(x)=\pi ^4(2x^3x^2x^4)/3`$ which provides
$$\sigma _1(x)=f_1(x)_0^1𝑑xf_1(x)=\frac{\pi ^4}{3}(\frac{1}{30}x^2+2x^3x^4)$$
(B17)
Using eq.(B2) we finally obtain
$$S_1(a,b)=\frac{\pi ^4}{4}(ab)^2\left[(a+b)\left(1\frac{a+b}{2}\right)\frac{|ab|}{3}\right]$$
(B18)
In the case $`k=2`$ the relation to be satisfied are
$`1`$ $`=`$ $`{\displaystyle \frac{2}{(2\pi )^6}}{\displaystyle _0^1}𝑑xf_1^𝚟(x)_xcos(2\pi px)`$ (B19)
$`0`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi p)^6}}{\displaystyle _0^1}𝑑xf_1^𝚟(x)_xsin(2\pi px)p1`$ (B20)
$`f_2(0)`$ $`=`$ $`f_2(1)`$ (B21)
A procedure similar to the previous cases provides $`f_2(x)=\pi ^6x^2(2x^46x^3+5x^21)/45`$ and therefore
$$\sigma _2(x)=\frac{\pi ^6}{45}\left[\frac{5}{7}+x^2(2x^46x^3+5x^21)\right]$$
(B22)
and
$$S_2(a,b)=\frac{\pi ^6}{240}(ab)^2\left\{|ab|^3+5(a+b)\left[(a+b)(2+a^2+b^2)3a^22ab3b^2\right]\right\}$$
(B23)
Corresponding author:
Carlo Pierleoni
Physics Department, University of L’Aquila
Via Vetoio, Coppito 67010-L’Aquila (Italy)
phone: +39+0862+433056
fax: +39+0862+433033
email: carlo.pierleoni@aquila.infn.it
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# Fractal Dimension of Higher-Dimensional Chaotic Repellors
## 1 Introduction
In addition to chaotic attractors, nonattracting chaotic sets (also called chaotic saddles or chaotic repellors) are also of great practical importance. In particular, such sets arise in the consideration of chaotic scattering, boundaries between basins of attraction, and chaotic transients. If a cloud of initial conditions is sprinkled in a bounded region including a nonattracting chaotic set, the orbits originating at these points eventually leave the vicinity of the set, and there is a characteristic escape time, $`\tau `$, such that, at late time, the fraction of the cloud still in the region decays exponentially at the rate $`\tau ^1`$.
The primary focus of this paper will be on studying the fractal dimension of nonattracting chaotic sets and their stable and unstable manifolds. Fractal dimension is of basic interest as a means of characterizing the geometric complexity of chaotic sets. In addition, a knowledge of the fractal dimension can, in some situations, provide quantitative information that is of potential practical use. For example, in the case of boundaries between different basins, the basin boundary is typically the stable manifold of a nonattracting chaotic set, and knowledge of the stable manifold’s box-counting dimension (also called its capacity) quantifies the degree to which uncertainties in initial conditions result in errors in predicting the type of long-term motion that results (e.g., which attractor is approached; see Sec. 5 and Ref. ). Our focus is on obtaining the information dimension of a suitable “natural measure” $`\mu `$ lying on the chaotic set (see Sec. 2 for definitions of the natural measures for nonattracting chaotic sets and their stable and unstable manifolds). The information dimension is a member of a one parameter ($`q`$) class of dimension definitions given by
$$D_q=\underset{ϵ0}{lim}\frac{1}{1q}\frac{\mathrm{ln}\mu _i^q}{\mathrm{ln}(1/ϵ)},$$
(1)
where $`q`$ is a real index, $`ϵ`$ is the grid spacing for a $`d`$-dimensional rectangular grid dividing the $`d`$-dimensional state space of the system, and $`\mu _i`$ is the natural measure of the $`i^{\text{th}}`$ grid cube. The box-counting dimension is given by (1) with $`q=0`$, and the information dimension is given by taking the limit $`q1`$ in (1),
$$D_1=\underset{ϵ0}{lim}\frac{I(ϵ)}{\mathrm{ln}(1/ϵ)},I(ϵ)=\underset{i}{}\mu _i\mathrm{ln}(1/\mu _i).$$
(2)
In general, the information dimension is a lower bound on the box-counting dimension, $`D_1D_0`$. In practice, in cases where $`D_1`$ and $`D_0`$ have been determined for chaotic sets, it is often found that their values are very close.
We are specifically concerned with investigating formulae conjectured in Ref. that give the information dimensions for the nonattracting chaotic set and its stable and unstable manifolds in terms of Lyapunov exponents and the decay time, $`\tau `$. These formulae generalize previous results for nonattracting chaotic sets of two-dimensional maps with one positive and one negative Lyapunov exponent , and for Hamiltonian systems of arbitrary dimensionality . In turn, these past results for nonattracting chaotic sets were motivated by the Kaplan-Yorke conjecture which gives the information dimension of a chaotic attractor in terms of its Lyapunov exponents . A rigorous result for the information dimension of an ergodic invariant chaotic set of a two-dimensional diffeomorphism has been given by L.-S. Young , and this result supports the Kaplan-Yorke conjecture for attractors and the two-dimensional map results of Refs. and for nonattracting chaotic sets. Furthermore, Ledrappier has proven that the conjectured Kaplan-Yorke dimension formula is an upper bound on the information dimension of chaotic attractors. Soluble examples for typical systems and numerical experiments tend to support the Kaplan-Yorke result. On the other hand, some specific examples violating the Kaplan-Yorke formula also exist. Thus, in the Kaplan-Yorke conjecture they claim that their formula for the dimension applies for “typical” systems. That is, if a specific system violates the formula then by making a “typical” arbitrarily small change to the system we create a new system for which the Kaplan-Yorke formula holds (this is discussed in Refs. and ). Specific systems violating the conjecture are called *atypical*, and, by Ledrappier’s result, have dimensions below the Kaplan-Yorke prediction. The conjectured results of Ref. for nonattracting chaotic sets are also of this character in that they are also claimed to hold only for “typical” systems. Thus, the issue of what constitutes a typical system (or a typical perturbation of an atypical system) is central. At present there is no rigorous formulation of “typical” for this purpose. For this reason it is important to address this question through examples. One of the purposes of this paper is to do just that. More generally, although many examples, both analytical and numerical, exist which support the dimension formula for chaotic attractors, for the case of nonattracting chaotic sets there are virtually no supporting examples (except for the case of two-dimensional maps with one positive and one negative Lyapunov exponent ). This paper provides such examples.
In Section 2 we review the results of Ref. . We define the *natural transient measure* on nonattracting chaotic sets and on their stable and unstable manifolds, and give the conjectured information dimension formulae for these measures.
In Section 3 we introduce a noninvertible two-dimensional map with two expanding (positive) Lyapunov exponents. This map is simple enough that it can be fully analyzed. This analysis is especially useful in elucidating certain aspects of the general problem, particularly in relation to the natural transient measure and the role of fluctuation in the finite time Lyapunov exponents.
In Section 4 we introduce our main example, a hard walled billiard chaotic scattering system in three spatial dimensions. For a particular symmetric configuration of the scatterer, we find (Secs. 5 and 6) that the situation is atypical (i.e., the conjectured formula does not apply), but that deviation of the configuration from the symmetric case immediately restores the validity of the conjectured dimension formula.
Section 5 presents numerical results for the dimension in the atypical (symmetric) and typical (asymmetric) cases obtaining good agreement with theoretical results in both cases.
Section 6 gives a theory for the structure of the nonattracting chaotic invariant set present in examples modeling our billiard scatterer. It is shown that the stable manifold of the invariant set is a continuous, nowhere-differentiable surface in both the typical and the atypical cases. Furthermore, we are able to obtain an explicit formula for the dimension in the atypical case (this formula is compared to numerical computations in Sec. 5).
In Sec. 7.1 we present a derivation of the dimension formulae of Sec. 2. Although these formulae have been previously derived in Ref. , the previous derivation assumed a specific type of mapping. The situation envisioned in Ref. does not apply to our billiard example. Thus, we are motivated to provide more general arguments for the applicability of the formulae in Sec. 2. In Sec. 7.2 we derive a formula for the dimension of the stable manifold of the chaotic repellor found in the typical (asymmetric) case. This derivation shows explicitly where the derivation of the typical formula fails for this nontrival atypical system.
In Sec. 8 we present conclusions and further discussion. In particular, we point out that the continuous, nowhere-differentiable surface forming the stable manifold of our billiard example may be amenable to experimental observation in a situation where the billiard walls are mirrors and the orbits are light rays. In general, we comment that we believe that such optical billiard experiments provide a particularly convenient arena for the experimental investigation of different types of fractal basin boundaries. This has so far proven to be very difficult in other experimental settings.
## 2 Dimension Formulae
A *chaotic* *saddle*, $`\mathrm{\Lambda }`$, is a *nonattracting*, ergodic, *invariant* set. By *invariant* we mean that all forward and reverse time evolutions of points in $`\mathrm{\Lambda }`$ are also in $`\mathrm{\Lambda }`$. The *stable* *manifold* of $`\mathrm{\Lambda }`$ is the set of all initial conditions which converge to *$`\mathrm{\Lambda }`$* upon forward time evolution. The *unstable* *manifold* of $`\mathrm{\Lambda }`$ is the set of all initial conditions which converge to $`\mathrm{\Lambda }`$ upon reverse time evolution. We say $`\mathrm{\Lambda }`$ is *nonattracting* if it does not completely contain its unstable manifold. In such a case there are points not in $`\mathrm{\Lambda }`$ that converge to it on backwards iteration.
To define the characteristic escape time, $`\tau `$, first define a bounded region, $`R`$, which contains $`\mathrm{\Lambda }`$ and no other chaotic saddle. Uniformly sprinkle a large number, $`N(0)`$, of initial conditions in $`R`$. (In this section we take the dynamical system to be a discrete time system, i.e., a map.) Iterate the sprinkled initial conditions forward $`n1`$ times and discard all orbits which are no longer in $`R`$. Denote the remaining number of orbits $`N(n)`$. We define $`\tau `$ as
$$e^{n/\tau }\frac{N(n)}{N(0)},$$
(3)
or, more formally, $`\tau =lim_n\mathrm{}lim_{N(0)\mathrm{}}\mathrm{ln}[N(0)/N(n)]/n`$. The Lyapunov exponents are defined with respect to the *natural transient measure* of the chaotic saddle . This measure is defined on an open set $`CR`$ as
$$\mu (C)=\underset{n\mathrm{}}{lim}\underset{N(0)\mathrm{}}{lim}\frac{N(\xi n,n,C)}{N(n)},$$
(4)
where $`0<\xi <1`$, and $`N(m,n,C)`$ is the number of sprinkled orbits still in $`R`$ at time $`n`$ that are also in $`C`$ at the earlier time $`m<n`$. The above definition of $`\mu (C)`$ is presumed to be independent of the choice of $`\xi `$ as long as $`0<\xi <1`$ (e.g., $`\xi =1/2`$ will do).
We take the system to be $`M`$-dimensional with $`U`$ positive and $`S`$ negative Lyapunov exponents measured with respect to $`\mu `$ (where $`U+S=M`$) which we label according to the convention,
$$h_U^+h_{U1}^+\mathrm{}h_1^+>0>h_1^{}\mathrm{}h_{S1}^{}h_S^{}.$$
Following we define a forward entropy,
$$H=\underset{i=1}{\overset{U}{}}h_i^+\tau ^1.$$
We now define a natural transient measure $`\mu _S`$ on the stable manifold and a natural transient measure $`\mu _U`$ on the unstable manifold. Using the notation of Eq. (4),
$$\mu _S(C)=\underset{n\mathrm{}}{lim}\underset{N(0)\mathrm{}}{lim}\frac{N(0,n,C)}{N(n)},$$
(5)
$$\mu _U(C)=\underset{n\mathrm{}}{lim}\underset{N(0)\mathrm{}}{lim}\frac{N(n,n,C)}{N(n)}.$$
(6)
Thus, considering the $`N(n)`$ orbits that remain in $`R`$ up to time $`n`$, the fraction of those orbits that initially started in $`C`$ gives $`\mu _S(C)`$, and the fraction that end up in $`C`$ at the final time $`n`$ gives $`\mu _U(C)`$. We use the measure (4), (5), (6) to define the information dimensions of the invariant set, the stable manifold, and the unstable manifold, respectively.
According to Ref. , the dimension of the unstable manifold is then
$$D_U=U+I+\frac{H(h_1^{}+\mathrm{}+h_I^{})}{h_{I+1}^{}},$$
(7)
where $`I`$ is defined by
$$h_1^{}+\mathrm{}+h_I^{}+h_{I+1}^{}Hh_1^{}+\mathrm{}+h_I^{}.$$
The dimension of the stable manifold is
$$D_S=S+J+\frac{H(h_1^++\mathrm{}+h_J^+)}{h_{J+1}^+},$$
(8)
where $`J`$ is defined by
$$h_1^++\mathrm{}+h_J^++h_{J+1}^+Hh_1^++\mathrm{}+h_J^+.$$
Considering the chaotic saddle to be the (generic) intersection of its stable and unstable manifolds, the generic intersection formula gives the dimension of the saddle,
$$D_\mathrm{\Lambda }=D_U+D_SM.$$
(9)
It is of interest to discuss some special cases of Eqs. (7)–(9). In the case of a chaotic attractor, the invariant set is the attractor itself, the stable manifold is the basin of attraction, and we identify the unstable manifold with the attractor. Thus $`D_S=M`$ and $`D_\mathrm{\Lambda }=D_U`$. Since points near the attractor never leave, we have $`\tau =\mathrm{}`$. Equation (7) then yields the Kaplan-Yorke formula ,
$$D_\mathrm{\Lambda }=U+I+\frac{(h_1^++\mathrm{}+h_U^+)(h_1^{}+\mathrm{}+h_I^{})}{h_{I+1}^{}},$$
(10)
where $`I`$ is the largest integer for which $`(h_1^++\mathrm{}+h_U^+)(h_1^{}+\mathrm{}+h_I^{})`$ is positive.
In the case of a two-dimensional map with one positive Lyapunov exponent $`h_1^+`$ and one negative Lyapunov exponent $`h_1^{}`$ with the exponents satisfying $`h_1^+h_1^{}1/\tau 0`$, Eqs. (7) and (8) give the result of Ref. and ,
$$D_U=1+\frac{h_1^+1/\tau }{h_1^{}},$$
$$D_S=1+\frac{h_1^+1/\tau }{h_1^+}.$$
Another case is that of a nonattracting chaotic invariant set of a one-dimensional map. In this case $`S=0`$ and $`U=1`$. The unstable manifold of the invariant set has dimension one, $`D_U=1`$. Recalling the definition of the stable manifold as the set of points that approach the invariant set as time increases, we can identify the stable manifold with the invariant set itself. This is because points in the neighborhood of the invariant set are repelled by it unless they lie precisely on the invariant set. Thus, $`D_S=D_\mathrm{\Lambda }`$, and from (8) and (9) we have $`D_S=D_\mathrm{\Lambda }=H/h_1^+`$, where $`H=h_1^+1/\tau `$.
Still another simple situation is the case of a two-dimensional map with two positive Lyapunov exponents. This case is particularly interesting because we will be able to use it (Sec. 3) to gain understanding of the nature of the natural measure whose dimension we are calculating. In this case $`U=2`$ and $`S=0`$. Thus $`D_U=2`$ and $`D_S=D_\mathrm{\Lambda }`$. There are two cases \[corresponding to $`J=0`$ and $`J=1`$ in Eq. (7)\]. For $`h_2^+\tau 1`$, we have that $`D_S=D_\mathrm{\Lambda }`$ is between zero and one,
$$D_S=D_\mathrm{\Lambda }=1+\frac{h_2^+}{h_1^+}\frac{1}{h_1^+\tau }.$$
(11)
For $`h_2^+\tau 1`$, we have that $`D_S=D_\mathrm{\Lambda }`$ is between one and two,
$$D_S=D_\mathrm{\Lambda }=2\frac{1}{h_2^+\tau }.$$
(12)
In the next section we will be concerned with testing and illustrating Eqs. (11) and (12) by use of a simple model.
## 3 Illustrative Expanding Two-Dimensional Map Model
We consider the following example,
$$x_{n+1}=2x_n\text{modulo}\mathrm{\hspace{0.17em}1},$$
(13)
$$y_{n+1}=\lambda (x_n)y_n+\frac{\eta }{2\pi }\mathrm{sin}(2\pi x_n),$$
(14)
where $`\lambda (x)>1`$, and the map is defined on the cylinder $`\mathrm{}y+\mathrm{}`$, $`1x0`$, with $`x`$ regarded as angle-like. We take $`\lambda (x)`$ to be the piecewise constant function,
$$\lambda (x)=\{\begin{array}{cc}\lambda _1\hfill & 0<x<1/2,\hfill \\ \lambda _2\hfill & 1/2<x<1,\hfill \end{array}$$
(15)
and, without loss of generality, we assume $`\lambda _1\lambda _2`$. \[Later, in Sec. 3.5, we will consider the problem with a general function $`\lambda (x)`$ and a general chaotic map $`x_{n+1}=M(x_n)`$ replacing (13), but for now we focus on (13)–(15).\]
For this map, almost every initial condition generates an orbit that either tends toward $`y=+\mathrm{}`$ or toward $`y=\mathrm{}`$. Figure 1 shows the regions where initial conditions generate these two outcomes, with the black (white) region corresponding to orbits that tend toward $`y=\mathrm{}`$ ($`y=+\mathrm{}`$). Initial conditions on the border of these two regions stay on the border forever. Thus, the border is an invariant set. It is also ergodic by virtue of the ergodicity of the map $`x_{n+1}=2x_nmod\mathrm{\hspace{0.17em}1}`$. We wish to apply Eqs. (11) and (12) to this invariant set and its natural measure.
The Jacobian matrix for our model is
$$𝒥(x)=\left[\begin{array}{cc}2\hfill & 0\hfill \\ \eta \mathrm{cos}2\pi x\hfill & \lambda (x)\hfill \end{array}\right].$$
Thus, for an ergodic invariant measure of the map, the two Lyapunov exponents are
$$h_a=p\mathrm{ln}\lambda _1+(1p)\mathrm{ln}\lambda _2$$
(16)
and
$$h_b=\mathrm{ln}2,$$
where $`p`$ is the measure of the region $`x<1/2`$. To find $`h_a`$ we thus need to know the measure of the invariant set. The measure we are concerned with is the natural transient measure introduced in Sec. 2.
### 3.1 The Decay Time and the Natural Measure
Consider a vertical line segment of length $`\mathrm{}_0`$ whose $`x`$ coordinate is $`x_0`$ and whose center is at $`y=y_0`$. After one iterate of the map (13)–(15), this line segment will have length $`\mathrm{}_1=\lambda (x_0)\mathrm{}_0`$ and be located at $`x=x_1`$ with its center at $`y=y_1`$, where $`(x_1,y_1)`$ are the iterates of $`(x_0,y_0)`$ using the map (13)–(15). Thus we see that vertical line segments are expanded by the multiplicative factor $`\lambda (x)\lambda _1>1`$. Now consider the strip, $`KyK`$, and sprinkle many initial conditions uniformly in this region with density $`\rho _0`$. A vertical line segment, $`x=x_0`$, $`KyK`$, iterates to $`x=x_1`$ and with its center at $`y_1=(\eta /2\pi )\mathrm{sin}2\pi x_0`$. We choose $`K>(\eta /2\pi )(\lambda _11)^1`$ so that the iterated line segment spans the strip $`KyK`$. After one iterate, the density will still be uniform in the strip: The region $`x<1/2`$ $`(x>1/2)`$, $`KyK`$, is expanded uniformly vertically by $`\lambda _1`$ $`(\lambda _2)`$ and horizontally by $`2`$. Thus, after one iterate, the new density in the strip is $`\rho _1=[(\lambda _1^1+\lambda _2^1)/2]\rho _0`$, and, after $`n`$ iterates, we have
$$\rho _n=[(\lambda _1^1+\lambda _2^1)/2]^n\rho _0.$$
Hence the exponential decay time for the number of orbits remaining in the strip is
$$\frac{1}{\tau }=\mathrm{ln}\left[\frac{1}{2}(\frac{1}{\lambda _1}+\frac{1}{\lambda _2})\right]^1.$$
(17)
To find the natural stable manifold transient measure of any $`x`$-interval $`s_m^{(n)}=[m/2^n,(m+1)/2^n]`$, where $`m=0,1,\mathrm{},2^n1`$, we ask what fraction of the orbits that were originally sprinkled in the strip and are still in the strip at time $`n`$ started in this interval. Let $`s_m^{(n)}`$ experience $`n_1(m)`$ vertical stretches by $`\lambda _1`$ and $`n_2(m)=nn_1(m)`$ vertical stretches by $`\lambda _2`$. Then the initial subregion of the $`s_m^{(n)}`$ still in the strip after $`n`$ iterates has vertical height $`K\lambda _1^{n_1(m)}\lambda _2^{n_2(m)}`$. Hence the natural measure of $`s_m^{(n)}`$ is
$$\mu (s_m^{(n)})=\frac{2^n\lambda _1^{n_1(m)}\lambda _2^{n_2(m)}}{\left[\frac{1}{2}(\lambda _1^1+\lambda _2^1)\right]^n}=\frac{\lambda _1^{n_2(m)}\lambda _2^{n_1(m)}}{(\lambda _1+\lambda _2)^n}.$$
(18)
(Note that this is consistent with $`\mu ([0,1])=_{m=0}^{n1}\mu (s_m^{(n)})=1`$.) Thus the measures of the intervals $`[0,1/2]`$ and $`[1/2,1]`$ are
$$p=\mu (s_0^{(1)})=\frac{\lambda _2}{\lambda _1+\lambda _2}$$
and
$$1p=\mu (s_1^{(1)})=\frac{\lambda _1}{\lambda _1+\lambda _2}.$$
It is important to note that our natural transient measures $`p`$ and $`(1p)`$ for the two-dimensional map are different from the natural measures of the same $`x`$-intervals for the one-dimensional map, $`x_{n+1}=2x_nmod\mathrm{\hspace{0.17em}1}`$, alone. In that case, with probability one, a random choice of $`x_0`$ produces an orbit which spends half its time in $`[0,1/2]`$ and half its time in $`[1/2,1]`$, so that in this case the natural measures of these regions are $`p=(1p)=1/2`$. The addition of the $`y`$-dynamics changes the natural measure of $`x`$-intervals.
From (16) we obtain
$$h_a=\frac{\lambda _2}{\lambda _1+\lambda _2}\mathrm{ln}\lambda _1+\frac{\lambda _1}{\lambda _1+\lambda _2}\mathrm{ln}\lambda _2.$$
(19)
For a general function $`f(Z)`$ with $`d^2f/dZ^2<0`$, averaging over different values of $`Z`$ gives the well-known inequality $`f(Z)f(Z)`$ where $`(\mathrm{})`$ denotes the average of the quantity $`(\mathrm{})`$. Using $`f(Z)=\mathrm{ln}Z`$ with $`Z=\lambda _1`$ with probability $`p=\lambda _2/(\lambda _1+\lambda _2)`$ and $`Z=\lambda _2`$ with probability $`(1p)`$, this inequality and Eqs. (17) and (19) yield the result that
$$h_a\frac{1}{\tau }.$$
(20)
In fact, we will see in Sec. 3.5 that (20) remains true for any choice of the function $`\lambda (x)>1`$ in (14). In (20) the equality sign applies if the vertical stretching is uniform ($`\lambda _1=\lambda _2`$) but, for any nonuniformity in the vertical stretching ($`\lambda _1\lambda _2`$), $`h_a`$ is strictly less than $`1/\tau `$.
### 3.2 Application of the Dimension Formulae
Let $`\lambda _2=r\lambda _1`$, $`r>1`$, and imagine that we fix $`r`$ and vary $`\lambda _1`$. Applying Eqs. (11) and (12) to our example we obtain three cases,
(a) $`h_b>1/\tau >h_a`$ ($`\lambda _1`$ small),
(b) $`1/\tau >h_b>h_a`$ ($`\lambda _1`$ moderate), and
(c) $`1/\tau >h_a>h_b`$ ($`\lambda _1`$ large).
Corresponding to these three cases (11) and (12) yield the following values for $`D_\mathrm{\Lambda }`$, the dimension of the invariant set,
$$D_a=1+\frac{\mathrm{ln}(1+r^1)\mathrm{ln}\lambda _1}{\mathrm{ln}2},\text{for}\lambda _1\lambda _a,$$
(21)
$$D_b=\frac{\mathrm{ln}(1+r^1)+(1+r)^1\mathrm{ln}r}{\mathrm{ln}\lambda _1+(1+r)^1\mathrm{ln}r},\text{for}\lambda _a\lambda _1\lambda _b,$$
(22)
$$D_c=\frac{\mathrm{ln}(1+r^1)+(1+r)^1\mathrm{ln}r}{\mathrm{ln}2},\text{for}\lambda _b\lambda _1,$$
(23)
where $`\mathrm{ln}\lambda _a=\mathrm{ln}(1+r^1)`$ and $`\mathrm{ln}\lambda _b=\mathrm{ln}2(1+r)^1\mathrm{ln}r`$. The solid line in Fig. 2 shows a plot of $`D_\mathrm{\Lambda }`$ versus $`\mathrm{ln}\lambda _1`$ for $`r=3`$. Note that for large $`\lambda _1`$, $`D_\mathrm{\Lambda }=D_c`$ is independent of $`\lambda _1`$.
It is also instructive to consider the case of uniform stretching $`(r=1)`$ for which $`\lambda _1=\lambda _2`$. In that case, $`h_a=1/\tau `$, and there is a rigorous known result for the dimension . For $`\lambda _1=\lambda _2`$, Eqs. (21)–(23) yield
$$D_\mathrm{\Lambda }=\{\begin{array}{ccc}2\frac{\mathrm{ln}\lambda _1}{\mathrm{ln}2}& \text{for}& 1\lambda _12,\hfill \\ 1& \text{for}& \lambda _12.\hfill \end{array}$$
(24)
(For $`r1`$ region (b), where $`D_\mathrm{\Lambda }=D_b`$, shrinks to zero width in $`\lambda _1`$.) For $`r=1`$ the natural transient measure is uniform; from (18) we have $`\mu (s_m^{(n)})=2^n`$ independent of the interval (i.e, independent of $`m`$). In this case there is no difference between the capacity dimension of the invariant set and the information dimension of its measure. Equation (24) agrees with the rigorous known result, thus lending support to the original conjecture.
### 3.3 Numerical Tests
The formulae (21)–(23) were verified by numerical measurements of the information dimension, $`D_1`$, of $`\mathrm{\Lambda }`$ at various values of $`\lambda _1`$ with $`r=\lambda _2/\lambda _1`$ fixed at $`r=3`$ (Fig. 2). Shown for comparison is the box-counting dimension, $`D_0`$. The values of the box-counting dimension are numerically indistinguishable from the values of the information dimension when $`D_0,\text{ }D_1>1`$, or $`\lambda _1<\lambda _a`$, the region corresponding to formula (21). For $`\lambda _1>\lambda _a`$, $`\mathrm{\Lambda }`$ is a smooth curve and so has a box-counting dimension of $`D_0=1`$. No points for $`D_1`$ are shown near $`\lambda _1=\lambda _b`$. It can be argued (Appendix A) that numerical convergence is too slow here to yield accurate measurements of the dimension.
To numerically determine the information dimension of $`\mathrm{\Lambda }`$ (the data shown as open squared in Fig. 2), we place a square $`xy`$ grid with a spacing $`\epsilon `$ between grid points over a region containing $`\mathrm{\Lambda }`$. Using the method dscribed in the next paragraph we compute the natural measure in each grid box and repeat for various $`\epsilon `$. The information dimension is then given by
$$D_1=\underset{\epsilon 0}{lim}\frac{I(\epsilon )}{\mathrm{ln}(1/\epsilon )},$$
where $`I(\epsilon )=_{i=1}^{N(\epsilon )}\mu _i\mathrm{ln}(1/\mu _i)`$ is a sum over the $`N(\epsilon )`$ grid boxes which intersect $`\mathrm{\Lambda }`$ and $`\mu _i`$ is the natural measure in the $`i`$th box. The slope of a plot of $`I(\epsilon )`$ versus $`\mathrm{ln}\epsilon `$ gives $`D_1`$. The box-counting dimension (the data shown as black dots in Fig. 2) is given by
$$D_0=\underset{\epsilon 0}{lim}\frac{\mathrm{ln}N(\epsilon )}{\mathrm{ln}1/\epsilon }$$
and calculated in an analogous way.
To determine which boxes intersect $`\mathrm{\Lambda }`$ and what measure is contained in each of them we take advantage of the fact that $`\mathrm{\Lambda }`$ is a function \[see Appendix A, Eq. (A2)\]. That is, for each value of $`x`$ there is only one corresponding value of $`y`$ in $`\mathrm{\Lambda }`$, which we denote $`y=y_\mathrm{\Lambda }(x)`$. We divide the interval $`0x<1`$ into $`2^n`$ intervals of width $`\delta 2^n`$. We wish to approximate $`y_\mathrm{\Lambda }(x_0)`$ for $`x_0`$ in the center of the $`x`$-interval. To do this we iterate $`x_0`$ forward using (13) $`m`$ times until the condition
$$\frac{\delta }{2}\lambda _1^{m_1}\lambda _2^{m_2}1$$
(25)
is first met, where $`m_1`$ ($`m_2`$) is the number of times the orbit lands in $`0x<1/2`$ ($`1/2x<1`$), and $`m_1+m_2=m`$ (we will see the reason for this condition below). All of the values, $`x_i`$, of the iterates are saved. Starting now from $`x_m`$ and taking $`y_m=0`$ we iterate backward $`m`$ times. For $`\eta `$ small enough, $`\mathrm{\Lambda }`$ is contained in $`1y1`$, so the point ($`x_m`$, $`y_m=0`$) is within a distance $`1`$ in the $`y`$-direction of $`\mathrm{\Lambda }`$. The $`m`$ reverse iterations shrink the segment $`y_myy_\mathrm{\Lambda }(x_m)`$ by a factor $`\lambda _1^{m_1}\lambda _2^{m_2}`$ so that $`\left|y_0y_\mathrm{\Lambda }(x_0)\right|\delta /2`$ by condition (25). Thus, we have found a point $`y_0`$ that approximates $`y_\mathrm{\Lambda }(x_0)`$ to within $`\delta `$. Since we will be using $`ϵ`$ boxes with $`ϵ\delta `$, we may regard $`y_0`$ as being essentially equal to $`y_\mathrm{\Lambda }(x_0)`$. The measure in the $`\delta `$ width interval containing $`x_0`$ is found by iterating $`x_0`$ forward $`n`$ times and using equation (18), $`\mu =\frac{\lambda _1^{n_2}\lambda _2^{n_1}}{(\lambda _1+\lambda _2)^n}`$, where $`n_1`$ ($`n_2`$) is the number of times the orbit lands in $`0x<1/2`$ ($`1/2x<1`$), and $`n_1+n_2=n`$. We associate this measure with the point $`(x_0,y_0)`$ (with $`y_0`$ found by the above procedure).
Note that for fractal $`y_\mathrm{\Lambda }(x)`$ the $`y`$ interval occupied by the curve $`y=y_\mathrm{\Lambda }(x)`$ in an $`x`$ interval of width $`\delta 1`$ is of order $`\delta ^{D_01}`$ which is large compared to $`\delta `$. We now cover the region with new grids having successively larger spacing, $`\epsilon _i=2^i\delta =2^{in}`$, and calculate $`I(\epsilon _i)`$ and $`N(\epsilon _i)`$ based on the data taken from the first $`\delta `$-grid. For $`i`$ large enough, such that the $`y`$ extent of the curve $`y_\mathrm{\Lambda }(x)`$ in a typical $`\delta `$ width interval is less than $`ϵ_i`$ (i.e., $`ϵ_i\stackrel{}{>}\delta ^{2D_0}`$ or $`i\stackrel{}{>}(D_01)n`$) we observe linear scaling of $`\mathrm{log}I(ϵ_i)`$ and $`\mathrm{log}N(ϵ_i)`$ with $`\mathrm{log}ϵ_i`$, and we use the slope of such plots to determine $`D_1`$ and $`D_0`$. The dimensions $`D_1`$ and $`D_0`$ are then determined as described above.
### 3.4 Atypical Case
The conjecture of Ref. is that the dimension formulae of Sec. 2 apply for “typical” systems. To see the need for this restriction consider Eqs. (14) for the case where $`\eta =0`$. It is easily shown that the dimension formulae can be violated in this case. The claim, however, is that $`\eta =0`$ is special, or “atypical”, in that, as soon as we give $`\eta `$ any nonzero value, the validity of the dimension formulae is restored. In this connection it is important to note that as long as $`\eta 0`$, the dimension of the invariant set is independent of the value of $`\eta `$. This follows since if $`\eta 0`$ we can always rescale the value of $`\eta `$ to one by the change of variables $`\stackrel{~}{y}=y/\eta `$. To see the violation of the dimension formulae for $`\eta =0`$, we note that in this case, by virtue of (14), the line $`y=0`$ is invariant. Thus the measure is distributed on a one-dimensional subspace, the $`x`$-axis. (This is very different from the picture in Fig. 1, where the invariant set, the boundary between black and white, appears to be fractal.) Using the definition of the information dimension and dividing the $`x`$-axis into intervals of width $`2^n`$, the information dimension of the natural measure is
$$D_\mathrm{\Lambda }=\underset{n\mathrm{}}{lim}\frac{_{m=0}^{n1}\mu (s_m^{(n)})\mathrm{ln}[1/\mu (s_m^{(n)})]}{\mathrm{ln}(2^n)}.$$
(26)
The quantity whose limit is taken in (26) is in fact independent of $`n`$. Thus, taking $`n=1`$ we obtain for $`D_\mathrm{\Lambda }`$ the result that, for $`\eta =0`$,
$$D_\mathrm{\Lambda }=D_c,$$
*for all* $`\lambda _1`$ and $`\lambda _2>1`$, where $`D_c`$ is given by (23). Thus, for $`h_a<h_b`$, $`D_\mathrm{\Lambda }`$ is greater when $`\eta 0`$ than when $`\eta =0`$, and, thus, the conjectured stable manifold dimension formula of Sec. 2 is violated. For $`h_a>h_b`$, $`D_\mathrm{\Lambda }`$ is the same in both cases.
In Sec. 4 we consider chaotic scattering in a three-dimensional billiard example for which the character of the atypical case is more interesting than in the above example. In particular, we find for our billiard that in both the typical and the atypical cases the stable manifold can have noninteger capacity dimension.
### 3.5 General Considerations
The previous considerations readily generalize to the case of an arbitrary smooth function $`\lambda (x)>1`$ and a general chaotic map, $`x_{n+1}=M(x_n)`$, which replaces (3.1). Consider the finite time vertical Lyapunov exponent,
$$\stackrel{~}{h}(x,n)=\frac{1}{n}\underset{m=1}{\overset{n}{}}\mathrm{ln}\lambda (M^{m1}(x))$$
computed for the initial condition $`x`$. Choosing $`x`$ randomly with uniform probability distribution in the relevant basin for chaotic motion \[e.g., $`x`$ in $`[0,1]`$ for Eq. (13)\], $`\stackrel{~}{h}(x,n)`$ can be regarded as a random variable. Let $`\stackrel{~}{P}(h,n)`$ denote its probability distribution function. For large $`n`$, we invoke large deviation theory to write $`\stackrel{~}{P}(h,n)`$ as
$$\mathrm{ln}\stackrel{~}{P}(h,n)=nG(h)+o(n),$$
or, more informally,
$$\stackrel{~}{P}(h,n)e^{nG(h)},$$
(27)
where the specific form of $`G(h)`$ depends on $`M(x)`$ and the specific $`\lambda (x)`$, and $`G(h)`$ is convex, $`d^2G(h)/dh^20`$. For the normalization, $`\stackrel{~}{P}(h,n)𝑑h=1`$, to hold for $`n\mathrm{}`$, we have that
$$\underset{h}{\mathrm{min}}G(h)=0;$$
see Fig. 3, where $`\overline{h}`$ denotes the value of $`h`$ for which the above minimum is attained. As $`n\mathrm{}`$ we see that $`\stackrel{~}{P}`$ approaches a delta function, $`\delta (h\overline{h})`$. Thus, $`\overline{h}`$ is the usual infinite time Lyapunov exponent for almost all initial conditions with respect to Lebesgue measure in $`0x1`$.
As described above $`\stackrel{~}{P}(h,n)`$ is the probability distribution of $`h(x,n)`$ for $`x`$ chosen randomly with respect to a uniform distribution in $`[0,1]`$. We now ask what the probability distribution of $`h(x,n)`$ is for $`x`$ chosen randomly with respect to the natural transient measure for our expanding map, $`x_{n+1}=M(x_n)`$ and (14). To answer this question we proceed as before and consider an initial vertical line segment $`|y|K`$ starting at $`x`$ \[with $`K>(\eta /2\pi )(\lambda _{\mathrm{min}}1)^1`$, $`\lambda _{\mathrm{min}}=\mathrm{min}_x\lambda (x)>1`$\]. After $`n`$ iterations, this line segment lengthens by the factor $`\mathrm{exp}[n\stackrel{~}{h}(x,n)]`$. Thus, the fraction of the line still remaining in the strip $`|y|<K`$ is $`\mathrm{exp}[n\stackrel{~}{h}(x,n)]`$. Hence, the fraction of points sprinkled uniformly in the strip that still remains after $`n`$ iterates is
$$e^{n/\tau }e^{nG(h)nh}𝑑h,$$
(28)
and the probability distribution of finite time vertical Lyapunov exponents for $`x`$ chosen randomly with respect to the natural transient measure is
$$P(h,n)\frac{e^{nG(h)nh}}{e^{nG(h)nh}𝑑h}.$$
(29)
Evaluating (28) for large $`n`$ we have $`e^{n[G(h)+h]}𝑑he^{n[G(h_{})+h_{}]}`$, where $`\mathrm{min}[G(h)+h]=G(h_{})+h_{}`$ and $`h_{}`$ is the solution of $`dG(h_{})/dh_{}=1`$. Thus ,
$$1/\tau =G(h_{})+h_{}.$$
(30)
See the construction in Fig. 3, in which the dotted line of slope $`1`$ is tangent to the graph of $`G(h)`$ at the point $`h=h_{}`$. The infinite time vertical Lyapunov exponent for the transient natural measure is
$$h_a=hP(h,n)𝑑h.$$
(31)
Using (29) and again letting $`n`$ be large (31) yields $`h_a=h_{}`$. Referring to Fig. 3 we see that
$$h_a1/\tau .$$
That is, Eq. (20) is valid for general $`M(x)`$ and $`\lambda (x)`$ and not just for $`M(x)`$ and $`\lambda (x)`$ given by (13) and (15).
## 4 A Three-Dimensional Billiard Chaotic Scatterer
We consider a three degree-of-freedom billiard. The billiard (Fig. 4) is formed by a hard ellipsoid of revolution, placed in a hard, infinitely long tube with cross-section as shown in Fig. 4(b). The center of the ellipsoid is placed at the center of the tube. We consider two cases: (a) the major axis of the ellipsoid coincides with the $`z`$-axis \[see Fig 4(a)\], (b) the major axis of the ellipsoid lies in the $`y`$-$`z`$ plane and makes an angle $`\xi `$ with the $`z`$-axis. The ratio of the minor radius of the ellipsoid ($`r_{}`$) to the width of a side of the tube \[dashed line in Fig. 4(b)\] is $`1/4`$. This leaves the major radius, $`r_{}`$, and the tilt angle, $`\xi `$, as parameters. A point particle injected into the system experiences specular reflection from the ellipsoid and the walls (i.e., the angle of reflection is equal to the angle of incidence, where both are taken with respect to the normal to the surface off of which the particle bounces). When the orbit has passed the top (bottom) of the ellipsoid, with positive (negative) $`z`$-velocity, we say that it has exited upward (downward). We fix the conserved energy so that $`\left|\stackrel{}{v}\right|=1`$.
### 4.1 Pictures of the Stable Manifold
We claim that the case of the untilted ellipsoid \[$`\xi =0`$ and case (a) above\] is atypical in the sense that it violates the formulae of Sec. 2. However, as soon as $`\xi 0`$, we claim that the formulae of Sec. 2 apply. We begin by discussing the atypical (untilted) case.
By the symmetry of the geometry of the billiard with the ellipsoid axis along $`z`$ \[Fig. 4(a)\], the chaotic saddle, $`\mathrm{\Lambda }`$, of this system is the collection of initial conditions satisfying $`z=v_z=0`$. Started with these initial conditions, an orbit will have $`z=v_z=0`$ for all forward and reverse time. The surface normals of the ellipsoid and walls at $`z=0`$ lie in the $`z=0`$ plane, and thus the particle cannot acquire a non-zero $`v_z`$. The $`z=0`$ slice through the three-dimensional billiard is a two-dimensional billiard with concave walls \[Fig. 4(b)\]. It is known that a typical orbit in this billiard will fill the phase space ergodically. Near $`z=0`$, we can picture a typical point on the stable manifold (denoted $`SM`$) as having, for example, $`z`$ slightly less than zero and $`v_z`$ slightly greater than zero. The particle will hit the ellipsoid below its equator and, thus, $`v_z`$ will be decreased with each bounce, yet remain positive. With successive bounces, the orbit on $`SM`$ slowly approaches $`z=v_z=0`$, the chaotic saddle.
To visualize $`SM`$, we note that it forms the boundary between initial conditions which escape upward and those which escape downward. We say that points which, when iterated, eventually escape upward (downward) are in the *basin* of upward (downward) escape. Points which are on the boundary between the two basins never escape at all, i.e. they are in $`SM`$. We initiate a (two-dimensional) grid of orbits (500x500) on the plane,$`3<x<3`$, $`y=5.1`$, $`2.5<z<0`$, $`v_x=0`$, $`v_z=.1`$, and $`v_y`$ is given by the condition $`\left|\stackrel{}{v}\right|=1`$. We iterate each of these initial conditions forward until it escapes. Then we plot a white (black) point for each orbit in the upward (downward) basin. The result is shown in Fig. 5(a). The boundary between the white and black regions is then the intersection of $`SM`$ lying in the phase space of the five-dimensional billiard ($`x,y,z,v_x,v_y,v_z`$ constrained by $`\left|\stackrel{}{v}\right|=1`$) with the specified two-dimensional $`x,z`$-plane. $`SM`$ appears to take the form of a nowhere-differentiable curve. This is true in various 2D slices, none of which are chosen specially, which suggests that $`SM`$ has this form in a typical slice. A similar procedure can be followed for the case of the tilted ellipsoid and the resulting picture is shown in Fig. 5(b) for the case of a tilt angle of $`\xi =2\pi /100`$ radians.
### 4.2 Lyapunov Exponents, Decay Times, and Approximate Formulae for the Stable Manifold
Again, we begin with the untilted case. To construct a map from this system we record the cylindrical coordinates ($`z`$, $`\varphi `$) and their corresponding $`z`$ and $`\varphi `$ velocity components, which we denote ($`v`$, $`\omega `$), each time the particle hits the ellipsoid. The coordinate $`r`$ is constrained, for a given $`z`$, by the shape of the ellipsoid surface, and $`v_r`$ is given by the energy conservation condition $`\left|\stackrel{}{v}\right|=1`$. The four components $`(`$$`z`$, $`v`$, $`\varphi `$, $`\omega `$$`)_n`$ give the state of the system at discrete time $`n`$, where $`n`$ labels the number of bounces from the ellipsoid. Let $`\stackrel{}{z}\left[\begin{array}{c}z\\ v\end{array}\right]`$ and $`\stackrel{}{\varphi }\left[\begin{array}{c}\varphi \\ \omega \end{array}\right]`$. For future reference, we express the map using the following notation,
$$\begin{array}{ccc}\hfill \stackrel{}{z}_{n+1}& =& M_z(\stackrel{}{z}_n,\stackrel{}{\varphi }_n),\hfill \\ \hfill \stackrel{}{\varphi }_{n+1}& =& M_\varphi (\stackrel{}{z}_n,\stackrel{}{\varphi }_n).\hfill \end{array}$$
(32)
In what follows, when we refer to an orbit, saddle, invariant set, stable manifold, etc., we are referring to these quantities for the discrete time map (rather than the original continuous time system).
In the case of the untilted ellipsoid, linearizing about an orbit on $`\mathrm{\Lambda }`$, (i.e., $`\stackrel{}{z}_n=0`$, $`\stackrel{}{\varphi }_n`$), we obtain, for the evolution of differential orbit perturbations $`\delta \stackrel{}{z}`$ and $`\delta \stackrel{}{\varphi }`$,
$$\begin{array}{ccc}\hfill \delta \stackrel{}{z}_{n+1}& =& DM_z(0,\stackrel{}{\varphi }_n)\delta \stackrel{}{z}_n,\hfill \\ \hfill \delta \stackrel{}{\varphi }_{n+1}& =& DM_\varphi (0,\stackrel{}{\varphi }_n)\delta \stackrel{}{\varphi }_n,\hfill \end{array}$$
where $`\stackrel{}{\varphi }_{n+1}=M_\varphi (0,\stackrel{}{\varphi }_n)`$ is the map for the two-dimensional billiard of Fig. 4(b), $`DM_z(0,\stackrel{}{\varphi })`$ is the tangent map for differential orbit perturbations in $`\stackrel{}{z}`$ evaluated at $`\stackrel{}{z}=0`$, and $`DM_\varphi (0,\stackrel{}{\varphi })`$ is the tangent map for differential perturbations lying in $`\mathrm{\Lambda }`$. Let $`\pm h_z`$ and $`\pm h_\varphi `$ denote the Lyapunov exponents with respect to the *natural transient measure* (Sec. 2) for perturbations in $`\stackrel{}{z}`$ and in $`\stackrel{}{\varphi }`$, respectively (these exponents occur in positive-negative pairs due to the Hamiltonian nature of the problem).
In principle, one could numerically evaluate $`h_z`$ and $`h_\varphi `$ by sprinkling a large number, $`N`$, of initial conditions in the vicinity of $`\mathrm{\Lambda }`$, iterating $`n1`$ times, evaluating $`h_z`$ and $`h_\varphi `$ over those orbits still near $`\mathrm{\Lambda }`$, and averaging $`h_z`$ and $`h_\varphi `$ over those orbits. We could also find $`\tau `$ by this procedure; it is the exponential rate of decay of the orbits from the vicinity of $`\mathrm{\Lambda }`$. For cases where the escape time $`\tau `$ is not long, this procedure, however, becomes problematic. For finite $`N`$ the number of retained orbits can be small or zero if $`n`$ is too large. Thus, we adopt an alternate procedure which we found to be less numerically demanding.
In particular we make user of the ideas presented in Sec. 3.5. We define the *uniform measure* as the measure generated by uniformly sprinkling many initial conditions in $`\mathrm{\Lambda }`$ (the hyperplane $`z=v_z=0`$). An average over these orbits of the tangent space stretching exponents would yield uniform measure Lyapunov exponents. We denote the distribution of finite-time Lyapunov exponents with respect to this measure by $`P(\stackrel{~}{h}_\varphi ,\stackrel{~}{h}_z,n)e^{nG(\stackrel{~}{h}_\varphi ,\stackrel{~}{h}_z)}`$ \[as in Eq. (27)\], where the tilde indicates finite time exponents for initial conditions distributed according to the uniform measure.
To compute the decay time and the Lyapunov exponents with respect to the natural transient measure we note that orbits near a point $`\stackrel{}{\varphi }`$ in $`\mathrm{\Lambda }`$ iterate away from $`\mathrm{\Lambda }`$ as $`\mathrm{exp}[n\stackrel{~}{h}_z(\stackrel{}{\varphi },n)]`$. Thus, the fraction of a large number of initial conditions sprinkled near $`\mathrm{\Lambda }`$ which remain near $`\mathrm{\Lambda }`$ after $`n`$ iterates is
$$P(\stackrel{~}{h}_\varphi ,\stackrel{~}{h}_z,n)e^{n\stackrel{~}{h}_z}𝑑\stackrel{~}{h}_z,$$
and $`h_\varphi `$ (the infinite time Lypunov exponent with respect to the natural transient measure) is
$$h_\varphi =\underset{n\mathrm{}}{lim}\frac{\stackrel{~}{h}_\varphi P(\stackrel{~}{h}_\varphi ,\stackrel{~}{h}_z,n)e^{n\stackrel{~}{h}_z}𝑑h}{P(\stackrel{~}{h}_\varphi ,\stackrel{~}{h}_z,n)e^{n\stackrel{~}{h}_z}𝑑h}.$$
This expression can be approximated numerically by choosing $`N`$ initial conditions, $`\stackrel{}{\varphi }_i`$, uniformly in $`\mathrm{\Lambda }`$ and calculating
$$n\stackrel{~}{h}_\varphi _n\frac{n_{i=1}^N\stackrel{~}{h}_\varphi (\stackrel{}{\varphi }_i,n)e^{n\stackrel{~}{h}_z(\stackrel{}{\varphi }_i,n)}}{N_{i=1}^Ne^{n\stackrel{~}{h}_z(\stackrel{}{\varphi }_i,n)}}.$$
We calculate the finite-time Lyapunov exponents $`\stackrel{~}{h}_\varphi (\stackrel{}{\varphi }_i,n)`$ and $`\stackrel{~}{h}_z(\stackrel{}{\varphi }_i,n)`$ using the QR decomposition method . Since we chose the $`\stackrel{}{\varphi }_i`$ uniformly in $`\mathrm{\Lambda }`$ these exponents are distributed according to $`P(n,\stackrel{~}{h}_\varphi ,\stackrel{~}{h}_z)`$. $`N`$ is taken to be large enough that there are at least $`100`$ terms contributing to $`90\%`$ of each sum. The range of $`n`$ is from $`10`$ to about $`40`$. We find that $`n\stackrel{~}{h}_\varphi _n`$ versus $`n`$ is well-fitted by a straight line and we take its slope as our estimate of $`h_\varphi `$.
To find $`\tau `$, first note that we numerically find that $`h_\varphi `$ does not vary much from $`\overline{h}_\varphi `$ ($`.91h_\varphi /\overline{h}_\varphi 1`$) as we change the system parameter (the height of the ellipsoid), where the overbar denotes an infinite-time Lypunov exponent with respect to the uniform measure on $`\mathrm{\Lambda }`$. Thus, we make the approximation $`P(\stackrel{~}{h}_\varphi ,\stackrel{~}{h}_z,n)P(\stackrel{~}{h}_z,n)`$, $`G(\stackrel{~}{h}_\varphi ,\stackrel{~}{h}_z)G(\stackrel{~}{h}_z)`$ and plot $`G(\stackrel{~}{h}_z)`$ versus $`\stackrel{~}{h}_z`$. The value of $`1/\tau `$ is given by Eq. (30), and $`h_z`$ is given by $`dG(n,h_z)/d\stackrel{~}{h}_z=1`$. (Note that $`h_z1/\tau \overline{h}_z`$.) A third order polynomial is fit to the data for $`G(\stackrel{~}{h}_z)`$ and used to find $`\tau `$ and $`h_z`$. This calculation is performed at a value of $`n`$ which allows a significant number of points in $`G(\stackrel{~}{h}_z)`$ versus $`\stackrel{~}{h}_z`$ to be collected near in the range $`h_z<\stackrel{~}{h}_z<\overline{h}_z`$.
For the case of the tilted ellipsoid, we will consider a very small tilt angle, $`\xi =2\pi /100`$. With this small tilt angle, the Lyapunov exponents and decay times for the tilted and the untilted cases are approximately the same. Thus, for the tilted case, we will use the same Lyapunov exponent values, $`\pm h_z`$ and $`\pm h_\varphi `$ and decay time $`\tau `$, that we numerically calculated for the untilted case. The dimension formula of Sec. 2 for $`SM`$ is
$$\begin{array}{cc}D_S=4(h_\varphi \tau )^1& \text{for}h_\varphi \tau 1,\\ D_S<3& \text{for}h_\varphi \tau <1.\end{array}$$
(33)
(As for cases (b) and (c) of Sec. 3.2, we could write explicit expressions for $`D_S`$ for $`h_\varphi \tau <1`$, but, in what follows we only need that $`D_S<3`$ in this case.)
If the tilt angle is made to be zero ($`\xi =0`$), we find that $`D_S`$ is not given by (33), but by the following formula (which we derive in Sec. 7.2),
$$\begin{array}{cc}D_S=4\frac{h_z+1/\tau }{h_\varphi }& \text{for}h_\varphi h_z+1/\tau ,\\ D_S<3& \text{for}h_\varphi <h_z+1/\tau .\end{array}$$
(34)
Since $`D_S`$ is given by (34) only if the tilt angle $`\xi `$ is precisely zero, we say that the untilted ellipsoid scattering system is atypical. As conjectured in , $`D_S`$ from (33) is greater than or equal to $`D_S`$ from (34). Note that $`D_S`$ from the first line of (33) is larger that $`D_S`$ from the first line of (34) by the factor $`h_z/h_\varphi `$. Although the transition of $`D_S`$ from $`\xi =0`$ to $`\xi 0`$ is strictly discontinuous, there is also a continuous aspect: In numerically calculating the dimension of a measure one typically plots $`\mathrm{ln}I(ϵ)`$ versus $`\mathrm{ln}(1/ϵ)`$, where $`I(ϵ)=\mu _i\mathrm{ln}[1/\mu _i]`$ and $`\mu _i`$ is the measure of the $`i^{\text{th}}`$ cube in an $`ϵ`$ grid. One then estimates the dimension as the slope of a line fitted to small $`ϵ`$ values in such a plot. In the case of very small tilt, such a plot is expected to yield a slope given by (34) for $`ϵ\stackrel{}{>}ϵ_{}`$ and subsequently, for $`ϵ\stackrel{}{<}ϵ_{}`$, to yield a slope given by (33). Here $`ϵ_{}`$ is a small tilt-dependent cross-over value, where $`ϵ_{}0`$ as $`\xi 0`$. In such a case, the dimension, which is defined by the $`ϵ0`$ limit, is given by (33).
## 5 Numerical Computations for the Three-Dimensional Billiard Scatterer
To verify that the untilted system is atypical we numerically calculated the box-counting dimension of $`SM`$ for various values of the parameter $`r_{}`$, then introduced a small tilt perturbation in the form of a $`2\pi /100`$ radian tilt of the ellipsoid about the $`x`$-axis and repeated the dimension calculations. The results are shown in Fig. 6. They confirm Eqs. (33) and (34). (We do not plot data points near the kink in the theoretical curve because of the numerical difficulty in obtaining reliable results in that parameter range. This is a result of slow converge near the kink. It is also found in the example of Sec. 3.3 and discussed in Appendix A.)
We compare measured values of the box-counting dimension to the predicted values of the information dimension in Fig. 6. The box-counting dimension gives an upper bound on the information dimension, but often the values of the two dimensions are very close. In particular, we can compare the result for the tilted ellipsoid system to the 2D map of Sec. 3. In the regime $`h_\varphi 1/\tau `$, our system is similar to case (a) studied in Section 3.2 (note that, as discussed in Sec. 4, $`1/\tau `$ is never smaller than $`h_z`$). Changing $`r_{}`$ while leaving $`r_{}`$ fixed changes $`\tau `$ while $`h_\varphi `$ changes only slightly. This is similar to varying $`\lambda _1`$ of the 2D map while keeping $`\lambda _2`$ fixed, i.e. varying $`r`$. Figure 2, shows (for a particular value of $`r`$) that the measured values of $`D_0`$ and $`D_1`$ are numerically indistinguishable in the region corresponding to case (a), $`\lambda <\lambda _a`$.
For $`h_\varphi \tau <1`$ ($`h_\varphi <h_z+1/\tau `$) the information dimension of $`SM`$ is predicted to be less than three for the tilted (untilted) ellipsoid system. This is analogous to cases (b) and (c) of Sec. 3.2; see Fig. 2. Since $`SM`$ divides the four-dimensional phase space, its box-counting dimension cannot be less than three, and, analogous to the result of Sec. 3 and Fig. 2, we expect that the box-counting dimension is three for $`h_\varphi \tau <1`$ ($`h_\varphi <h_z+1/\tau `$). This expectation is borne out by the numerical results, Fig. 6.
The box-counting dimension of $`SM`$ was computed using the uncertainty dimension method . This method gives the box-counting dimension of the basin boundary which, as discussed above, coincides with $`SM`$.
The uncertainty dimension method was carried out as follows:
1. Choose a point, $`\stackrel{}{x}`$, at random in a region of a 2D plane intersecting the basin boundary (e.g., Fig. 6) and determine by iteration in which basin it lies.
2. Determine in which basins the perturbed initial points $`\stackrel{}{x}\pm \stackrel{}{\delta }`$ lie ($`\stackrel{}{\delta }`$ is some small vector).
3. If the three points examined in (1) and (2) do not all lie in the same basin, then $`\stackrel{}{x}`$ is called “uncertain”.
4. Repeat 1 to 3 for many points $`\stackrel{}{x}`$ randomly chosen in the 2D plane, and obtain the fraction of these that are uncertain.
5. The fraction of points which is uncertain for a given $`\stackrel{}{\delta }`$, denoted $`f`$, scales like $`f|\stackrel{}{\delta }|^{2d_0},`$ where $`d_0`$ is the box-counting dimension of the intersection of $`SM`$ with the 2D plane. The box-counting dimension of $`SM`$ in the full 4D state space of the map is $`D_0=2+d_0`$. (The dimension of a generic intersection of a 2D plane with a set having dimension $`D_0`$ in a 4D space is $`d_0=2+D_04`$, which gives $`D_0=2+d_0`$.) Thus, plot $`\mathrm{ln}f`$ versus $`\mathrm{ln}\left|\stackrel{}{\delta }\right|`$, fit a straight line to the plot, and estimate $`D_0`$ as $`4`$ minus the slope of this line. See Fig. 8 for an example of such a plot.
## 6 Structure of the Stable Manifold
### 6.1 Untilted Ellipsoid
Due to the symmetry to the untilted ellipsoid billiard, the chaotic saddle of the untilted ellipsoid system has a special geometry (i.e., it lies in $`z=v=0`$), which, as we show, accounts for the dimension being lower than the predicted value for a typical system. Similarly, the symmetry induces a special geometry on the stable manifold ($`SM`$). Figure 7 shows a 3D slice of $`SM`$ in the atypical system. The slice is at a fixed value of $`\omega `$. The axes are $`z`$, $`v`$, and $`\varphi `$, but one could have chosen an arbitrary line through ($`\varphi `$,$`\omega `$) as the third axes and seen a plot which was qualitatively the same. The stable manifold is organized into rays emanating from the $`\varphi `$-axis with oscillations along the $`\varphi `$-direction. The magnitude of the oscillations decreases to zero as the $`\varphi `$-axis is approached.
To understand the structure of $`SM`$ in more detail, assume that $`\left|\stackrel{}{z}\right|`$ is small. Then, since $`\left|\stackrel{}{z}\right|=0`$ is invariant, we can approximate the dynamics by expanding to first order in $`\stackrel{}{z}`$,
$$\stackrel{}{z}_{n+1}DM_z(0,\stackrel{}{\varphi }_n)\stackrel{}{z}_n,$$
(35)
$$\stackrel{}{\varphi }_{n+1}M_\varphi (0,\stackrel{}{\varphi }_n).$$
(36)
Say $`(\stackrel{}{z}_{SM},\stackrel{}{\varphi }_{SM})`$ is a point on $`SM`$. As this point is iterated we have that $`\left|\stackrel{}{z}\right|0`$ with increasing $`n`$. However, since (35) is linear in $`\stackrel{}{z}_n`$, for any constant $`\alpha `$, and the initial condition, $`(\alpha \stackrel{}{z}_{SM},\stackrel{}{\varphi }_{SM})`$, the subsequent orbit must also have $`\left|\stackrel{}{z}\right|0`$. Consequently, if $`(\stackrel{}{z}_{SM},\stackrel{}{\varphi }_{SM})`$ lies in $`SM`$, so does $`(\alpha \stackrel{}{z}_{SM},\stackrel{}{\varphi }_{SM})`$. Thus, for the system (35), (36), the stable manifold at any point $`\stackrel{}{\varphi }`$ lies on a straight line through the origin of the two-dimensional $`\stackrel{}{z}`$-space. Put another way, in the approximation (35), (36), the stable manifold can be specified by an equation giving the *angle* of $`\stackrel{}{z}`$ as a function of $`\stackrel{}{\varphi }`$. Thus, decomposing into polar coordinates $`(\rho ,\chi )`$, the stable manifold for $`\left|\stackrel{}{z}\right|0`$ approaches the form
$$\chi =\chi (\stackrel{}{\varphi }).$$
(37)
(This explains the structure seen in Fig. 7.) For $`\left|\stackrel{}{z}\right|`$ finite the linearity of (35) is not exact, and we expect that the behavior of the stable manifold at constant $`\stackrel{}{\varphi }`$ is not a straight line through the origin of the $`\stackrel{}{z}`$ plane. Rather, as $`\left|\stackrel{}{z}\right|`$ becomes larger we expect (and numerically observe) the straight line for small $`\left|\stackrel{}{z}\right|`$ to appear as a smooth curve through $`\stackrel{}{z}=0`$.
Since for fixed $`\stackrel{}{\varphi }`$ $`SM`$ varies smoothly with increasing $`\rho =\left|\stackrel{}{z}\right|`$, the dimension of $`SM`$ is not affected by the approximation (35) and (36). That is, to find the dimension of $`SM`$, we can attempt to find it in the region of small $`\left|\stackrel{}{z}\right|`$ where (35) and (36) are valid, and that determination will apply to the whole of $`SM`$.
The task of analytically determining $`D_S`$ is too hard for us to accomplish in a rigorous way for the system, (35), (36), applying to our billiard \[a heuristic analysis yielding (34) is given in Sec. 7.2\]. Thus, to make progress, we adopt a model system with the same structure as (35), (36). In particular, we wish to replace the two-dimensional billiard map (36) by a simpler map, $`M_\varphi \overline{M}_\varphi `$, that, like the original two-dimensional billiard map, is chaotic and describes a Hamiltonian system. For this purpose we choose the cat map,
$$\stackrel{}{\varphi }_{n+1}=\overline{M}_\varphi (\stackrel{}{\varphi }_n)C\stackrel{}{\varphi }_n\text{modulo}\mathrm{\hspace{0.17em}1},$$
(38)
where $`C`$ is the cat map matrix,
$$C=\left[\begin{array}{cc}2& 1\\ 1& 1\end{array}\right].$$
Similarly, we replace $`DM_z(0,\stackrel{}{\varphi })`$ in (35) by a simple symplectic map depending on $`\stackrel{}{\varphi }`$,
$$\stackrel{}{z}_{n+1}=\overline{M}_z(\stackrel{}{z}_n)=\left[\begin{array}{cc}\lambda & f(\stackrel{}{\varphi }_n)\\ 0& \lambda ^1\end{array}\right]\stackrel{}{z}_n,$$
(39)
where $`\lambda >1`$, and $`f(\stackrel{}{\varphi })`$ is a smooth periodic function with period one in $`\varphi `$ and $`\pi `$. \[Note that vertical (i.e., parallel to $`z`$) line segments are uniformly expanded by the factor $`\lambda `$, and thus, by the same argument as in Sec. 3.1, we have $`1/\tau =\mathrm{ln}\lambda `$ and $`h_z=1/\tau `$.\]
Since only the angle of $`\stackrel{}{z}`$ is needed to specifiy the stable manifold, we introduce the variable $`\nu =z/v=\mathrm{tan}\chi `$. We can then derive a map for $`\nu `$: From (39) we have $`\nu _{n+1}v_{n+1}=\lambda \nu _nv_n+v_nf(\stackrel{}{\varphi }_n)`$ and $`v_{n+1}=\lambda ^1v_n`$. Dividing the first equation by the second, $`v_{n+1}`$ and $`v_n`$ cancel and we obtain
$$\nu _{n+1}=\lambda ^2\nu _n+\lambda f(\stackrel{}{\varphi }).$$
(40)
We now consider the dynamical system consisting of (38) and (40). Note that the system (38), (40) is a three-dimensional map, unlike the system (38), (39), which is a four -dimensional map.
For (38), (40), the stable manifold is given by
$$\nu =\nu _S(\stackrel{}{\varphi })=\lambda ^1\underset{i=0}{\overset{\mathrm{}}{}}\lambda ^{2i}f(C^i\stackrel{}{\varphi }_0),$$
(41)
where $`C`$ is the cat map matrix.
To verify that this is $`SM`$ we note that points above $`SM`$, $`\nu >\nu _S(\stackrel{}{\varphi })`$ \[below $`SM`$, $`\nu <\nu _S(\stackrel{}{\varphi })`$\] are repelled toward $`\nu \mathrm{}`$ ($`\nu \mathrm{}`$). Thus, on backwards iteration, points go toward $`SM`$. We take advantage of this behavior to determine $`SM`$. Imagine that we iterate $`\stackrel{}{\varphi }`$ forward $`n`$ iterates to $`\stackrel{}{\varphi }_n=C^n\stackrel{}{\varphi }\text{modulo}\mathrm{\hspace{0.17em}1}`$, then choose a value of $`\nu _n`$, and iterate it backwards using (40). We find that the initial value of $`\nu `$ at time zero giving the chosen value $`\nu _n`$ at time $`n`$ is
$$\nu _0=(\nu _n/\lambda ^{2n})\lambda ^1\underset{i=0}{\overset{n1}{}}\lambda ^{2i}f(C^i\stackrel{}{\varphi }).$$
Keeping $`\nu _n`$ fixed and letting $`n+\mathrm{}`$, the value of $`\nu _0`$ approaches $`\nu _S(\stackrel{}{\varphi })`$, given by (41).
Results proven in show that the box-counting dimension (the capacity) of the graph of the function $`\nu =\nu _S(\stackrel{}{\varphi })`$ given by (41) is
$$\widehat{D}_S=\{\begin{array}{cc}32\frac{\mathrm{ln}\lambda }{\mathrm{ln}B},& \text{for }\lambda B,\\ 2,& \text{for }\lambda >B,\end{array}$$
(42)
where $`B=\frac{3+\sqrt{5}}{2}>1`$ is the larger eigenvalue of the matirx $`C`$. The formula for $`\lambda <B`$ holds for almost all (with respect to Lebesgue measure) values of $`\lambda `$. Since $`\lambda >1`$, the sum in (41) converges absolutely, implying that $`\nu _S(\stackrel{}{\varphi })`$ is a continuous function of $`\stackrel{}{\varphi }`$. Thus, when the first result in Eq. (42) applies (i.e., the surface is fractal with $`\widehat{D}_S>2`$), the stable manifold is a continuous nondifferentiable surface.
For example, evaluating (41) on the surface $`\stackrel{}{\varphi }=s\widehat{u}_+`$ where $`\widehat{u}_+`$ is the unit vector in the eigendirection of $`C`$ corresponding to the expanding eigenvalue $`B`$, we have that $`\nu `$ versus $`s`$ is of the form
$$\nu =\underset{i=0}{\overset{\mathrm{}}{}}\lambda ^{2i}g(B^is).$$
Thus, the graph of $`\nu `$ versus $`s`$ for $`B>\lambda `$ has the form of Weierstra$`\beta `$’s famous example of a continuous, nowhere-differentiable curve.
To obtain (42) in another way, we again consider the map (38), (40). We claim that (38), (40) can be regarded as a *typical* system, and that the formulae of Sec. 2 should apply to it. That is, in contrast to the existence of a symmetry for (39) \[namely, $`\stackrel{}{z}\stackrel{}{z}`$ leaves (39) invariant\], (40) has no special symmetry. The Lyapunov exponent corresponding to (40) is $`h_\nu =2\mathrm{ln}\lambda `$. Note that for the system (38), (40) \[and also for the system (38), (39)\] there are *no* fluctuations in the finite time Lyapunov exponents and thus the decay time for the system (38), (40) is given by $`\tau _\nu ^1=h_\nu `$. Noting that the Lyapunov exponents for the cat map are $`\pm \mathrm{ln}B`$ and applying (3) to the three -dimensional map (38), (40), we immediately recover Eq. (42). As discussed in Appendix B, this point of view can also be exploited for the original ellipsoid system \[rather than just for the model system (38) and (39)\].
Returning now to the full four-dimensional system, (38), (39), and noting that $`SM`$ is smooth along the direction that we eliminated when we went from (38), (39) to (38), (40), we have that the dimension $`D_S`$ of the stable manifold of the invariant set ($`\stackrel{}{z}=0`$) for (38), (39) is $`D_S=\widehat{D}_S+1`$. The Lyapunov exponents for $`\stackrel{}{z}`$ motion in the four-coordinate system are $`\pm h_z=\pm \mathrm{ln}\lambda `$ and $`\pm h_\varphi =\pm \mathrm{ln}B`$ for $`\stackrel{}{\varphi }`$ motion. In terms of the Lyapunov exponents, $`D_S`$ is
$$D_S=\{\begin{array}{cc}42\frac{h_z}{h_\varphi },& \text{ for }h_z/h_\varphi 1/2,\\ 3,& \text{ for }h_z/h_\varphi >1/2.\end{array}$$
(43)
Since $`h_z=1/\tau `$ for (38) and (39) we see that, for $`h_z/h_\varphi 1/2`$, Eq. (43) is the same as Eq. (34). Also, the system (38), (39) has no finite time Lyapunov exponent fluctations, and, thus, the information dimension and the box-counting dimensions are the same. Hence, $`D_S=3`$ when $`SM`$ is smooth ($`h_z/h_\varphi >1/2`$). This is analogous to the situation $`r=1`$ and Eq. (24) of Sec. 3. \[The lack of finite time Lyapunov exponent fluctuations for our model system (38), (39) is reflected in the fact that all periodic orbits in $`\mathrm{\Lambda }`$ have precisely the same Lyapunov exponents, namely $`\pm \mathrm{ln}\lambda `$ and $`\pm \mathrm{ln}B`$.\]
### 6.2 Basin Boundary for a Map Modeling the Tilted Ellipsoid Billiard
We now wish to investigate the structure of the stable manifold when we give the ellipsoid a small tilt. Again, we adopt the approach of Sec. 4.2: we obtain a rigorous result by utilizing a simpler map model that preserves the basic features of the tilted ellipsoid case. Here we again use (38) but we now modify (39) to incorporate the main effect of a small tilt. The effect of this modification is to destroy the invariance of $`\stackrel{}{z}=0`$. Thinking of the first non-zero term in a power series expansion for small $`\left|\stackrel{}{z}\right|`$, this invariance results because the first expansion term is linear in $`\stackrel{}{z}`$; i.e., the $`\stackrel{}{z}`$-independent term in the expansion is exactly zero. When there is tilt this is not so. Thus we replace (39) by
$$\left[\begin{array}{c}z_{n+1}\\ v_{n+1}\end{array}\right]=\left[\begin{array}{cc}\lambda & f(\stackrel{}{\varphi }_n)\\ 0& \lambda ^1\end{array}\right]\left[\begin{array}{c}z_n\\ v_n\end{array}\right]+\left[\begin{array}{c}f_z(\stackrel{}{\varphi }_n)\\ f_v(\stackrel{}{\varphi }_n)\end{array}\right].$$
(44)
The simplest version of (44) which still has the essential breaking of $`\stackrel{}{z}\stackrel{}{z}`$ symmetry is the case where $`f=f_v=0`$. Because setting $`f=f_v=0`$ greatly simplifies the analysis, we consider this case in what follows (we do not expect our conclusion to change if $`f,f_v0`$). Thus, we have
$$\left[\begin{array}{c}z_{n+1}\\ v_{n+1}\end{array}\right]=\left[\begin{array}{cc}\lambda & 0\\ 0& \lambda ^1\end{array}\right]\left[\begin{array}{c}z_n\\ v_n\end{array}\right]+\left[\begin{array}{c}f_z(\stackrel{}{\varphi }_n)\\ 0\end{array}\right].$$
(45)
The problem of finding the stable manifold for the invariant set of the map, (38) and (45), is now the same as for the previously considered case of Eqs. (38) and (40) \[compare the equation $`z_{n+1}=\lambda z_n+f_z(\stackrel{}{\varphi }_n)`$ with (40)\]. Thus, making use of this equivalence we can immediately write down the equation for the stable manifold in the four-dimensional state space $`(z,v,\varphi ,\omega )`$ as
$$z=v\underset{i=0}{\overset{\mathrm{}}{}}\lambda ^if_z(C^i\stackrel{}{\varphi }),$$
(46)
which is obtained from (41) using the replacements $`\nu z/v`$, $`\lambda ff_z`$, and $`\lambda ^2\lambda `$. The rigorous results of Ref. again show that this is a continuous, nowhere-differentiable surface for almost all $`\lambda `$ in $`1<\lambda <B`$, and, furthermore, when this is so ($`\mathrm{ln}\lambda /\mathrm{ln}B=h_z/h_\varphi 1`$) we have
$$D_S=4\frac{h_z}{h_\varphi }.$$
(47)
Also, $`D_S=3`$ when $`h_z/h_\varphi >1`$. Since $`h_z=1/\tau `$, this is the same as Eq. (33).
## 7 Derivation of Dimension Formulae
### 7.1 $`D_S`$ for Typical Systems
While the dimension formulae presented in Sec. 2 and derived in were found to be valid for the two-dimensional map of Sec. 3 and the ellipsoid scatterer of Secs. 4-6, the derivation given in does not apply to these systems. In the derivation of , a higher-dimensional generalization of horseshoe-like dynamics is assumed. When this type of dynamics is responsible for a fractal basin boundary, the boundary looks, locally, like a Cantor set of smooth surfaces. The examples of this paper do not possess horseshoe-like dynamics. Their dynamics produces basin boundaries which consist of a single continuous surface; see Figs. 1, 5, and 7. We proceed to derive dimension formulae using a more general argument not restricted to an assumed particular type of boundary. The resulting formulae are identical to those derived in and stated in Sec. 2.
Let $`R`$ be the portion of the state space which contains all points within $`ϵ`$ of a nonattracting, ergodic, invariant set, $`\mathrm{\Lambda }`$ of an $`M`$-dimensional map, $`P`$. Let $`sm`$ be the portion of the stable manifold of $`\mathrm{\Lambda }`$ which is contained in $`R`$. The points in $`R`$ are within $`ϵ`$ of $`sm`$ since $`\mathrm{\Lambda }`$ is a subset of $`sm`$. If we sprinkle a large number, $`N_0`$, of orbit initial conditions in $`R`$, then the number of orbits left in $`R`$ after $`n1`$ iterates is assumed to scale like (3) $`N_n/N_0e^{n/\tau }`$. Let the map, $`P`$, have Lyapunov exponents
$$h_U^+h_{U1}^+\mathrm{}h_1^+>0>h_1^{}\mathrm{}h_{S1}^{}h_S^{},$$
where $`U+S=M`$.
We will use the box-counting definition of dimension
$$N(ϵ)ϵ^D,$$
(48)
where $`N(ϵ)`$ is the minimum number of $`M`$-dimensional hypercubes (“boxes”) needed to cover $`sm`$ and $`D`$ is its box-counting dimension.
We wish to develop a covering of $`sm`$ using small boxes and determine how the number of boxes in that covering scales as the size of the boxes is decreased. We will look at how the linearized system dynamics distorts a typical small box. This will help us determine how the Lyapunov exponents are related to the dimension of $`sm`$. Since we assume a smooth map, the dimension of the stable manifold of the map is equal to that of $`sm`$.
Cover $`sm`$ with boxes which are of length $`ϵ`$ on each of their $`M`$ sides. Call this set of boxes $`C`$. The number of boxes in $`C`$ is denoted $`N_0`$. Iterate each box forward $`n`$ steps, where $`n`$ is large, but not so large that the linearized dynamics do not apply to the boxes. Call the set of iterated boxes $`P^n(C)`$. A typical iterated box in $`P^n(C)`$ is a distorted $`M`$-dimensional parallelopiped and has dimensions
$$ϵe^{nh_U^+}\times \mathrm{}\times ϵe^{nh_1^+}\times ϵe^{nh_1^{}}\times \mathrm{}\times ϵe^{nh_S^{}}.$$
Each parallelopiped intersects $`sm`$ since its preimage (a box) did. (The set $`P^n(C)`$ does not, however, cover all of $`sm`$ since points on $`sm`$ move closer to $`\mathrm{\Lambda }`$ on forward iteration.)
To construct a refined covering of $`sm`$ we begin by covering each parallelopiped of $`P^n(C)`$ with slabs of size
$$\stackrel{U\text{ factors}}{\overline{ϵ\times \mathrm{}\times ϵ}}\times ϵe^{nh_1^{}}\times \mathrm{}\times ϵe^{nh_S^{}}.$$
There are roughly $`\mathrm{exp}[n(h_U^++\mathrm{}+h_1^+)]`$ such slabs. Only $`e^{n/\tau }`$ of these slab are within $`ϵ`$ of $`sm`$. Let $`C^{}`$ denote the set of $`N^{}N_0\mathrm{exp}[n(h_U^++\mathrm{}+h_1^+1/\tau )]`$ slabs needed to cover the part of $`sm`$ lying in $`P^n(C)`$. Let us define $`H=h_U^++\mathrm{}+h_1^+1/\tau `$ so that $`N^{}N_0e^{nH}`$.
Iterate each of the $`N^{}`$ slabs backward $`n`$ steps. We now have the set $`P^n(C^{})`$. It contains $`N^{}`$ parallelopipeds of size
$$ϵe^{nh_U^+}\times \mathrm{}\times ϵe^{nh_1^+}\times \stackrel{\text{S factors}}{\overline{ϵ\times \mathrm{}\times ϵ}}.$$
The set $`P^n(C^{})`$ forms a covering of $`sm`$. To calculate the dimension (which is defined in terms of boxes) we cover the $`N^{}`$ parallelopipeds in $`P^n(C^{})`$ with boxes which are $`ϵ_j=ϵe^{nh_{j+1}^+}`$ on each side. (We will choose the value of the index $`j`$ below.) The number of boxes needed to cover $`sm`$, for boxes of size $`ϵ_j`$, scales as
$$N^{\prime \prime }N^{}\frac{ϵe^{nh_j^+}}{ϵ_j}\times \frac{ϵe^{nh_{j1}^+}}{ϵ_j}\times \mathrm{}\times \frac{ϵe^{nh_1^+}}{ϵ_j}\times \left(\frac{1}{ϵ_j}\right)^S.$$
To cover the slabs with these boxes we need a factor of $`1`$ boxes along each direction of a slab which is shorter than the edge length of a box and a factor of $`ϵe^{nh_k^+}/ϵ_j`$ boxes along each direction (here, the $`k^{\text{th}}`$ direction) which is longer than the edge length of a box. In terms of $`N_0`$,
$$N^{\prime \prime }N_0\mathrm{exp}\{n[(S+j)h_{j+1}+Hh_1^+h_2^+\mathrm{}h_j^+1/\tau ]\}.$$
To compute the dimension of $`sm`$, we compare $`N(ϵ)N_0`$ to $`N(ϵ_j)N^{\prime \prime }`$ using (48). This gives $`N(ϵ_j)/N(ϵ)(ϵ/ϵ_j)^D\mathrm{exp}(nDh_{j+1}^+)`$ which yields the following $`j`$-dependent dimension esitmate,
$$D(j)=S+j+\frac{Hh_1^++h_2^++\mathrm{}h_j^+}{h_{j+1}^+}.$$
(49)
Our definition of box-counting dimension, (48), requires us to find the minimum number of boxes needed to cover the set. Since we are certain that the set is covered, (49) yields an upper bound for the dimension for any $`j`$. Thus, to find the best estimate of those given by (49), we minimize $`D(j)`$ over the index $`j`$. Comparing $`D(j)`$ to $`D(j+1)`$ yields the condition
$$h_1^++\mathrm{}+h_J^++h_J^+Hh_1^++\mathrm{}+h_J^+,$$
where $`J`$ is the best choice for $`j`$ (i.e., the choice giving the minimum upper bound). The conjecture is that this minimum upper bound $`D(J)`$ from (49) is the actual dimension of $`sm`$ for typical systems. $`D(J)`$ is the same as the result Eq. (8) presented in Sec. 2 for $`D_S`$. The derivation of the dimension formula for the unstable manifold, $`D_U`$, is similar to that just presented for the stable manifold.
\[The derivation just presented gives the *information* dimension of the stable manifold, not the box-counting dimension. We were considering the sizes of typical boxes in the system and covered only those. This leaves boxes that have atypical stretching rates for large $`n`$ (a box containing a periodic point, for example, will, in general, not stretch at the rates given by the Lyapunov exponents) unaccounted for. What we have actually computed is the box-counting dimension of most of the measure, which is the information dimension , of the stable manifold. See for discussion of this point.\]
As derived, $`D(J)`$ gives an upper bound on the dimension. We saw in the ellipsoid example of Secs. 4-6 that the $`zz`$ symmetry of the system lead to a special geometry for its stable manifold (Fig. 7) and the formula just derived did not apply (although it was an uppder bound). It is the conjecture of that this formula gives not the upper bound, but the exact dimension of the stable manifold for typical systems. This is supported by our results for the tilted ellipsoid example.
### 7.2 $`D_S`$ for the Untilted (Atypical) Case
The derivation above for typical systems gives the wrong dimension formula for the stable manifold of the (atypical) untilted ellipsoid system. This typical-system derivation consistently overcounts the boxes needed to cover the atypical stable manifold of the untilted ellipsoid system. Sec. 6.1 showed that the stable manifold had a ray-like structure (cf. Fig. 7) induced by the $`z`$-direction reflection symmetry of the system. If we account for this structure we can modify the typical-system derivation so that we do not overcount boxes. The modified derivation produces the correct formula.
Recall that our untilted ellipsoid map (which we call $`P`$ here) has Lyapunov exponents $`\pm h_z`$ and $`\pm h_\varphi `$ and decay time $`\tau `$. We consider the case
$$h_\varphi >h_z+\tau >0$$
Let us cover the region of state space which is within $`ϵ/2`$ of $`\mathrm{\Lambda }`$ with $`N_0`$ boxes with edge lengths $`ϵ\times ϵ\times ϵ\times ϵ`$. Call this set $`C`$. The map, $`P`$, is linear in $`\stackrel{}{z}`$ (in particular, the $`\stackrel{}{z}`$ portion is of the form $`\stackrel{}{z}_{n+1}=DM(\stackrel{}{\varphi })\stackrel{}{z}_n`$) so that the graph of $`SM`$ in $`\stackrel{}{z}`$ space for a given value of $`\stackrel{}{\varphi }`$ is a straight line through $`\stackrel{}{z}=0`$, i.e. $`z=vg(\stackrel{}{\varphi })`$. We denote by $`sm`$ the portion of the stable manifold which is contained in $`C`$. Since $`SM`$ contains $`\mathrm{\Lambda }`$, each box in $`C`$ intersects $`sm`$.
Iterate these boxes forward $`n1`$ steps. They become a set distorted of parallelopipeds which we call $`P^n(C)`$. Each parallelopiped has dimensions
$$ϵe^{nh_\varphi }\times ϵe^{nh_z}\times ϵe^{nhz}\times ϵe^{nh_\varphi }$$
and intersects $`sm`$. We cover $`P^n(C)`$ by $`N_0\mathrm{exp}[n(h_z+h_\varphi )]`$ parallelopipeds with dimensions
$$ϵ\times ϵ\times ϵe^{nh_z}\times ϵe^{nh_\varphi },$$
and discard all of these parallelopipeds which do not intersect $`sm`$. There are $`N^{}N_0\mathrm{exp}[n(h_z+h_\varphi 1/\tau )]`$ parallelopipeds remaining which cover $`sm`$. The portions of these paralleopipeds which are in $`ϵ>\left|z\right|>ϵe^{nh_z}`$ do not contain $`sm`$. Suppose some portion of $`sm`$ did fall in $`ϵ>\left|z\right|>ϵe^{nh_z}`$. Since $`P`$ is linear in $`\stackrel{}{z}`$, this outlying portion of $`sm`$ would, upon $`n`$ reverse iterations, map to the region $`ϵe^{nh_z}>\left|z\right|>ϵ`$, which contradicts the definition of $`sm`$ given above (i.e. $`sm`$ is within $`ϵ/2`$ of $`\mathrm{\Lambda }`$). Therefore we can discard the portion of each parallelopiped which lies is $`ϵ>\left|z\right|>ϵe^{nh_z}`$. These $`N^{}`$ parallelopipeds now have dimensions
$$ϵ\times ϵe^{nh_z}\times ϵe^{nhz}\times ϵe^{nh_\varphi },$$
and are denoted $`C^{}`$. We now iterate the parallelopipeds in $`C^{}`$ backward $`n`$ times and call the resulting set $`P^n(C^{})`$. This set contains $`N^{}`$ parallelopipeds with dimensions
$$ϵe^{nh_\varphi }\times ϵe^{2nh_z}\times ϵ\times ϵ.$$
We cover $`P^n(C^{})`$ (which covers $`sm`$) with hypercubes which have as their edge length $`ϵe^{nh_\varphi }`$. The number of hypercubes needed is $`N^{\prime \prime }N_0\mathrm{exp}[n(h_z+h_\varphi 1/\tau 2h_z+3h_\varphi )]`$. The (information) dimension is again found by comparing the number of $`ϵ`$-sized hypercubes needed to cover $`sm`$ to the number of $`ϵe^{nh_\varphi }`$ sized hypercubes needed to cover $`sm`$,
$$\frac{N^{\prime \prime }}{N_0}\left(\frac{ϵe^{nh_z}}{ϵ}\right)^D,$$
which yields
$$D=4\frac{h_z+1/\tau }{h_\varphi }.$$
This expression is between three and four for $`(h_z+1/\tau )/h_\varphi <1`$.
## 8 Discussions and Conclusions
In summary, this paper has tested and illustrated formulae giving the fractal dimension of nonattracting chaotic sets in terms of their Lyapunov exponents and decay time. We have done this using two examples, one, a noninvertible, two-dimensional map with two expanding (positive) Lyapunov exponents, and the other, a chaotic scattering billiard in three spatial dimensions. The first example is particularly useful for understanding the transient measure. The second example provides a striking illustration of the issue of atypicality, and also has the benefit of potential physical realization (see below). Another point is that our second example provides the first known case of a chaotic scatterer with a fractal basin boundary that is a continuous, nowhere-differentiable surface (such structures have been previously discussed for basins of attraction of dissipative systems but not in chaotic scattering). Finally, in Sec. 7 we have provided arguments that are more general than those previously given for the dimension formulae of Sec. 2.
To conclude, we now offer some discussion on the possible experimental realization of our billiard example. Our point here is that such billiard systems offer a particularly nice way of experimentally realizing and studying basin boundaries. In the past there have been very few experiments that have attempted to study basin boundaries. The problem is the experimental difficulty of carrying out the technique used in numerical work: namely, choosing a grid of initial conditions and seeing where each orbit from each initial condition goes. To do this experimentally, one would have to precisely prepare an initial condition, run the experiment to see where the orbit with that initial condition goes, and do this for each initial condition. The difficulties with this approach are that it is often not possible to prepare initial conditions either on a fine enough grid or sufficiently precisely to observe small scale basin boundary structures and that running the experiment many times can be time consuming. (In addition, experimental parameters may drift over the course of the multiple experimental runs.) Nevertheless, in one case this program was successfully carried out. In that paper an electrical circuit was used as the experimental chaotic system, and the authors were able to experimentally demonstrate a riddled basin of attraction (a type of basin for which final outcomes are particularly difficult to predict; see Refs. , ). In another work, Cusumano and Kimble have devised a new experimental procedure for studying basin boundaries. This procedure has advantages over the straightforward method of preparing initial conditions on a grid. Using this method, Cusumano and Kimble have successfully mapped out basin boundaries for experimental dissipative mechanical systems. To our knowledge, these two works (and a billiard we discuss below) are so far the only experimental investigations of fractal basin boundaries. This is in contrast to the large number of papers that have experimentally realized fractal structure in chaotic attractors. This disparity in the situations of attracting and nonattracting chaotic sets seems to be largely due to the disparity in the ease of experimental realization for the two cases.
We believe that billiard systems offer a convenient avenue for experimental investigation of the various general types of basin boundaries. In particular, we suggest an optical realization of chaotic scattering billiards in which the billiard surface is made to be optically reflecting, and light rays play the role of orbits. For the scatterer of Sec. 4-6, for example, imagine that we have a tube with cross section as shown in Fig. 4(b) and mirrored inner walls. The length of the tube is $`L>2r_{}`$, where $`r_{}`$ is the major radius of an ellipsoid with a mirrored surface. Imagine that the tube is oriented vertically with its bottom, open end placed on a red surface, and that the ellipsoid is suspended in the tube. In this configuration, an observer looking in the top of the tube sees multiple reflections of the red surface that is at the bottom of the tube. Since rays whose directions are reversed retrace the same path, we can think of orbits from initial conditions starting on the retina of the observer’s eye passing through his pupil, bouncing around in the scatterer, and then exiting either through the bottom (red) or the top (not red). Thus, the observed boundary of the red region is precisely the basin boundary that we have been discussing in Secs. 4-6. Note that to observe that boundary we did not have to prepare many initial conditions and repeat the experiment many times. Rather a single image giving a global picture of the basin can be formed (replacing the observer’s eye by a camera). We have already applied this approach to a chaotic scatterer formed from four mirrored spheres to experimentally study a different general type of basin boundary structure known as a Wada boundary (a Wada boundary is a boundary separating three or more regions such that every boundary point is a boundary point for all regions). In the case of our ellipsoid billiard, such an experiment would represent the first experimental realization of a fractal basin boundary with the structure of a continuous, nowhere-differentiable surface, one of the small number of typical possible basin boundary types .
We thank J. A. Yorke and B.R. Hunt for discussion. This work was supported by ONR and DOE.
## Appendix A: Slow Convergence of the Dimension Near the Transition Between Fractal and Nonfractal Behavior
It is difficult to numerically measure the box-counting dimension of the chaotic repellors studied in Sec. 3 and Secs. 4-6 for some range of the system parameter. When the parameter is set near the value at which the repellor makes the transition from a fractal to a non-fractal function the number of boxes needed to cover the set scales as
$$N(ϵ)\frac{[\mathrm{ln}(1/ϵ)]^{\frac{1}{2}}}{ϵ^{D_0}}$$
(A1)
with $`ϵ`$ (the width of a box) down to very small $`ϵ`$. In (A1), $`D_0=1`$ in the case of the two-dimensional map and $`D_0=3`$ in the case of the ellipsoid system. Due to the logarithmic term in (A1) accurate numerical determination of $`D_0`$ can be very demanding. In the following we demonstrate that (A1) holds near the transition point between a fractal and a nonfractal. In particular, we consider (13)-(15) \[$`D_0=1`$ in Eq. (A1)\] in the case where $`\lambda _1=\lambda _2=\lambda `$. Using the same reasoning as for the derivation of Eq. (41) we obtain for the invariant set the curve
$$y(x)=\frac{\eta }{2\pi }\underset{i=0}{\overset{\mathrm{}}{}}\lambda _i^1\mathrm{sin}(2\pi 2^ix).$$
(A2)
This result is also derived in Ref. . The box-counting dimension of the graph of $`y(x)`$ versus $`x`$ is $`D_0=2[\mathrm{ln}\lambda /\mathrm{ln}2]`$ for $`\lambda 2`$ and $`D_0=1`$ for $`\lambda 2`$. Thus, the graph of $`y(x)`$ is a fractal for $`\lambda 2`$ and is non-fractal for $`\lambda 2`$. We are interested in the transition region, $`\lambda 2`$, for which we wish to show that (A1) applies.
To calculate the box-counting dimension of (A2) we need to know how the number of $`ϵ`$-width boxes needed to cover the set scales with $`ϵ`$. If we let the sum in (A2) run from $`i=0`$ to $`i=n`$, the resulting approximation to $`y(x)`$ is a smooth, finite length curve:
$$y_n(x)=\frac{\eta }{2\pi }\underset{i=0}{\overset{n}{}}\lambda ^i\mathrm{sin}(2\pi 2^ix).$$
(A3)
As $`n`$ increases, the ratio of the height to the width of the oscillations added by each term in the sum increases. For the purpose of estimating $`N(ϵ)`$ for some fixed value of $`ϵ`$ it suffices to consider (A3) with $`n`$ sufficiently large. A practical size for the boxes in a covering is one which scales as the width of the smallest oscillations of (A2): $`ϵ_n2^n`$. This way we can be sure to resolve the smallest details of the approximation curve (A3) and be sure that higher order approximations will not significantly alter our count of the number of boxes in the covering. (Oscillations introduced by the $`n+1^{\text{st}}`$ term, for example, will be roughly half the width of a box of size $`ϵ_n2^n`$, or $`2^{(n+1)}`$ .)
The length of the curve (A3) is
$$\mathrm{}_n=_0^1𝑑x\sqrt{1+[y_n^{}(x)]^2}\sqrt{1+[y_n^{}(x)]^2},$$
where $``$ denotes an average over $`x`$. For large $`n`$ and $`\lambda 2`$, $`\left|y_n^{}(x)\right|`$ is typically large and $`\mathrm{}_n\left|y^{}(x)\right|`$.
To determine $`\mathrm{}_n`$ we examine
$$y_n^{}(x)=\eta \underset{i=0}{\overset{n}{}}\left(\frac{2}{\lambda }\right)^i\mathrm{cos}(2\pi 2^ix).$$
For $`\frac{2}{\lambda }`$ sufficiently large, the largest term in this sum dominates and $`y_n^{}(x)\left(\frac{2}{\lambda }\right)^n`$. In this case, $`\mathrm{}_n\left(\frac{2}{\lambda }\right)^n`$ and $`N(ϵ_n)\frac{\mathrm{}_n}{ϵ_n}`$. Taking $`ϵ=2^n`$, or $`n=\mathrm{ln}ϵ/\mathrm{ln}2`$ we have that $`N(ϵ)ϵ^{(2\mathrm{ln}\lambda /\mathrm{ln}2)}`$. The scaling is a power law in $`ϵ`$. (It gives a box-counting dimension of $`D_0=2\mathrm{ln}\lambda /\mathrm{ln}2`$.)
For $`\lambda <2`$, but $`\lambda 2`$, the last term will not dominate the sum, but the last several, up to $`n_x`$, terms will dominate. To estimate $`n_x`$, set $`\left(\frac{2}{\lambda }\right)^{n_x}`$ equal to a constant factor $`K>1`$ where $`K1`$ is order one (e.g., $`K=2`$). This gives $`n_x(2\lambda )^11`$. If $`nn_x`$, then all the coefficients $`(2/\lambda )^i`$ in the sum $`y_n^{}(x)`$ are of the same order, $`y_n^{}(x)_{i=0}^n\mathrm{cos}(2\pi 2^ix).`$ For a given value of $`x`$ and different values of $`i1`$ the arguments of the cosine terms in the sum are very different and vary rapidly with $`x`$. Thus, we regard the terms in the sum as random. In this case $`y_n^{}\sqrt{n}`$, and, thus, $`\mathrm{}_n\sqrt{n}`$. Again, we say that $`N(ϵ_n)\frac{\mathrm{}_n}{ϵ_n}`$, but now $`N(ϵ_n)\sqrt{n}/ϵ_n\sqrt{\mathrm{ln}ϵ_n}/ϵ_n`$, which is (A1). If $`n>n_x`$, then $`\left|y_n^{}(x)\right|n_x^{1/2}\left(\frac{2}{\lambda }\right)^n`$.
For the two-dimensional example studied numerically in Sec. 3.3, the numerical results are shown in Fig. 2. Due to the effect discussed in this appendix we have not plotted data near the kink at $`\lambda _1=\lambda _b`$. The results for the ellipsoid system are discussed in Sec. 5 and shown in Fig. 6. A typical two-dimensional slice of the basin boundary for this system is a continuous, nowhere-diffferentiable curve similar to the chaotic saddle of the two-dimensional system. We thus have the same difficulty in measuring the dimension near $`h_\varphi \tau =1`$, Fig. 6(a) and $`h_\varphi =h_z+\tau ^1`$ Fig. 6(b).
## Appendix B: Dynamics of the Orientation of $`\stackrel{}{z}`$
Rewriting (35) as,
$`z_{n+1}`$ $`=`$ $`m_{11}(\stackrel{}{\varphi }_n)z_n+m_{12}(\stackrel{}{\varphi }_n)v`$ (B1)
$`v_{n+1}`$ $`=`$ $`m_{21}(\stackrel{}{\varphi }_n)z_n+m_{22}(\stackrel{}{\varphi }_n)v_{n,}`$ (B2)
eliminating $`z`$ in favor of $`\nu =z/v`$, and dividing the first of the above equations by the second, $`v_{n+1}`$ and $`v_n`$ cancel to yield an evolution equation for $`\nu `$,
$$\nu _{n+1}=\frac{m_{11}(\stackrel{}{\varphi }_n)\nu _n+m_{12}(\stackrel{}{\varphi }_n)}{m_{21}(\stackrel{}{\varphi }_n)\nu _n+m_{22}(\stackrel{}{\varphi }_n)}.$$
(B3)
Equations (36) and (B3) now constitute a 3D map. Our transformation from (B1), (B2) eliminates the symmetry of the original system, and it is reasonable to now assume that the information dimension of the natural transient measure for the 3D map (36) and (B3) is given by the result for a typical system. Thus, we have that the information dimension of the stable manifold of the invariant set for the 3D map is
$$\widehat{D}_S=3(\widehat{h}_\varphi \widehat{\tau })^1,\text{for}\widehat{h}_\varphi \widehat{\tau }1.$$
(B4)
Where the superscribed circumflex denotes quantities calculated for the natural transient measure of the 3D map, (36) and (B3). Since the information dimension $`\widehat{D}_S`$ is a lower bound on the box-counting dimension for the map (36) and (B3), and since the box-counting dimension of $`SM`$ for the original 4D map is one plus that for the 3D map (this follows from the structure shown in Fig. 7), we have that $`\widehat{D}_S+1`$ \[with $`\widehat{D}_S`$ given by (B4)\] is a lower bound for the box-counting dimension of $`SM`$ for the 4D map. We emphasize, however, that $`\widehat{D}_S+1`$ is different from $`D_S`$ for the 4D system. This is because of the difference in the natural transient measures for the 3D and 4D maps. (In contrast, the box-counting dimensions do not depend on the measures.) In particular, as shown below, as compared to the 4D map natural transient measure, the natural transient measure for the 3D map more strongly weights $`\stackrel{}{\varphi }`$ values which experience slower repulsion from the invariant set $`\mathrm{\Lambda }`$.
To see this we consider the Lyapunov exponent corresponding to the evolution of a differential perturbation in $`\nu `$ which we denote $`\delta \nu `$. Differentiating (B3) we have
$$\delta \nu _{n+1}/\delta \nu _n=\overline{d}(\stackrel{}{\varphi }_n)(m_{21}\nu _n+m_{22})^2,$$
(B5)
where $`\overline{d}(\stackrel{}{\varphi }_n)`$ the determinant of (B1), (B2), $`\overline{d}(\stackrel{}{\varphi }_n)=m_{11}m_{22}m_{12}m_{21}=1`$. Noting that $`\nu _n=z_n/\nu _n`$ and comparing (B5) with (B2), we have
$$\delta \nu _{n+1}/\delta \nu =\overline{d}(\stackrel{}{\varphi }_n)(v_n/v_{n+1})^2.$$
(B6)
For points on $`SM`$, $`v_n`$ decreases exponentially with time (i.e., it approaches the invariant set $`\mathrm{\Lambda }`$). Hence,
$$\delta \nu _n/\delta \nu _0=Q(v_n/v_0)^2\mathrm{exp}[2n\stackrel{~}{h}_z(\stackrel{}{\varphi }_0,n)],$$
(B7)
where $`Q=\overline{d}(\stackrel{}{\varphi }_0)\overline{d}(\stackrel{}{\varphi }_1)\mathrm{}\overline{d}(\stackrel{}{\varphi }_{n1})`$ and $`\stackrel{~}{h}_z(\stackrel{}{\varphi },n)`$ with $`\stackrel{}{\varphi }=\stackrel{}{\varphi }_0`$ is the expanding $`\stackrel{}{z}`$, finite-time Lyapunov exponent for the 4D system. The estimate, $`\mathrm{exp}(2n\stackrel{~}{h}_z)`$, in (B7) follows since, due to the Hamiltonian nature of the problem, $`Q^{1/n}1`$ as $`n\mathrm{}`$. Thus, for a given value of $`\stackrel{}{\varphi }`$, in the 3D system, the expansion away from the invariant set is at the rate $`2\stackrel{~}{h}_z(\stackrel{}{\varphi },n)`$ rather than at the rate $`\stackrel{~}{h}_z(\stackrel{}{\varphi },n)`$, applying to the 4D system. For the example of (39) there are no finite time Lyapunov exponent fluctuations \[$`\stackrel{~}{h}_z(\stackrel{}{\varphi },n)=\mathrm{ln}\lambda `$ independent of $`n`$ and $`\stackrel{}{\varphi }`$\], so that the 3D and 4D natural transient measures are the same, both are uniform in $`\stackrel{}{\varphi }`$.
In this case,
$$\widehat{\tau }=\tau /2=(\mathrm{ln}\lambda )^1.$$
(B8)
More generally, the finite time exponents, $`\stackrel{~}{h}_z(\stackrel{}{\varphi },n)`$, fluctuate as $`\stackrel{}{\varphi }`$ varies. The faster expansion rate away from $`\mathrm{\Lambda }`$ (namely $`2\stackrel{~}{h}_z`$) in the 3D case means that $`\stackrel{}{\varphi }`$ values with slower expansion rates are more strongly weighted. Thus, for a case such as our ellipsoid problem we expect that
$$\widehat{\tau }>\tau /2\widehat{h}_\nu <2h_z,$$
where $`\widehat{h}_\nu `$ is the infinite time Lyapuonv exponent for the 3D map measure and for perturbations in $`\nu `$.
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# Dynamics of Nonlocality for A Two-Mode Squeezed State in Thermal Environment
## I INTRODUCTION
Quantum nonlocality is one of the most profound features of quantum mechanics. It enables current developments of quantum information theory encompassing quantum teleportation, quantum computation, and quantum cryptography . There have been studies on tests of quantum nonlocality versus local realism. Bell suggested an inequality which any local hidden variable theory must obey . Several types of Bell’s inequalities were derived in terms of two-body correlation functions of two measurement variables at distant places for the test of quantum nonlocality of a spin-1/2 or SU(2) system.
A spin-1/2 system can be utilized as a qubit for quantum computation. Quantum nonlocality of the spin-1/2 system is required as a quantum channel to teleport an unknown qubit state . In fact, it is possible to teleport not only a two-dimensional spin-1/2 quantum state but also a $`N`$-dimensional state and a continuous-variable state . The type of the quantum channel depends on the dimension of Hilbert space of an unknown quantum state. For the teleportation of a continuous-variable state, the quantum channel should be in an entangled continuous-variable state such as the two-mode squeezed state . Recently, practical implementation of the quantum teleportation for the continuous-variable state has been realized experimentally using the two-mode squeezed field . In the quantum teleportation, the most important ingredient is the quantum nonlocality of the channel which can be easily destroyed in the nature. In this paper we are interested in how the thermal environment affects the quantum nonlocality of the two-mode squeezed field.
Quantum nonlocality of an entangled continuous-variable state has been discussed using the Schmidt form for entangled nonorthogonal states and the quadrature-phase homodyne measurement . A given state is nonlocal when it violates any Bell’s inequality. In fact, a state does not have to violate all the possible Bell’s inequalities to be considered quantum nonlocal. A state is quantum nonlocal for the given Bell’s inequality which is violated by the measurement of the state. Banaszek and Wódkiewicz defined Bell’s inequality based on the parity measurement and they found that the two-mode squeezed state violates the Bell’s inequality, showing quantum nonlocality .
In this paper we study the dynamic behavior of the quantum nonlocality based on the parity measurement for the two-mode squeezed state in the thermal environment. A measurement of the degree of quantum nonlocality is defined here by the maximal violation of Bell’s inequality. The nonlocality is stronger for the squeezed state with a larger degree of squeezing. It is found that the nonlocality disappears more rapidly in the thermal environment as the initial state is squeezed more.
This paper is organized as follows. In Sec. II, Bell’s inequality based on the parity measurement is discussed. The parity measurement is directly related to the Wigner function. In Sec. III, the two-mode master equation is solved for the dynamics of the Wigner function of the initial two-mode squeezed state. The convolution theory is utilized in the solution . We investigate the dynamic behavior of the quantum nonlocality measured by the maximum violation of Bell’s inequality for the two-mode squeezed state in the thermal environment in Sec. IV.
## II Bell’s inequality by parity measurement
It is important to choose the type of measurement variables when testing nonlocality for a given state. In the original EPR gedanken experiment , Einstein, Podolsky, and Rosen considered the positions (or the momenta) of two particles as the measurement variables to discuss the two-body correlation. Bell argued that the EPR wave function does not exhibit nonlocality because its Wigner function $`W(x_1,p_1;x_2,p_2)`$ is positive everywhere, allowing the description by a local hidden variable theory. Munro showed that various types of Bell’s inequalities are not violated in terms of the homodyne measurements of two particles . To the contrary, Banaszek and Wódkiewicz examined even and odd parities as the measurement variables and showed that the EPR state and the two-mode squeezed state are nonlocal in the sense that they violate Bell’s inequalities such as Clauser and Horne inequality and Clauser-Horne-Shimony-Holt inequality.
The even and odd parity operators, $`\widehat{O}_e`$ and $`\widehat{O}_o`$, are the projection operators to measure the probabilities of the field having even and odd numbers of photons, respectively:
$$\widehat{O}_e=\underset{n=0}{\overset{\mathrm{}}{}}|2n2n|;\widehat{O}_o=\underset{n=0}{\overset{\mathrm{}}{}}|2n+12n+1|.$$
(1)
The Wigner function at the origin of phase space for a state of the density operator $`\widehat{\rho }`$ is proportional to the mean parity :
$$W(0)=(2/\pi )\text{Tr}\left[(\widehat{O}_e\widehat{O}_o)\widehat{\rho }\right].$$
(2)
Further, the Wigner function $`W(\alpha )`$ at the phase point $`\alpha `$ is the mean parity for the displaced original state
$$W(\alpha )=(2/\pi )\text{Tr}[(\widehat{O}_e\widehat{O}_o)\widehat{D}(\alpha )\widehat{\rho }\widehat{D}^{}(\alpha )]$$
(3)
where $`\widehat{D}(\alpha )`$ is the displacement operator .
So far the argument has been confined to the parity measurement of a single-mode field. As the quantum nonlocality can be discussed for two-mode fields, we thus define the quantum correlation operator based on the joint parity measurements:
$`\widehat{\mathrm{\Pi }}^{ab}(\alpha ,\beta )=\widehat{\mathrm{\Pi }}_e^a(\alpha )\widehat{\mathrm{\Pi }}_e^b(\beta )\widehat{\mathrm{\Pi }}_e^a(\alpha )\widehat{\mathrm{\Pi }}_o^b(\beta )`$ (4)
$`\widehat{\mathrm{\Pi }}_o^a(\alpha )\widehat{\mathrm{\Pi }}_e^b(\beta )+\widehat{\mathrm{\Pi }}_o^a(\alpha )\widehat{\mathrm{\Pi }}_o^b(\beta )`$ (5)
where the superscripts $`a`$ and $`b`$ denote the modes and the displaced parity operator, $`\mathrm{\Pi }_{e,o}(\alpha )`$, is defined as
$$\widehat{\mathrm{\Pi }}_{e,o}(\alpha )=\widehat{D}(\alpha )\widehat{O}_{e,o}\widehat{D}^{}(\alpha ).$$
(6)
The displaced parity operator acts like a rotated spin projection operator in the spin measurement. We can easily derive that the local hidden variable theory imposes the following Bell’s inequality
$`|B(\alpha ,\beta )||\widehat{\mathrm{\Pi }}^{ab}(\alpha ,\beta )+\widehat{\mathrm{\Pi }}^{ab}(\alpha ,\beta ^{})`$ (7)
$`+\widehat{\mathrm{\Pi }}^{ab}(\alpha ^{},\beta )\widehat{\mathrm{\Pi }}^{ab}(\alpha ^{},\beta ^{})|2`$ (8)
where we call $`B(\alpha ,\beta )`$ as the Bell function.
By a simple extension of the relation (3), the two-mode Wigner function is found to be proportional to the mean of $`\widehat{\mathrm{\Pi }}_{ab}`$ such that $`W(\alpha ,\beta )=(4/\pi ^2)`$ Tr$`[\widehat{\rho }_{ab}\widehat{\mathrm{\Pi }}^{ab}(\alpha ,\beta )]`$ for the two-mode state of the density operator $`\widehat{\rho }_{ab}`$. The Bell function (7) can then be written in terms of the Wigner functions at different phase-space points,
$$B(\alpha ,\beta )=\frac{\pi ^2}{4}\left[W(0,0)+W(\alpha ,0)+W(0,\beta )W(\alpha ,\beta )\right].$$
(9)
The type of Bell’s inequality in Eq. (7) was first discussed by Clauser, Horne, Shimony, and Holt . Clauser and Horne later found another type of inequality which can be also expressed in phase space using the quasiprobability $`Q`$ function . The $`Q`$ function is related to the probability of the state having no photons. The lower and upper critical values of the Clauser-Horne Bell’s inequality is -1 and 0.
We have seen that the two-mode Wigner function is useful to test quantum nonlocality of the given field so that, in the next section, we find the evolution of the Wigner function for the initial two-mode squeezed state coupled with the thermal environment.
## III Time evolution of two-mode squeezed states in thermal environment
The two-mode squeezed state is the correlated state of two field modes $`a`$ and $`b`$ that can be generated by a nonlinear $`\chi ^{(2)}`$ medium . The two-mode pure squeezed state is obtained by applying the unitary operator on the two-mode vacuum
$$|\mathrm{\Psi }_{ab}(\sigma )=\mathrm{exp}\left(\sigma \widehat{a}\widehat{b}+\sigma ^{}\widehat{b}^{}\widehat{a}^{}\right)|0_a,0_b$$
(10)
where $`\sigma =s\mathrm{exp}(i\phi )`$ and $`\widehat{a}`$ ($`\widehat{b}`$) is an annihilation operator for the mode $`a`$ ($`b`$). The value of $`s`$ determines the degree of squeezing. The larger $`s`$ is, the more the state is squeezed.
The Wigner function corresponding to the squeezed state is the Fourier transform of its characteristic function $`C_W(\zeta ,\eta )`$ ,
$$C_W(\zeta ,\eta )=\mathrm{Tr}\left\{\widehat{\rho }\mathrm{exp}(\zeta \widehat{a}^{}\zeta ^{}\widehat{a})\mathrm{exp}(\eta \widehat{b}^{}\eta ^{}\widehat{b})\right\}.$$
(11)
For the two-mode squeezed state of the density matrix $`\widehat{\rho }=|\mathrm{\Psi }_{ab}(\sigma )\mathrm{\Psi }_{ab}(\sigma )|`$, the Wigner function is written as
$`W_{ab}(\alpha ,\beta )={\displaystyle \frac{4}{\pi ^2}}\mathrm{exp}[2\mathrm{cosh}(2s)(|\alpha |^2+|\beta |^2)`$ (12)
$`+2\mathrm{sinh}(2s)(\alpha \beta +\alpha ^{}\beta ^{})].`$ (13)
The correlated nature of the two-mode squeezed state is exhibited by the $`\alpha \beta `$ cross-term which vanishes when $`s=0`$.
The Fokker-Planck equation (in Born-Markov approximation) describing the time evolution of the Wigner function in the interaction picture can be written as
$`{\displaystyle \frac{W_{ab}(\alpha ,\beta ,\tau )}{\tau }}={\displaystyle \frac{\gamma }{2}}{\displaystyle \underset{\alpha _i=\alpha ,\beta }{}}[{\displaystyle \frac{}{\alpha _i}}\alpha _i+{\displaystyle \frac{}{\alpha _i^{}}}\alpha _i^{}`$ (14)
$`+2({\displaystyle \frac{1}{2}}+\overline{n}){\displaystyle \frac{^2}{\alpha _i\alpha _i^{}}}]W_{ab}(\alpha ,\beta ,\tau ),`$ (15)
where we have assumed that the two modes of the environment are independent each other and the energy decay rates of the two modes are same and denoted by $`\gamma `$. The two modes have the same average thermal photon number $`\overline{n}`$. By solving the Fokker-Planck equation (14), we get the time evolution of the Wigner function at time $`\tau `$ to be given by the convolution of the original function and the thermal environment ,
$`W_{ab}(\alpha ,\beta ,\tau )={\displaystyle \frac{1}{t(\tau )^4}}{\displaystyle d^2\zeta d^2\eta W_a^{th}(\zeta )W_b^{th}(\eta )}`$ (16)
$`\times W_{ab}\left({\displaystyle \frac{\alpha r(\tau )\zeta }{t(\tau )}},{\displaystyle \frac{\beta r(\tau )\eta }{t(\tau )}},\tau =0\right),`$ (17)
where the parameters $`r(\tau )=\sqrt{1e^{\gamma \tau }}`$ and $`t(\tau )=\sqrt{e^{\gamma \tau }}`$. $`W^{th}(\zeta )`$ is the Wigner function for the thermal state of the average thermal photon number $`\overline{n}`$:
$$W^{th}(\zeta )=\frac{2}{\pi (1+2\overline{n})}\mathrm{exp}\left(\frac{2|\zeta |^2}{1+2\overline{n}}\right).$$
(18)
Performing the integration in Eq. (16), the Wigner function for the initial two-mode squeezed state evolving in the thermal environment is obtained as
$`W_{ab}(\alpha ,\beta ,\tau )=𝒩\mathrm{exp}[E(\tau )(|\alpha |^2+|\beta |^2)`$ (19)
$`+F(\tau )(\alpha \beta +\alpha ^{}\beta ^{})]`$ (20)
where
$`E(\tau )`$ $`=`$ $`{\displaystyle \frac{2r(\tau )^2(1+2\overline{n})+2t(\tau )^2\mathrm{cosh}2s}{D(\tau )}}`$ (21)
$`F(\tau )`$ $`=`$ $`{\displaystyle \frac{2t(\tau )^2\mathrm{sinh}2s}{D(\tau )}}`$ (22)
$`D(\tau )`$ $`=`$ $`t(\tau )^4+2r(\tau )^2t(\tau )^2(1+2\overline{n})\mathrm{cosh}2s`$ (24)
$`+r(\tau )^4(1+2\overline{n})^2`$
and $`𝒩`$ is the normalization factor. In the limit of $`s=0`$, the $`\alpha \beta `$-cross term vanishes and the state can be represented by the direct product of each mode states such that $`W_{ab}(\alpha ,\beta ,\tau )=W_a(\alpha ,\tau )W_b(\beta ,\tau )`$. It is obvious that the Wigner function (19) exhibits the local characteristics in this limit.
The system will eventually assimilate with the environment which can be seen in the Wigner function, at the limit of $`\tau \mathrm{}`$,
$$W_{ab}(\alpha ,\beta )=\frac{4}{\pi ^2(1+\overline{n})^2}\mathrm{exp}[\frac{2}{(1+2\overline{n})}(|\alpha |^2+|\beta |^2)].$$
(25)
This is the direct product of two thermal states in modes $`a`$ and $`b`$.
## IV Evolution of quantum nonlocality
Substituting Eq. (19) into Eq. (9), we find the evolution of the nonlocality for the initial two-mode squeezed state in the thermal environment. The Bell function $`B`$ at time $`\tau `$ is written by
$`B(\alpha ,\beta ,\tau )={\displaystyle \frac{\pi ^2𝒩}{4}}\mathrm{exp}\{1+\mathrm{exp}[E(\tau )|\alpha |^2]`$ (26)
$`+\mathrm{exp}[E(\tau )|\beta |^2]\mathrm{exp}[E(\tau )(|\alpha |^2+|\beta |^2)`$ (27)
$`+2F(\tau )|\alpha \beta |\mathrm{cos}\theta ]\},`$ (28)
where $`\theta _\alpha `$ and $`\theta _\beta `$ are the phases of $`\alpha `$ and $`\beta `$ and $`\theta =\theta _\alpha +\theta _\beta `$. When $`\mathrm{cos}\theta =1`$, the Bell function $`B_m(|\alpha |,|\beta |,\tau )`$ is described by the absolute values $`|\alpha |`$ and $`|\beta |`$. $`B_m`$ is symmetric in exchanging $`\alpha `$ and $`\beta `$ such that $`B_m(|\alpha |,|\beta |,\tau )=B_m(|\beta |,|\alpha |,\tau )`$. It is straightforward to show that $`BB_m`$ at any instance of time $`\tau `$. In order to find the evolution of the nonlocality, the maximal value $`|B|_{max}`$ of the Bell function $`B`$ is calculated by the steepest decent method and using the properties of $`B_m(|\alpha |,|\beta |,\tau )`$. We say the field is quantum-mechanically nonlocal as $`|B|_{max}`$ is larger than 2 and the nonlocality is stronger as $`|B|_{max}`$ gets larger.
The initial two-mode squeezed state is always nonlocal as $`|B|_{max}>2`$ for $`s>0`$. $`|B|_{max}`$ increases monotonously as the degree $`s`$ of squeezing increases. The state becomes maximally nonlocal with $`|B|_{max}2.19055`$ as $`s\mathrm{}`$. In an intermediate time $`0<\tau <\mathrm{}`$, the pure squeezed state evolves to a two-mode mixed squeezed state and nonlocality is lost at a certain evolution time. Figs. 1 and 2 show $`|B|_{max}`$ versus the dimensionless time $`r(\tau )`$ defined in Eq. (16). We find that the nonlocality initially prepared persists until the characteristic time $`\tau _c(s,\overline{n})`$ depending on the temperature of the thermal environment and the initial squeezing. In Fig. 1 it is found that, when the environment is the vacuum, $`|B|_{max}`$ decreases as time proceeds. After reaching at the minimum value, $`|B|_{max}`$ increases to 2 which is the value of $`|B|_{max}`$ for the vacuum. Even though it is not clearly seen in the figure due to the scale of the figure, for any $`\overline{n}0`$ thermal environment, $`|B|_{max}`$ increases to its value for the thermal field after it decreases to a minimum. In Fig. 1, as $`\overline{n}`$ gets larger $`|B|_{max}`$ decreases much faster and further.
In Figs. 2, we identify an interesting phenomenon that the larger the initial degree of squeezing the more rapidly $`|B|_{max}`$ decreases. We analyze the reason why $`|B|_{max}`$ decreases more rapidly as the initial squeezing is larger as follows.
The two-mode squeezed state (10) can be represented by the continuous superposition of two-mode coherent states (A similar analysis was done for a single-mode squeezed state )
$$|\mathrm{\Psi }_{ab}(\sigma )=d^2\alpha G(\alpha ,\sigma )|\alpha ,\alpha ^{}\text{e}^{i\phi }$$
(29)
where the Gaussian weight function
$$G(\alpha ,\sigma )=(\pi \mathrm{sinh}s)^1\mathrm{exp}\left[\left(\frac{1\mathrm{tanh}s}{\mathrm{tanh}s}\right)|\alpha |^2\right].$$
(30)
As $`s`$ gets larger, the weight of a large $`\alpha `$ state is greater so that the contribution of $`|\alpha ,\alpha ^{}\text{e}^{i\phi }`$ of a large $`\alpha `$ becomes more important in the continuous superposition (29).
The quantum interference between coherent component states is the key of quantum nature of the field. The quantum interference is destroyed by the environment. The speed of destruction depends on the distance between the coherent component states and the average thermal energy of the environment . This is a reason why the macroscopic quantum superposition state is not easily seen in the nature. In the continuous superposition (29) we find that as the degree of squeezing is larger, the superposition extends further so that the quantum interference can be destroyed more easily. The quantum nonlocality in the two-mode squeezed state is also originated from the quantum interference between the coherent component states which can be destroyed easily as the contribution of the large amplitude coherent state becomes important.
In fact the uncertainty increases to its maximum and decreases to the value of the environment when a single-mode squeezed state is influenced by the thermal environment . The uncertainty increases faster as the degree of squeezing is larger. This can be explained using the same argument as the lost of quantum nonlocality.
In Fig. 2(a), when the environment is in the vacuum, it is found that the characteristic time $`\tau _c(s,\overline{n})`$ to lose the quantum nonlocality is shorter as the initial degree of squeezing is larger. In Fig. 2(b), when the non-zero temperature thermal environment ($`\overline{n}0`$) is concerned, we find that the larger degree of squeezing does not necessarily result in the shorter characteristic time $`\tau _c(s,\overline{n})`$. This clearly shows that the characteristic time is a function of the average number of thermal photons as well as the degree of squeezing. However, it is still true that $`|B|_{max}`$ decreases faster (the slope of its curve is steeper) when $`s`$ is larger. It is also found that $`|B|_{max}`$ decreases faster for $`\overline{n}0`$ than for $`\overline{n}=0`$.
We have studied the dynamic behavior of the nonlocality for the two-mode squeezed state in the thermal environment. The two-mode squeezed state can be used for the quantum channel in quantum teleportation of a continuous variable state. The two-mode squeezed state is found to be a nonlocal state regardless of its degree of squeezing and the higher degree of squeezing brings about the larger quantum nonlocality. As the squeezed state is influenced by the thermal environment the nonlocality is lost. The rapidity of the loss of nonlocality depends on the initial degree of squeezing and the average thermal energy of the environment. The more strongly the initial field is squeezed, the more rapidly the maximum nonlocality decreases. This has been analyzed extensively.
###### Acknowledgements.
We thank Professor J.W. Noh for bringing Ref. to our attention. This work is supported by the BK21 Grant of the Korea Ministry of Education and by the Sogang University Research Grants in 1999.
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# QED correction to radiative tail from elastic peak in DIS
## I Introduction
Deep inelastic scattering (DIS) is one of the powerful tools in investigating a nucleon nature. Recent experiments at CERN, DESY and SLAC have provided very precise data in the wide kinematical region. The modern level of data analysis in experiments on DIS requires careful consideration of the QED radiative effects which can give substantial contribution to measured quantities. Usually the radiative photon cannot be registered in a detector. As is well understood the corrections due to soft photons and loop effects cannot be separated from observables in principle. Hence their contribution has to be calculated theoretically and subtracted from observed data. The lowest order radiative corrections were first calculated by Mo and Tsai . Covariant approach utilized in was applied to unpolarized and polarized DIS. Second order correction to DIS cross section in the leading approximation is discussed in and recently in . For completeness we cite the papers in which the correction was calculated within a framework of electroweak theory basically for HERA kinematics. The excellent review of the lowest order RC can be found, say, in .
One of the important contribution to RC comes from a so called radiative tail from elastic peak (or simply elastic radiative tail), when final hadronic state is pure nucleon but the invariant mass of unobserved particle system (a radiated photon and this nucleon) is the same as in the main Born-level process. Consequently events detected can correspond to the main process as well as to the elastic radiative tail. The calculation in the lowest order of QED is standard and the results are included in many FORTRAN codes intending to perform the RC procedure of experimental data (see review and more recently developed codes ). Several papers were devoted to electroweak correction to elastic radiative tail. Numerical analysis of the elastic radiative tail shows that its contribution is very important and can exceed the main measured process at the Born level. Therefore the next step is to calculate QED correction to the elastic radiative tail with the maximal possible accuracy. So far only the leading correction to elastic radiative tail due to double bremsstrahlung, which is part of the total second order correction, was treated numerically and the attempt to calculate the correction exactly was done in .
The structure of the cross section of elastic radiative tail is the following
$$\sigma ^{ERT}\underset{Q_{hmin}^2}{}\text{d}Q_h^2𝒦(Q_h^2,Q^2,W^2)^2(Q_h^2),$$
(1)
where $`Q_h^2`$ is a momentum square transferred to hadronic system, and $`Q^2`$ and $`W^2`$ are leptonic kinematical variables measured. The quantity $`𝒦`$ is a kinematical factor known exactly and $``$ is a nuclear form factor. Due to rapid fall of the form factor squared as a function of $`Q_h^2`$ the main integration region is close to the lower integration limit. In papers this fact was used to construct an approximation where $`Q_h^2`$ is considered as a small parameter of order of the proton mass squared. In this paper we will also use this approximation to analyze the correction to elastic radiative tail. Application of Sudakov technique will allow us to obtain compact explicit formulae for processes considered. The first effect which has to be considered is the one-loop correction and the emission of additional real photon. We will analyze both of them at leading and next-to-leading levels. Another effect contributing to the cross section is a lepton pair creation. We will calculate it in the leading log approximation.
Obtaining a second order correction to deep inelastic scattering is the main motivation of this paper. However our results can be used in other cases. For instance, they can be considered as a radiative correction in measurements with hard photon detected in coincidence with scattered lepton (see , for example), that allows one to reach kinematical regions otherwise unreachable. That is why we do not concretize our notation usually used in DIS but instead try to keep it as general as possible. In the next section we introduce our notation and obtain results for the cross section of single bremsstrahlung using Sudakov technique. In section III we give results for one-loop corrections. Double bremsstrahlung and contributions due to pair production are considered in sections IV and V and final remarks are given in section VI. Some technical details are discussed in Appendices.
## II Single bremsstrahlung
We study the process
$`e(p_1)+p(p_2)e(p_1^{})+\gamma (k_1)+p(p_2^{}),s=2p_1p_2,`$ (2)
$`Q_h^2=(p_2p_2^{})^2,Q^2=2p_1p_1^{},k_1^2=0,`$ (3)
$`p_1^2=p_1^{}_{}{}^{}2=m^2,p_2^2=p_2^2^{}=M^2,q^2=O_h^2,`$ (4)
in the kinematical region
$$sQ^2>Q_h^2M^2,2p_2p_1^{}s.$$
(5)
The expression for differential cross section in Born approximation looks (details are given in the Appendix A):
$`2\epsilon _1^{}{\displaystyle \frac{\text{d}^3\sigma _0^\gamma }{\text{d}^3p_1^{}}}={\displaystyle \frac{4\alpha ^3}{\pi ^2}}{\displaystyle \frac{\text{d}^2𝐪}{(𝐪^2+Q_{\mathrm{min}}^2)^2}\frac{1}{1b}\mathrm{\Phi }^\gamma \mathrm{\Phi }^{\mathrm{prot}}},`$ (6)
with $`b=2p_2p_1^{}/s`$ the energy fraction of the scattered electron. We imply the Sudakov parameterization of the 4-momenta in the problem (see Appendix A).
Note that due to gauge invariance condition
$$q^\rho J_\rho ^{(1)}(\alpha _qp_2+q_{})^\rho J_\rho ^{(1)}=0,$$
(7)
the quantity $`\mathrm{\Phi }^\gamma `$ is constructed out of $`(1/s)p_2J^{(1)}`$ which may be rearranged as follows:
$`{\displaystyle \frac{1}{s}}p_2^\mu J_\mu ^{(1)}`$ $`=`$ $`{\displaystyle \frac{s}{s_1}}|𝐪|e_q^\mu J_\mu ^{(1)},e_q={\displaystyle \frac{𝐪}{|𝐪|}},`$ (8)
$`s_1`$ $`=`$ $`s\alpha _q=(p_1^{}+k_1)^2+𝐪^2m^2.`$ (9)
Thus $`\mathrm{\Phi }^\gamma `$ vanishes for small $`𝐪^2`$. The explicit expression for $`\mathrm{\Phi }^{\mathrm{prot}}`$ is found to be
$$\mathrm{\Phi }^{\mathrm{prot}}=2(F_1^2+\frac{𝐪^2}{M^2}F_2^2).$$
(10)
For $`\mathrm{\Phi }^\gamma `$ we have (we refer for further details to the Appendix A):
$`\mathrm{\Phi }^\gamma ={\displaystyle \frac{(1b)^2b(1+b^2)𝐪^2}{n_1n}},`$ (11)
with
$`n=(𝐩_1^{}b𝐪)^2,n_1=(𝐩_1^{}𝐪)^2.`$ (13)
Another fact is that both $`\mathrm{\Phi }^\gamma /𝐪^2`$ and $`\mathrm{\Phi }^{\mathrm{prot}}`$ do not vanish in the limit of small momentum transfer $`|𝐪|`$, thus providing the logarithmic enhancement upon performing the $`Q_h^2𝐪^2`$ integration (Weizsäcker-Williams approximation). Indeed, the quantity $`Q_{\mathrm{min}}^2`$ entering the cross section is a small quantity,
$$Q_{\mathrm{min}}^2=M^2\left(\frac{Q^2}{(1b)s}\right)^2M^2.$$
(14)
For completeness we put the phase volume of the scattered electron in terms of Sudakov variables:
$`{\displaystyle \frac{\text{d}^3p_1^{}}{2\epsilon _1^{}}}={\displaystyle \frac{\text{d}b}{2b}}\text{d}^2𝐩_1^{},Q^2=2p_1p_1^{}={\displaystyle \frac{𝐩_1^{}_{}{}^{}2}{b}}.`$ (15)
Note that the requirement $`Q^2>Q_h^2`$ provides the absence of singularities while doing an integration over $`\text{d}^2𝐪`$.
## III Virtual and soft photons emission contribution
The correction coming from the emission of virtual and soft photons (in the cms reference frame) can be drawn out of paper , in which the radiative corrections to the Compton tensor were calculated
$`2\epsilon _1^{}{\displaystyle \frac{\text{d}^3\sigma _{B+V+S}}{\text{d}^3p_1^{}}}`$ $`=`$ $`2\epsilon _1^{}{\displaystyle \frac{\text{d}^3\sigma _0^\gamma }{\text{d}^3p_1^{}}}[1+{\displaystyle \frac{\alpha }{2\pi }}\stackrel{~}{\rho }`$ (16)
$`+`$ $`{\displaystyle \frac{\alpha }{4\pi }}{\displaystyle \frac{1}{1+b^2}}(\tau _{11}+b(\tau _{12}+\stackrel{~}{\tau }_{12})+b^2\stackrel{~}{\tau }_{11})],`$ (17)
with
$`\stackrel{~}{\rho }`$ $`=`$ $`2(L1)(2\mathrm{ln}\mathrm{\Delta }\mathrm{ln}b)+3L_h\mathrm{ln}^2b`$ (18)
$``$ $`{\displaystyle \frac{9}{2}}{\displaystyle \frac{\pi ^2}{3}}+2\text{Li}_2\left(\mathrm{cos}^2{\displaystyle \frac{\theta }{2}}\right),`$ (19)
$`L`$ $`=`$ $`\mathrm{ln}{\displaystyle \frac{Q^2}{m^2}},L_h=\mathrm{ln}{\displaystyle \frac{Q_h^2}{m^2}},`$ (20)
$`\mathrm{\Delta }`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }E}{E}},\text{Li}_2(x)={\displaystyle \underset{0}{\overset{x}{}}}{\displaystyle \frac{\mathrm{ln}(1y)}{y}}\text{d}y.`$ (21)
The Born cross section after substitution of Eq. (11) into Eq. (6) and neglect of sub-leading terms becomes
$`2\epsilon _1^{}{\displaystyle \frac{\text{d}^3\sigma _0^\gamma }{\text{d}^3p_1^{}}}={\displaystyle \frac{4\alpha ^3}{\pi ^2}}{\displaystyle \frac{\text{d}^2\mathrm{𝐪𝐪}^2}{(𝐪^2+Q_{\mathrm{min}}^2)^2}\frac{(1b)(1+b^2)}{b(Q^2)^2}\mathrm{\Phi }^{\mathrm{prot}}},`$
where $`\mathrm{\Delta }E,E`$ are the upper bound on the undetectable soft photon energy, and the energy of the initial electron, respectively; $`\theta `$ is the angle in the laboratory reference frame between the initial and the scattered electron momenta. Somewhat cumbersome functions $`\tau _{ij}`$ are explicitly given in the Appendix D. It should be noted that they do not contain any large logarithms but include the quantity $`Q_h^2`$ which is small in our approximation. If one keeps only non-zero terms in the expansion over $`Q_h^2`$ then
$`{\displaystyle \frac{1}{2}}\left(\tau _{11}+b(\tau _{12}+\stackrel{~}{\tau }_{12})+b^2\stackrel{~}{\tau }_{11}\right)=\left[3\mathrm{log}{\displaystyle \frac{Q^2}{Q_h^2(1b)}}1\right]`$ (22)
$`\times (1+b^2)+4b\mathrm{log}(1b)+[b^2+(1b)^2]`$ (23)
$`\times \left[\mathrm{log}^2{\displaystyle \frac{(1b)}{b}}+\pi ^2\right]+[1+(1b)^2]\mathrm{log}^2(1b)`$ (24)
$`+(32b)\mathrm{log}b,`$ (25)
The logarithms $`\mathrm{log}Q_h^2`$ cancel out exactly in the sum of (22) and $`\stackrel{~}{\rho }`$.
## IV Two hard photons emission contribution
We will consider now the process of two hard photons emission:
$$e(p_1)+p(p_2)e(p_1^{})+\gamma (k_1)+\gamma (k_2)+p(p_2^{}).$$
(26)
The relevant contribution to the cross section looks
$`2\epsilon _1^{}{\displaystyle \frac{\text{d}^3\sigma }{\text{d}^3p_1^{}}}`$ $`=`$ $`{\displaystyle \frac{\alpha ^4}{8\pi ^4}}{\displaystyle \frac{\text{d}^2𝐪}{(𝐪^2+Q_{\mathrm{min}}^2)^2}\frac{\text{d}x_1\text{d}^2𝐤_1}{x_1x_2}\mathrm{\Phi }^{\gamma \gamma }\mathrm{\Phi }^{\mathrm{prot}}},`$ (27)
$`Q_{\mathrm{min}}^2`$ $`=`$ $`M^2\left({\displaystyle \frac{s_1}{s}}\right)^2,s_1=(p_1^{}+k_1+k_2)^2,`$ (28)
with the expression for $`\mathrm{\Phi }^{\mathrm{prot}}`$ given earlier. The explicit form of $`\mathrm{\Phi }^{\gamma \gamma }`$ can be found in the Appendix B. The integration over $`\text{d}^2𝐤_1`$ may be performed using the integrals given in the Appendix C.
Concerning the region $`Q_h^2Q^2`$, the result is found to be <sup>*</sup><sup>*</sup>*Upon applying the crossing transformation to the amplitude of $`e\overline{e}`$ annihilation to $`\gamma \gamma \gamma `$ presented in paper .
$`\mathrm{\Phi }^{\gamma \gamma }`$ $`=`$ $`16𝐪^2\{{\displaystyle \frac{Q^2}{s_1^2}}[{\displaystyle \frac{s_1^2(1+b^2)}{d_1d_2d_1^{}d_2^{}}}+{\displaystyle \frac{d_1^2+d_1^{}_{}{}^{}2}{bs_1^2d_2d_2^{}}}+{\displaystyle \frac{d_2^2+d_2^{}_{}{}^{}2}{bs_1^2d_1d_1^{}}}]`$ (29)
$``$ $`{\displaystyle \frac{2}{Q^4}}(1+𝒫_{12})[{\displaystyle \frac{m^2}{d_1^2}}{\displaystyle \frac{x_2^2(b^2+(1x_1)^2)}{b(1x_1)^3}}`$ (30)
$`+`$ $`{\displaystyle \frac{m^2}{d_1^{}_{}{}^{}2}}{\displaystyle \frac{x_2^2b^2(1+(1x_2)^2)}{(1x_2)^3}}]\}.`$ (31)
with the notations introduced
$`s_1={\displaystyle \frac{𝐤_1^2}{x_1}}+{\displaystyle \frac{𝐤_2^2}{x_2}}+{\displaystyle \frac{𝐩_1^{}_{}{}^{}2}{b}},d_i={\displaystyle \frac{1}{x_i}}(m^2x_i^2+𝐤_i^2),`$ (32)
$`d_i^{}={\displaystyle \frac{1}{x_ib}}[m^2x_i^2+(x_i𝐩_1^{}b𝐤_i)^2],`$ (33)
where $`x_{1,2}`$ are the energy fractions of hard photons, $`x_1+x_2+b=1`$. Besides we use the relations
$$𝐤_1+𝐤_2+𝐩_1^{}=0,2qp_1^{}=s_1b,s_1=2qp_1=s\alpha _q.$$
An integration over $`\text{d}^2𝐤_1`$ may be performed analytically and to a logarithmic accuracy it boils down to
$$\frac{\text{d}^2𝐤_1}{\pi }[\frac{1}{d_1};\frac{1}{d_2};\frac{1}{d_1^{}};\frac{1}{d_2^{}}]=L[x_1;x_2;\frac{x_1}{b};\frac{x_2}{b}].$$
(34)
The resulting contribution (again to a logarithmic accuracy) takes the following form
$`{\displaystyle \text{d}^2𝐤_1\mathrm{\Phi }^{\gamma \gamma }}`$ $`=`$ $`{\displaystyle \frac{16\pi 𝐪^2L}{b(Q^2)^2}}(1+𝒫_{12})x_2^2`$ (35)
$`\times `$ $`[(1+{\displaystyle \frac{1}{(1x_1)^2}}+{\displaystyle \frac{b^2}{(1x_2)^2}})(1+b^2)`$ (36)
$`+`$ $`{\displaystyle \frac{b^2}{(1x_1)^4}}+{\displaystyle \frac{b^4}{(1x_2)^4}}],`$ (37)
$`\mathrm{\Delta }`$ $`<`$ $`x_i<1b\mathrm{\Delta }.`$ (38)
Carrying out the integration of Eq. (29) over $`\stackrel{}{k}_1`$ and $`x_1`$ to a next-to-leading accuracy we obtain for the contribution to the differential cross section coming from emission of two hard photons,
$$2\epsilon _1^{}\frac{\text{d}^3\sigma ^{\gamma \gamma }}{\text{d}^3p_1^{}}=\frac{2\alpha ^4}{\pi ^3}\frac{\text{d}^2𝐪}{(𝐪^2+Q_{\mathrm{min}}^2)^2}\frac{𝐪^2(T_{LL}+T_{NLO})}{b(Q^2)^2}\mathrm{\Phi }^{\mathrm{prot}}$$
(39)
where the leading and next-to-leading contributions read
$`T_{LL}`$ $`=`$ $`(L1)[4(1b)(1+b^2)\mathrm{ln}{\displaystyle \frac{1b}{\mathrm{\Delta }}}`$ (40)
$`+`$ $`(1b)(1b^2)\mathrm{ln}b{\displaystyle \frac{2}{3}}(1b)(72b+7b^2)],`$ (41)
$`T_{NLO}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{b^4+6b^2+1}{1+b}}\mathrm{log}^2b`$ (46)
$`{\displaystyle \frac{1}{3}}(3b^2)(3b)\mathrm{log}b`$
$`+{\displaystyle \frac{8}{3}}(1b)(b^2+b+1)\mathrm{log}(1b)`$
$`(1b)[{\displaystyle \frac{1}{3}}(15b^22b+15)`$
$`+2(\text{Li}_2(b){\displaystyle \frac{\pi ^2}{6}}){\displaystyle \frac{b^4+6b^2+1}{1b^2}}].`$
There are two possible experimental setups we concern with: the first one in which a recoil proton is registered, and the second — pure inclusive setup — with only a final lepton observed. Definitely, NLO contribution obtained can be counted valid only for the former experimental setup, while in the latter case one can use the expression given above only to a LL accuracy.
The general answer for the cross section in Born approximation with the lowest order correction to the leading approximation is a sum of the contributions coming from virtual and real soft photons emission given above as well as from two hard photons emission and is free from dependence on the auxiliary parameter $`\mathrm{\Delta }`$.
The graphs given below illustrate behavior featured by the complete QED RC contribution to the cross section of DIS as well as the comparative contributions of the LL, NLO terms and of the correction due to pair production.
## V Contribution of lepton pair production
Consider now the hard pair production process that takes place at the same order of perturbation theory as the two hard photon emission. In the same way we may conclude that the soft pair case as well as the case of double collinear kinematics does not contribute to the radiative tail. Therefore we may consider only semi-collinear kinematics of hard pair production of which there exist two different mechanisms . One of these is the two photon mechanism of pair creation. An electron from that pair having momentum $`p_1^{}`$ is detected in experiment and the scattered electron moves close to the initial electron direction. This kinematics permits us to apply the Weizsäcker-Williams approximation,
$`2\epsilon _1^{}{\displaystyle \frac{\text{d}^3\sigma _{\mathrm{pair}}^{(1)}}{\text{d}^3p_1^{}}}`$ $`=`$ $`{\displaystyle \frac{2\alpha ^4}{\pi ^3}}{\displaystyle \frac{\text{d}^2𝐪}{(𝐪^2+Q_{\mathrm{min}}^2)^2}\frac{𝐪^2L}{b(Q^2)^2}\mathrm{\Phi }^{\mathrm{prot}}}`$ (47)
$`\times `$ $`{\displaystyle \frac{\text{d}\beta _{}}{(1\beta _{})^4}}((1\beta _{}b)^2+b^2)(1+\beta _{}^2),`$ (48)
$`s_1`$ $`=`$ $`Q^2{\displaystyle \frac{1\beta _{}}{\beta _+}},b+\beta _{}+\beta _+=1.`$ (49)
The second mechanism is characterized by the bremsstrahlung mechanism of pair creation, with an electron from a pair to be detected. Leaving details to the Appendix E let us present here the result
$`2\epsilon _1^{}{\displaystyle \frac{\text{d}^3\sigma _{\mathrm{pair}}^{(2)}}{\text{d}^3p_1^{}}}`$ $`=`$ $`{\displaystyle \frac{2\alpha ^4}{\pi ^3}}{\displaystyle \frac{\text{d}^2𝐪}{(𝐪^2+Q_{\mathrm{min}}^2)^2}\frac{𝐪^2L}{(Q^2)^2}\mathrm{\Phi }^{\mathrm{prot}}}`$ (50)
$`\times `$ $`{\displaystyle \frac{b(1+\beta _{}^2)\text{d}\beta _{}}{(1\beta _{})^4}}[(1b\beta _{})^2+b^2]`$ (51)
$`s_1`$ $`=`$ $`Q^2{\displaystyle \frac{1\beta _{}}{b\beta _{}}}.`$ (52)
The integration over $`\beta _{}`$ can be performed analytically with additional assumption that $`Q_{\mathrm{min}}^2`$ has no $`\beta _{}`$ dependence. The result for the sum of these contributions is found to be
$`2\epsilon _1^{}{\displaystyle \frac{\text{d}^3\sigma _{\mathrm{pair}}}{\text{d}^3p_1^{}}}`$ $`=`$ $`{\displaystyle \frac{2\alpha ^4}{\pi ^3}}{\displaystyle \frac{\text{d}^2𝐪}{(𝐪^2+Q_{\mathrm{min}}^2)^2}\frac{𝐪^2L(1+b^2)}{b(Q^2)^2}\mathrm{\Phi }^{\mathrm{prot}}}`$ (53)
$`\times `$ $`\left(1b+2(1+b)\mathrm{log}b+{\displaystyle \frac{4}{3b}}(1b^3)\right).`$ (54)
## VI Discussion and conclusion
In the paper presented the correction to radiative tail from elastic peak is studied in the kinematics when a final lepton is measured. Using Sudakov technique the contributions of loops (16), double photon bremsstrahlung (39) and a pair production (47,50) are calculated.
In this section we analyze obtained contributions numerically. Both the relative contributions of the processes considered and the total correction to the lowest order process are investigated within kinematical conditions of experiments on electron DIS at TJNAF and DESY (both for HERA and for HERMES). It is convenient to define the following quantities:
$$\delta =\frac{\sigma _L+\sigma _N+\sigma _p}{\sigma _0},\delta _{L,N,p}=\frac{\sigma _{L,N,p}}{\sigma _0}.$$
(55)
Here $`\sigma _0`$ stands for the cross section of radiative tail from elastic peak (6). Other $`\sigma `$’s constitute the next order results. The quantity $`\sigma _p`$ is a direct sum of two mechanisms of pair creations (47,50), whereas $`\sigma _L`$ and $`\sigma _N`$ are the leading (including mass singularities terms $`\mathrm{log}(Q^2/m^2)`$) and next-to-leading (independent of leptonic mass) terms.
They are obtained upon summing up expressions given in Eqs. (16) and (39) after cancellation of infrared divergence.
The quantities $`\delta `$ and $`\delta _{L,N,p}`$ are presented in Figs. 1 and 2. As is clearly seen the corrections have a steep dependence on $`y`$ and a poor one on $`x`$ (or $`Q^2`$) and $`s`$. Figure 2 shows results for $`\sqrt{s}=7GeV`$, however very similar plots can be produced for the two other cases considered. Having analyzed results obtained the following conclusions could be drawn.
The relative radiative correction to elastic radiative tail is important practically everywhere. The modern level of data analysis and very high experimental accuracies achieved in current experiments on DIS require that a generalization of standard radiative correction procedure be made in order to include a second order radiative correction. An extremely interesting region where the correction considered is important is actually high $`y`$ domain. Remind, that this one (up to $`y`$0.95) is under investigation at TJNAF.
The main contribution to a second order radiative correction comes from the effect of pair creation. Asymptotical behavior of $`\sigma _p`$ for small $`b=1y`$ is $`1/b^2`$ whereas the other cross sections feature only $`1/b`$ behavior. That is in fact a reason of the large correction in the region of high $`y`$. In the paper presented this particular contribution is calculated in the leading approximation only, therefore a study of the correction, induced by a pair production, at the next-to-leading level is highly desirable.
The relative contribution of the next-to-leading correction $`\sigma _N`$ is not small with respect to the leading log contribution $`\sigma _L`$. In the region of large $`y`$ the relative contribution $`\sigma _N/\sigma _L`$ does not exceed 5%, whereas for small $`y`$ it can reach as much as 20–30%. From the other hand the next-to-leading contribution completely fixes all uncertainties of leading log approximation thus leaving unknown only terms proportional to lepton mass squared and $`Q_h^2`$, which is effectively small due to behavior of form factors.
## Acknowledgements
We would like to thank P.Kuzhir for useful discussions and comments. The work of IA was partially supported by the U.S. Department of Energy under contract DE–AC05–84ER40150. EAK and BGS acknowledge support of RFBR via grant No. 99-02-17730.
## Appendix A. Details of matrix element calculus: the case of single photon bremsstrahlung
Using the Sudakov decomposition of the 4-vectors in the problem
$`p_1^{}=\alpha _1^{}\stackrel{~}{p}_2+b\stackrel{~}{p}_1+p_1^{},k_1=\alpha _1\stackrel{~}{p}_2+x_1\stackrel{~}{p}_1+k_1,`$ (A.1)
$`q=p_2p_2^{}=\alpha _q\stackrel{~}{p}_2+\beta _q\stackrel{~}{p}_1+q_{},`$ (A.2)
$`p_2^{}=\alpha _2^{}\stackrel{~}{p}_2+\beta _2^{}\stackrel{~}{p}_1+p_2^{},v_{}p_1=v_{}p_2=0,`$ (A.3)
$`\stackrel{~}{p}_1=p_1p_2{\displaystyle \frac{m^2}{s}},\stackrel{~}{p}_2=p_2p_1{\displaystyle \frac{M^2}{s}},`$ (A.4)
we have excluded parameters $`\alpha _1,\alpha _1^{},\beta _q`$ using the on–shell conditions
$`p_2^{}_{}{}^{}2M^2=s\beta _q(1\alpha _q)𝐪^2\alpha _qM^2=0,`$ (A.5)
$`p_1^{}_{}{}^{}2=sb\alpha _1^{}𝐩_1^{}_{}{}^{}2=0,k_1^2=sx_1\alpha _1𝐤_1^2=0,`$ (A.6)
besides
$`\mathrm{\Phi }^{\mathrm{prot}}`$ $`=`$ $`{\displaystyle \frac{1}{s^2}}\mathrm{Sp}\{(\widehat{p}_2^{}+M)\mathrm{\Gamma }_\rho (\widehat{p}_2+M)\stackrel{~}{\mathrm{\Gamma }}_\sigma p_1^\rho p_1^\sigma \},`$ (A.7)
$`\mathrm{\Gamma }_\rho `$ $`=`$ $`F_1(q^2)\gamma _\rho +{\displaystyle \frac{\sigma _{\mu \rho }q^\mu }{2M}}F_2(q^2).`$ (A.8)
Here $`F_{1,2}(q^2)`$ are the Dirac and Pauli form factors of a proton. For $`\mathrm{\Phi }^\gamma `$ we have:
$`\mathrm{\Phi }^\gamma ={\displaystyle \frac{1}{s^2}}\mathrm{Sp}\{\widehat{p}_1^{}O_\mu \widehat{p}_1\stackrel{~}{O}^\mu \},`$ (A.9)
$`O_\mu =\widehat{p}_2{\displaystyle \frac{\widehat{p}_1\widehat{k}_1}{2p_1k_1}}\gamma _\mu +\gamma _\mu {\displaystyle \frac{\widehat{p}_1^{}+\widehat{k}_1}{2p_1^{}k_1}}\widehat{p}_2,`$ (A.10)
and then
$$q^2=Q_h^2=\frac{1}{1\alpha _q}[𝐪^2+M^2\alpha _q^2][𝐪^2+Q_{\mathrm{min}}^2],$$
(A.11)
with $`Q_{\mathrm{min}}^2`$ given in the text. The matrix element
$$M=\frac{1}{q^2}J_\sigma ^{(1)}\overline{u}(p_2^{})\mathrm{\Gamma }_\rho u(p_2)g^{\rho \sigma },$$
(A.12)
using the Gribov representation for the metric tensor
$$g^{\rho \sigma }=g_{}^{\rho \sigma }+\left(\frac{2}{s}\right)(\stackrel{~}{p}_2^\rho \stackrel{~}{p}_1^\sigma +\stackrel{~}{p}_2^\sigma \stackrel{~}{p}_1^\rho )\left(\frac{2}{s}\right)\stackrel{~}{p}_2^\sigma \stackrel{~}{p}_1^\rho ,$$
(A.13)
may be put in a form
$$M=\frac{2s}{q^2}\left(\frac{1}{s}p_2^\sigma J_\sigma ^{(1)}\right)\left(\frac{1}{s}\overline{u}(p_2^{})\mathrm{\Gamma }_\rho u(p_2)p_1^\rho \right).$$
(A.14)
Note that each expressions in the parentheses on the RHS of Eq. (A.14) do not depend on $`s`$ in the limit $`s\mathrm{}`$. The expression for $`\mathrm{\Phi }^\gamma `$ may be transformed using the following reduced expression
$`O_\mu =x_1[sb\gamma _\mu ({\displaystyle \frac{1}{n}}{\displaystyle \frac{1}{n_1}})`$ $`+`$ $`{\displaystyle \frac{1}{n_1}}b\gamma _\mu \widehat{q}\widehat{p}_2{\displaystyle \frac{1}{n}}\gamma _\mu \widehat{p}_2\widehat{q}],`$ (A.15)
$`x_1`$ $`=`$ $`1b.`$ (A.16)
to take the form given in Eq. (11).
## Appendix B. Details of matrix element calculus: the case of double photon bremsstrahlung
Let’s first demonstrate that the matrix element of the process
$$\gamma ^{}(q)+e(p_1)e(p_1^{})+\gamma (k_1)+\gamma (k_2)$$
(B.1)
is explicitly proportional to $`𝐪`$ for small values of the latter, which is in fact the requirement of gauge invariance with respect to the virtual photon. The matrix element is described by six diagrams. With regard to the gauge invariance this set can be separated out to the two subsets in each of which the gauge condition is satisfied independently. Introducing the photon-permutating operator $`𝒫_{12}`$ we bring the matrix element to the form:
$`=(1+𝒫_{12})Q,Q=_1+_2+_3,`$ (B.2)
where
$`_1`$ $`=`$ $`{\displaystyle \frac{1}{dd_1}}\overline{u}(p_1^{})\widehat{p}_2(\widehat{p}_1\widehat{k}_1\widehat{k}_2+m)`$ (B.4)
$`\times \widehat{e}_2^{}(\widehat{p}_1\widehat{k}_1+m)\widehat{e}_1^{}u(p_1),`$
$`_2`$ $`=`$ $`{\displaystyle \frac{1}{d_1d_2^{}}}\overline{u}(p_1^{})\widehat{e}_2^{}(\widehat{p}_1\widehat{k}_1+\widehat{q}+m)\widehat{p}_2`$ (B.5)
$`\times `$ $`(\widehat{p}_1\widehat{k}_1+m)\widehat{e}_1^{}u(p_1),`$ (B.6)
$`_3`$ $`=`$ $`{\displaystyle \frac{1}{d^{}d_2^{}}}\overline{u}(p_1^{})\widehat{e}_2^{}(\widehat{p}_1\widehat{k}_1+\widehat{q}+m)`$ (B.8)
$`\times \widehat{e}_1^{}(\widehat{p}_1+\widehat{q}+m)\widehat{p}_2u(p_1),`$
and
$`d`$ $`=`$ $`d_1+d_2{\displaystyle \frac{1}{x_1x_2}}(x_1\stackrel{}{k}_2x_2\stackrel{}{k}_1)^2,`$
$`d^{}`$ $`=`$ $`d_1^{}+d_2^{}+{\displaystyle \frac{1}{x_1x_2}}(x_1\stackrel{}{k}_2x_2\stackrel{}{k}_1)^2.`$
The permutation operator $`𝒫_{12}`$ for the photons acts the following way
$`𝒫_{12}f(k_1,e_1;k_2,e_2)=f(k_2,e_2;k_1,e_1),𝒫_{12}^2=1.`$
The quantity $`Q`$ is gauge invariant regarding the virtual photon $`k`$ since all permutations of this photon have been taken into account. Therefore $`Q`$ is proportional to $`q_{}`$ in the limit of $`q_{}0`$. Indeed, making use of the relations
$`Q=p_{2\mu }Q^\mu ,q_\mu Q^\mu =(\alpha _q\stackrel{~}{p}_2+q_{})_\mu Q^\mu =0,`$ (B.9)
we immediately obtain (neglecting the small contribution $`\beta _qp_\mu Q^\mu 1/s`$)
$`Q={\displaystyle \frac{q_\mu }{\alpha _q}}Q^\mu .`$ (B.10)
Then transform the quantities $`_j`$ to such a form that the noticed low $`q_{}`$ behavior is present in their sum $`Q`$ explicitly. The reason is that in this case all individual large (compared to $`q_{}`$) contributions are mutually cancelled. The first step is to use the Dirac equations $`\widehat{p}_1u(p_1)=mu_1`$, $`\overline{u}(p_1^{})\widehat{p}_1^{}=m\overline{u}(p_1^{})`$ and to rearrange the amplitudes $`_j`$ of Eq. (B.4),
$`_1`$ $`=`$ $`\overline{u}(p_1^{})\{{\displaystyle \frac{s\beta _1^{}}{d_1}}\widehat{e}_2^{}(\widehat{p}_1\widehat{k}_1+m)\widehat{e}_1^{}`$ (B.11)
$``$ $`{\displaystyle \frac{1}{d_1}}\widehat{p}_2\widehat{q}\widehat{e}_2^{}(\widehat{p}_1\widehat{k}_1+m)\widehat{e}_1^{}\}u(p_1),`$ (B.12)
$`_2`$ $`=`$ $`\overline{u}(p_1^{})\{+{\displaystyle \frac{s(1x_1)}{d_1d_2^{}}}\widehat{e}_2^{}(\widehat{p}_1\widehat{k}_1+m)\widehat{e}_1^{}{\displaystyle \frac{1}{d_2^{}}}\widehat{e}_2^{}\widehat{p}_2\widehat{e}_1^{}`$ (B.14)
$`+{\displaystyle \frac{1}{d_1d_2^{}}}\widehat{e}_2^{}\widehat{q}\widehat{p}_2(\widehat{p}_1\widehat{k}_1+m)\widehat{e}_1^{}\}u(p_1),`$
$`_3`$ $`=`$ $`\overline{u}(p_1^{})\{{\displaystyle \frac{s}{d^{}d_2^{}}}\widehat{e}_2^{}(\widehat{p}_1\widehat{k}_1+m)\widehat{e}_1^{}+{\displaystyle \frac{s}{d^{}d_2^{}}}\widehat{e}_2^{}\widehat{q}\widehat{e}_1^{}`$ (B.16)
$`+{\displaystyle \frac{1}{d^{}d_2^{}}}\widehat{e}_2^{}(\widehat{p}_1^{}+\widehat{k}_2+m)\widehat{e}_1^{}\widehat{q}\widehat{p}_2\}u(p_1).`$
From these formulae it can be noted that the last terms in $`_1,_2,_3`$, up to terms of the order of
$`{\displaystyle \frac{m^2}{E^2}},\theta ^2,{\displaystyle \frac{m}{E}}\theta ,`$
are proportional to $`q_{}`$,
$`\widehat{\stackrel{~}{p}}_2\widehat{q}=\widehat{\stackrel{~}{p}}_2(\alpha _q\widehat{\stackrel{~}{p}}_2+\beta _q\widehat{p}+\widehat{q}_{})=\widehat{\stackrel{~}{p}}_2\widehat{q}_{}=\widehat{q}\widehat{\stackrel{~}{p}}_2.`$ (B.17)
Next, one can see that the sum of the first three terms in Eqs. (B.11) is also proportional to $`q_{}`$ since (for more details see )
$`A{\displaystyle \frac{b}{d_1}}+{\displaystyle \frac{1x_1}{d_1d_2^{}}}+{\displaystyle \frac{1}{d^{}d_2^{}}},A|_{q_{}0}=0.`$ (B.18)
Finally we consider the sum of the second terms of the quantities $`_2,_3`$ given in Eqs. (B.11). Using the relations (, Eq.(21)) and
$`(p_1^{}+k_1+k_2)^2=(p_1k)^2=m^2𝐤^2s\alpha _k,`$
one immediately gets
$`{\displaystyle \frac{\widehat{\stackrel{~}{p}}_2}{d_2^{}}}+{\displaystyle \frac{s(\alpha _q\widehat{\stackrel{~}{p}}_2+\widehat{q}_{})}{d^{}d_2^{}}}={\displaystyle \frac{s\widehat{q}_{}}{d^{}d_2^{}}}+{\displaystyle \frac{\widehat{\stackrel{~}{p}}_2𝐪^2}{d^{}d_2^{}}}.`$ (B.19)
Therefore, from Eqs. (B.17), (B.18), (B.19) it is clearly seen that the property illustrated by Eq. (B.10)
$`\left(_1+_2+_3\right)|_{q_{}0}=0`$
is evidently satisfied and consequently the quantity $`Q=_{j=1}^3_j`$ became a sum of terms explicitly proportional to $`q_{}`$,
$`Q`$ $`=`$ $`\overline{u}(p_1^{})\{As\widehat{e}_2^{}(\widehat{p}_1\widehat{k}_1+m)\widehat{e}_1^{}`$ (B.20)
$``$ $`{\displaystyle \frac{1}{d_1}}\widehat{\stackrel{~}{p}}_2\widehat{q}_{}\widehat{e}_2^{}(\widehat{p}_1\widehat{k}_1+m)\widehat{e}_1^{}`$ (B.21)
$``$ $`{\displaystyle \frac{𝐪^2}{d^{}d_2^{}}}\widehat{e}_2^{}\widehat{e}_1^{}\widehat{\stackrel{~}{p}}_2+{\displaystyle \frac{s}{d^{}d_2^{}}}\widehat{e}_2^{}\widehat{q}_{}\widehat{e}_1^{}`$ (B.22)
$`+`$ $`{\displaystyle \frac{1}{d_1d_2^{}}}\widehat{e}_2^{}\widehat{q}_{}\widehat{\stackrel{~}{p}}_2(\widehat{p}_1\widehat{k}_1+m)\widehat{e}_1^{}`$ (B.23)
$`+`$ $`{\displaystyle \frac{1}{d^{}d_2^{}}}\widehat{e}_2^{}(\widehat{p}_1^{}+\widehat{k}_2+m)\widehat{e}_1^{}\widehat{q}_{}\widehat{\stackrel{~}{p}}_2\}u(p_1).`$ (B.24)
Calculating the contribution of the trace $`\mathrm{Sp}\{p_1^{}Qp_1\stackrel{~}{Q}\}`$ we neglect masses whose contribution to the quantity $`\mathrm{\Phi }^{\gamma \gamma }`$ may be restored using the general prescription . The corresponding correction has the form:
$`\mathrm{\Delta }_m\mathrm{\Phi }^{\gamma \gamma }`$ $`=`$ $`(1+𝒫_{12})\{{\displaystyle \frac{4m^2}{d_1^{}_{}{}^{}2}}{\displaystyle \frac{x_2^2y_1(1+y_1^2)}{(1x_2)^2}}`$ (B.25)
$`\times `$ $`{\displaystyle \frac{𝐪^2}{(𝐪y_1𝐩_1^{})^2(𝐪𝐩_1^{}/b)^2}}`$ (B.26)
$``$ $`{\displaystyle \frac{4m^2}{d_1^2}}{\displaystyle \frac{\beta _2^2z_1(1+z_1^2)𝐪^2}{(𝐪𝐩_1^{})^2(𝐩_1^{}(1\beta _2)𝐪)^2}}\},`$ (B.27)
where
$`y_1={\displaystyle \frac{1x_2}{b}},\beta _2={\displaystyle \frac{x_2}{1x_1}},z_1={\displaystyle \frac{b}{1x_1}}.`$ (B.28)
## Appendix C. Evaluation of 2-dimensional integrals
The azimuthal integration may be performed making use of the following equality:
$`J_{12\mathrm{}n}`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle _0^{2\pi }}𝑑\varphi {\displaystyle \underset{i}{}}[a_i+b_i\mathrm{cos}(\varphi \varphi _i)]^1`$ (C.1)
$`=`$ $`{\displaystyle \underset{k=1}{\overset{n}{}}}{\displaystyle \frac{1}{r_k}}{\displaystyle \underset{jk}{\overset{n}{}}}{\displaystyle \frac{b_k}{b_{kj}+\text{i}r_k\mathrm{sin}(\varphi _k\varphi _j)}},`$ (C.2)
with
$`r_i`$ $`=`$ $`\sqrt{a_i^2b_i^2},|a_i|>|b_i|,`$
$`b_{ij}`$ $`=`$ $`b_ia_jb_ja_i\mathrm{cos}(\varphi _i\varphi _j).`$
It is curious to note that the absence of the imaginary part provides an interesting algebraic identity. For $`n=2,n=3`$ it looks
$`J_{12}`$ $`=`$ $`{\displaystyle \frac{1}{d_{12}}}({\displaystyle \frac{b_1}{r_1}}b_{12}+{\displaystyle \frac{b_2}{r_2}}b_{21}),d_{12}=a_{12}^2r_1^2r_2^2,`$ (C.3)
$`a_{12}`$ $`=`$ $`a_1a_2b_1b_2\mathrm{cos}(\varphi _1\varphi _2),`$ (C.4)
$`J_{123}`$ $`=`$ $`{\displaystyle \frac{b_1^2}{r_1}}{\displaystyle \frac{a_{12}a_{13}r_1^2a_{23}}{d_{12}d_{13}}}+{\displaystyle \frac{b_2^2}{r_2}}{\displaystyle \frac{a_{21}a_{23}r_2^2a_{13}}{d_{12}d_{23}}}`$ (C.5)
$`+`$ $`{\displaystyle \frac{b_3^2}{r_3}}{\displaystyle \frac{a_{31}a_{32}r_3^2a_{12}}{d_{31}d_{32}}}.`$ (C.6)
This form is convenient for a subsequent integration over $`\text{d}𝐤_1^2`$.
## Appendix D. NLO contributions from virtual and soft photon emission
To avoid the misprints we use here the notations of the paper
$`s`$ $`=`$ $`d_1^{},t=d_1,u=Q^2,`$ (D.1)
$`s+t+u`$ $`=`$ $`q^2,\stackrel{~}{f}(s,t)=f(t,s),a=s+t,`$ (D.2)
$`b`$ $`=`$ $`s+u,c=u+t.`$ (D.3)
The quantities $`\tau _{ij}`$ encountered in the text (see Eq. (16)) may be written as
$`\tau _{11}`$ $`=`$ $`G(1+{\displaystyle \frac{u^2}{s^2}})\stackrel{~}{G}(2+{\displaystyle \frac{b^2}{t^2}})+2[{\displaystyle \frac{b^2}{st}}+{\displaystyle \frac{2u}{a}}`$ (D.4)
$`+`$ $`{\displaystyle \frac{2}{a^2}}(u^2bt)]l_{qu}+{\displaystyle \frac{b^2}{tc^2}}(2c+t)l_{qs}+{\displaystyle \frac{2us}{s}}l_{qt}`$ (D.5)
$`+`$ $`{\displaystyle \frac{1}{q^2}}\left[{\displaystyle \frac{4}{a}}(btu^2)4u2q^2+t{\displaystyle \frac{b^2}{c}}\right],`$ (D.6)
$`\tau _{12}`$ $`=`$ $`{\displaystyle \frac{c}{s^2}}(us)G+{\displaystyle \frac{1}{t^2}}(uq^2st)\stackrel{~}{G}2[{\displaystyle \frac{uq^2}{st}}+{\displaystyle \frac{2us+t}{a}}`$ (D.7)
$`+`$ $`{\displaystyle \frac{2}{a^2}}(u^2cs)]l_{qu}+{\displaystyle \frac{2c+t}{c^2}}(s{\displaystyle \frac{u}{t}}q^2)l_{qs}`$ (D.8)
$``$ $`{\displaystyle \frac{c}{bs}}(2us)l_{qt}+{\displaystyle \frac{1}{q^2}}[{\displaystyle \frac{4}{a}}(u^2cs)+8u+3t`$ (D.9)
$``$ $`s+{\displaystyle \frac{2}{c}}us],`$ (D.10)
and the additional notations look
$`l_{qu}`$ $`=`$ $`\mathrm{ln}{\displaystyle \frac{q^2}{u}},l_{qs}=\mathrm{ln}{\displaystyle \frac{q^2}{s}},l_{qt}=\mathrm{ln}{\displaystyle \frac{q^2}{t}},l_{ut}=\mathrm{ln}{\displaystyle \frac{u}{t}},`$ (D.11)
$`G`$ $`=`$ $`l_{qu}(l_{qt}+l_{ut})+2\text{Li}_2\left(1{\displaystyle \frac{t}{q^2}}\right)`$ (D.13)
$`2\text{Li}_2\left(1{\displaystyle \frac{q^2}{u}}\right)2\text{Li}_2(1).`$
## Appendix E. Semi-collinear kinematics of pair creation
The matrix element in the kinematics (2) may be put in a form (we extract the coupling constant):
$$M^{(1)}=\frac{1}{q_1^2}J_\nu I_\mu g^{\mu \nu },J_\nu =\overline{u}(p_{})\gamma _\nu u(p_1),$$
(E.1)
where the current $`I`$ describes a pair production by the photon with momentum $`q_1`$ off a proton. Using the Sudakov form of the 4-vectors $`p_{}`$ and $`q`$ with basic 4-vectors $`p_1`$ and $`p_2`$,
$$p_{}=\alpha _{}\stackrel{~}{p}_2+\beta _{}\stackrel{~}{p}_1+p_{},q=\alpha _q\stackrel{~}{p}_2+\beta _q\stackrel{~}{p}_1+q_{},$$
the representation of the metric tensor
$$g_{\nu \mu }=g_{\nu \mu }+\frac{2}{s}p_{2\nu }p_{1\mu }$$
and the gauge condition
$$Iq=I(\beta _qp_1+q_{})=0,\beta _q+\beta _{}=1,$$
we obtain for the matrix element squared and summed over spin states of electron:
$$|M^{(1)}|^2=\frac{1}{(q_1^2)^2}\left[2q_1^2𝐈^2+\frac{8}{\beta _q^2}\left(𝐩_{}𝐈\right)^2\right].$$
(E.2)
To calculate the quantity $`𝐈^2`$, we again present it in a form
$`I`$ $`=`$ $`e_{q_1}𝐈=e_{q_1}^\mu e_q^\nu {\displaystyle \frac{2s|\stackrel{}{q}|}{q^2s_1}}p_{2\rho }Y_\rho \overline{u}(p_1^{})O_{\mu \nu }v(p_+),`$ (E.3)
$`s_1`$ $`=`$ $`(p_2+q_1)^2,Y_\rho =\overline{u}(p_2)\mathrm{\Gamma }_\rho u(p_2^{}).`$ (E.4)
The phase volume is transformed the way to take the following form
$$\text{d}\mathrm{\Gamma }_4=(2\pi )^8\frac{1}{8s\beta _{}\beta _+b}\text{d}^2q\text{d}^2p_{}\text{d}\beta _{}.$$
(E.5)
Using
$$|\overline{u}(p_1^{})O_{\mu \nu }v(p_+)e_{q_1}^\mu e_q^\nu |^2=8\left[\frac{b}{\beta _+}+\frac{\beta _+}{b}\right],$$
we obtain the result for the cross section given in the text.
For the kinematics of bremsstrahlung mechanism the matrix element has a form
$$M^{(2)}=\frac{1}{k_1^2}I_\mu J_\nu g^{\mu \nu },k_1=p_++p_1^{}.$$
(E.6)
Here it is suitable to use alternative basis vectors of Sudakov parameterization
$`p_+`$ $`=`$ $`\alpha _+q+b_+\stackrel{~}{p}_1^{}+p_+,k_1=a_1q+b_1\stackrel{~}{p}_1^{}+k_1,`$ (E.7)
$`g_{\mu \nu }`$ $`=`$ $`g_{\mu \nu }+{\displaystyle \frac{2}{\stackrel{~}{s}}}q^\nu p_1^{}_{}{}^{}\mu ,k_1^2={\displaystyle \frac{𝐩_+^2+m^2b_1^2}{b_11}}>0.`$ (E.8)
Quite the same manipulations give
$$|M^{(2)}|^2=2k_1^2𝐈^2\frac{8}{b_1^2}\left(𝐤_1𝐈\right)^2.$$
Performing the integration over $`\text{d}^2(p_+)_{}`$ to a logarithmic accuracy and expressing the parameter $`b_1`$ in terms of the standard Sudakov decomposition with basic 4-vectors $`p_1,p_2`$
$$b_1=\frac{1\beta _{}}{b},$$
we immediately obtain the result given in the text.
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# Disorder Effects in Superconducting Multiple Loop Quantum Interferometers
\[
## Abstract
A theoretical study is presented on a number $`N`$ of resistively shunted Josephson junctions connected in $`\mathrm{𝑝𝑎𝑟𝑎𝑙𝑙𝑒𝑙}`$ as a disordered $`1D`$ array by superconducting wiring in such a manner that there are $`N1`$ individual SQUID loops with arbitrary shape formed. Under a constant current bias $`I>I_c`$, and irrespective of the degree of the disorder, but depending on the strength of magnetic field $`B`$, all junctions in the array oscillate at the $`\mathrm{𝑠𝑎𝑚𝑒}`$ frequency $`\nu _B`$. Computer simulations of the full non linear dynamics of a disordered junction array reveal: (i) the frequency $`\nu _B`$ is $`\mathrm{𝑛𝑜𝑡}`$ a periodic function of $`B`$, (ii) in the overdamped junction regime $`\nu _B`$ displays a sharp $`\mathrm{𝑔𝑙𝑜𝑏𝑎𝑙}`$ minimum around $`B=0`$. For zero inductive coupling the problem becomes equivalent to a $`\mathrm{𝑣𝑖𝑟𝑡𝑢𝑎𝑙}`$ single junction model.
PACS numbers: 85.25.Dq, 85.25.Am, 85.25.Cp
\]
So-called weak links, or Josephson junctions, are the basic active elements of superconductor quantum electronics. A key feature of a weak link between two superconductors, $`1`$ and $`2`$, is the property that there can flow a dissipationless macroscopic supercurrent $`I_s\left(\phi \right)`$ due to the tunneling of Cooper pairs with charge $`2e`$. This supercurrent depends on the gauge invariant phase difference $`\phi =\mathrm{\Theta }_1\mathrm{\Theta }_2+\frac{2e}{\mathrm{}c}_1^2<d𝐬,𝐀>`$ of the macroscopic BCS pairing wavefunctions on either side of the weak link. Josephson junctions made with modern fabrication techniques often have a sandwich type layered geometry, with a thin non superconducting tunneling barrier in the middle between two thick superconducting electrodes. In recent time also other types of weak links, for example of the bi-crystal type, became important in high-temperature superconductors. For an ideal S-I-S junction the supercurrent is connected to the phase difference $`\phi `$ across the tunneling barrier by $`I_s\left(\phi \right)=I_c\mathrm{sin}\phi .`$ It is important to realize that the supercurrent $`I_s`$ flows stationary provided it does not exceed a characteristic critical current $`I_c`$, the so-called Josephson critical current, which determines the maximum dissipationless current that can flow across a tunneling barrier. In general, $`I_c`$ depends on the material properties of the junction, on temperature $`T`$, and also on magnetic field $`𝐁=rot𝐀`$. Applying to a Josephson junction a bias current $`I`$ with a constant strength $`I>`$ $`I_c`$, there appears a rapidly oscillating voltage signal $`V\left(t\right)`$ across the junction, which determines the rate of change of the time dependent phase difference $`\phi \left(t\right)`$ according to
$$\mathrm{}_t\phi \left(t\right)=2eV\left(t\right)$$
(1)
This is the fundamental non stationary Josephson relation which governs the physics of weak superconductivity. So, for $`I>`$ $`I_c`$ there flows, besides the dissipationless supercurrent $`I_s`$, also a dissipative normal current $`I_n`$ in the junction, whose physical origin is the transfer of single (unpaired) electrons.
Within the range of validity of the RCSJ model, the dissipative current may be described with sufficient accuracy as a superposition of an ohmic current, characterized by a parallel ohmic shunt resistance $`R`$, and a displacement current, which is characterized by a parallel geometric shunting capacitance $`C`$ describing electric polarization inside the tunnelling barrier. The total junction current $`I`$ is then:
$$I=I_c\mathrm{sin}\phi \left(t\right)+\frac{\mathrm{}}{2eR}_t\phi \left(t\right)+\frac{\mathrm{}C}{2e}_t^2\phi \left(t\right)$$
(2)
The time average
$$V=\underset{t\mathrm{}}{lim}\frac{1}{t}_0^t𝑑t^{}V(t^{})=\frac{\mathrm{}}{2e}\underset{t\mathrm{}}{lim}\frac{\phi \left(t\right)\phi \left(0\right)}{t}$$
(3)
is the dc voltage part of the in general not sinusoidal voltage signal $`V(t)`$ across the electrodes of a Josephson junction. For a strongly overdamped junction, $`C=0`$, one finds, assuming a constant bias current $`I>I_c`$, a relatively simple formula: $`V=R\sqrt{I^2I_c^2}`$. The dc voltage $`V`$ is connected to the oscillation frequency $`\omega =2\pi \nu `$ of the voltage signal $`V\left(t\right)`$ by:
$$h\nu =2eV$$
(4)
This result for the voltage response function $`V`$ of a weak link suggests a spectroscopic interpretation. When a Cooper pair is transferred from the superconducting side $`1`$ to the superconducting side $`2`$ of the junction, under conditions where $`I>I_c`$, a microwave photon with energy $`2eV`$ is released in the form of one quantum of electromagnetic radiation (Josephson radiation).
As far as macroscopic quantum interference is concerned, it was actually known long before the discovery of the Josephson effects, that magnetic flux threading the area of a superconducting ring, made out of a material that is thick compared to the magnetic penetration depth $`\lambda `$, should be quantized in units of the flux quantum $`\mathrm{\Phi }_0=\frac{hc}{2e}`$.
Technical applications of the physics of weak superconductivity include ultrasensitive quantum interferometers, which indeed combine the afore mentioned Josephson effects with flux quantization. Consider, as indicated schematically in Fig.(1a), a standard two junction SQUID (for simplicity with symmetric junction parameters) under the dc current bias $`I>`$ $`2I_c`$. Such a device is actually a flux-to-voltage transducer. Let $`\mathrm{\Phi }=𝐁,𝐚=`$ $`\left|𝐁\right|\left|𝐚\right|\mathrm{cos}\alpha `$ be the magnetic flux threading the orientated area element $`𝐚`$ of the superconducting SQUID loop, where $`\alpha `$ is the angle between the normal vector of the orientated area element and the magnetic field vector $`𝐁`$, as depicted schematically in Fig.(1a). The total magnetic field, $`𝐁=𝐁^{\left(1\right)}+𝐁^{\left(2\right)}`$, is then a superposition of the primary external magnetic field $`𝐁^{\left(1\right)}`$, which generates the flux $`\mathrm{\Phi }^{\left(1\right)}=𝐁^{\left(1\right)},𝐚`$ one wants to detect, and a secondary magnetic field $`𝐁^{\left(2\right)}`$ that results, for example, from the inductance $`L`$ (or other impedance effects) in the circuit. The screening current $`I_{sc}`$ circulating in the SQUID loop leads to a total flux $`\mathrm{\Phi }=\mathrm{\Phi }^{\left(1\right)}+\mathrm{\Phi }^{\left(2\right)}`$. Dependent on the secondary flux term, $`\mathrm{\Phi }^{\left(2\right)}=LI_{sc}`$, there exists an optimal size $`\left|𝐚_L\right|`$ for any SQUID loop . A dimensionless measure for the inductance of such a loop is $`\beta _L=\frac{LI_c}{\mathrm{\Phi }_0}`$. Note also that a two junction SQUID cannot be directly employed as a detector of absolute strength of external magnetic field. This is because the voltage response function $`V_{xy}`$ of the SQUID, i.e. the time average of the rapidly oscillating voltage signal $`V_{xy}\left(t\right)`$ across the nodes $`x`$ and $`y`$ of the circuit, is a periodic function of the strength of external magnetic field, see Fig.(1a).
A straightforward extension of the standard two junction SQUID is sketched in Fig.(1b). This is a $`1D`$ array of $`N`$ adjacent Josephson junctions connected in parallel. In particular, the area elements of the $`N1`$ SQUID loops formed in this manner are all equal, e.g. $`𝐚_n=𝐚_L`$ for all $`n`$. The voltage response signal $`V_{xy}`$ vs. strength $`\left|𝐁^{\left(1\right)}\right|`$ of external magnetic field of this array has the same period than a standard two junction SQUID with loop area $`\left|𝐚_L\right|`$, see Fig.(1b).
A more general quantum interference device is obtained when the area elements $`𝐚_n`$ of the $`N1`$ loops in the array differ in size and, possibly, in orientation, as depicted schematically in Fig.(1c). If the sizes $`\left|𝐚_n\right|`$ of the orientated area elements $`𝐚_n`$ of the individual superconducting loops are chosen in a disordered fashion the voltage response function $`V_{xy}`$ vs. $`\left|𝐁^{\left(1\right)}\right|`$ becomes nonperiodic, see Fig.(1c). Taking into account inductive couplings among the currents in the circuit, the maximum loop size in the disordered array should coincide with the corresponding optimal loop size of a standard two junction SQUID, i.e. $`\mathrm{max}\left|𝐚_n\right|=\left|𝐚_L\right|`$. The voltage response signal $`V_{xy}`$ vs. strength of magnetic field of a disordered junction array is, under a suitable dc current bias $`I`$, indeed a unique function of $`\left|𝐁\right|`$ around its narrow global minimum at $`\left|𝐁\right|=0`$. This suggests that it should be possible, e.g. by measuring control current(s) flowing through the wires of a set of suitably orientated compensation coil(s), to reconstruct absolute strength, orientation and even the phase of an incident primary magnetic field signal, i.e. to determine the full vector $`𝐁^{\left(1\right)}(t)`$.
The $`n`$-th Josephson junction in the array has, within the range of validity of the RCSJ model, optional individual junction parameters $`R_n`$, $`C_n`$ and $`I_{c,n}`$. The corresponding current $`I_n`$ flowing through the $`n`$-th Josephson junction is, according to Eq.(2), determined by the gauge invariant phase difference $`\phi _n\left(t\right)`$ across that junction. The total current $`I`$ flowing through the nodes $`x`$
and $`y`$, respectively, of the circuit is then obtained from Kirchhoff’s rule as the phase sensitive superposition of the individual junction currents $`I_n`$:
$`I={\displaystyle \underset{n=0}{\overset{N1}{}}}\left[I_{c,n}\mathrm{sin}\phi _n(t)+({\displaystyle \frac{\mathrm{}C_n}{2e}}_t^2+{\displaystyle \frac{\mathrm{}}{2eR_n}}_t)\phi _n\left(t\right)\right]`$ (5)
Note that the gauge invariant phase differences $`\phi _n`$ of adjacent Josephson junctions in the array are not independent, but are connected to each other by the condition of flux quantization:
$$\phi _n\phi _{n1}=\frac{2\pi }{\mathrm{\Phi }_0}𝐁,𝐚_nmod2\pi $$
(6)
Here $`\left|𝐚_n\right|`$ is the area of the superconducting loop connecting adjacent Josephson junctions numbered as $`n`$ and $`n1`$, respectively, and $`𝐁`$ denotes the magnetic field threading the orientated area element $`𝐚_n`$ of this loop. Note that Eq.(6) applies quite generally, provided the superconducting material, out of which the connecting loops are made, is thick compared to the magnetic penetration depth $`\lambda `$. In this case there exists a path inside the wire connecting, say, junction $`n`$ with its neighbor junction $`n1`$, on which the superfluid velocity field $`𝐯_s`$ becomes negligible small. So, $`\mathrm{}\mathrm{\Theta }=\frac{2e}{c}𝐀`$ along this path. Since all junctions in the array are connected in parallel, the rapidly oscillating voltage $`V_n(t)`$ at the electrodes of a particular Josephson junction, numbered as $`n`$in the array, is related to the signal $`V_{xy}(t)`$ between the nodes $`x`$ and $`y`$ of the circuit by
$$V_{xy}(t)=V_n(t)+_{xny}d𝐬,𝐄\left(t\right)$$
(7)
By Faraday’s law the electric field $`𝐄`$ along an integration path $`xny`$, that starts at node $`x`$, traverses the tunneling barrier of the $`n`$-th Josephson junction just once, and then terminates at node $`y`$, is directly connected to the time derivative of the flux threading the area elements of the 1D array. Once the signal $`V_0(t)`$ $`=`$ $`\frac{\mathrm{}}{2e}_t\phi _0(t)`$ is known the other voltage signals $`V_n(t)`$ across the electrodes of the $`n`$-th junction follow from
$$V_n(t)V_{n1}(t)=\frac{1}{c}_t𝐁\left(t\right),𝐚_n$$
(8)
Taking into account the Biot-Savart type inductive couplings among the currents flowing in the circuit prohibits further simplification. However, it follows directly from Eq.(6) that one can indeed eliminate from Eq.(5) all phase variables $`\phi _n(t)`$ in favor of a single phase, for example $`\phi _0(t)`$, provided the extremely simplifying assumption is made, that all inductive couplings vanish. In this case the problem of $`N`$ coupled Josephson junctions is mapped onto a virtual single Josephson junction model. With $`\frac{1}{R}=\frac{1}{N}_{n=0}^{N1}\frac{1}{R_n}`$, $`T_N=\frac{\mathrm{}}{2e}\frac{1}{I_cR}`$ and $`\varphi \left(t\right)=\phi _0(t)`$ there results a scalar differential equation determining the phase difference $`\varphi \left(t\right)`$:
$`\left|S_N\left(𝐁\right)\right|\mathrm{sin}\left[\varphi \left(t\right)+\delta _N\left(𝐁\right)\right]+T_N\left(RC_t^2+_t\right)\varphi \left(t\right)`$ (9)
$`=J_N{\displaystyle \frac{2\pi }{\mathrm{\Phi }_0}}T_N\left(RC_t^2𝐁\left(t\right),𝐚_C+_t𝐁\left(t\right),𝐚_R\right)`$ (10)
where we have defined ($`𝐚_0=\mathrm{𝟎}`$)
$`S`$ $`{}_{N}{}^{}(𝐁)={\displaystyle \frac{1}{N}}{\displaystyle \underset{n=0}{\overset{N1}{}}}{\displaystyle \frac{I_{c,n}}{I_c}}\mathrm{exp}\left[{\displaystyle \frac{2\pi i}{\mathrm{\Phi }_0}}{\displaystyle \underset{m=0}{\overset{n}{}}}𝐁,𝐚_m\right]`$ (11)
$`𝐚`$ $`{}_{R}{}^{}={\displaystyle \frac{1}{N}}{\displaystyle \underset{n=0}{\overset{N1}{}}}{\displaystyle \frac{R}{R_n}}{\displaystyle \underset{m=0}{\overset{n}{}}}𝐚_m,𝐚_C={\displaystyle \frac{1}{N}}{\displaystyle \underset{n=0}{\overset{N1}{}}}{\displaystyle \frac{C_n}{C}}{\displaystyle \underset{m=0}{\overset{n}{}}}𝐚_m`$ (12)
and $`J_N=\frac{I}{NI_c}`$, $`I_c=\frac{1}{N}_{n=0}^{N1}I_{c,n}`$, $`C=\frac{1}{N}_{n=0}^{N1}C_n`$. The complex function $`S_N\left(𝐁\right)=\left|S_N\left(𝐁\right)\right|e^{i\delta _N\left(𝐁\right)}`$ denotes the characteristic structure factor of the 1D Josephson junction array, as defined in Eq.(11). It is an extremely responsive function of strength and orientation of magnetic field, and it is strongly affected by the choice of the individual area elements $`𝐚_m`$. In general $`\left|S_N\left(𝐁\right)\right|`$ is very sensitive to permutations among the $`𝐚_m`$’s.
In the overdamped junction regime, $`C=0`$, under conditions where a constant current $`I`$ is biased such that $`1\left|S_N\left(𝐁\right)\right|/J_N`$ $`\mathrm{sin}\alpha _B`$, and assuming for simplicity a homogeneous static magnetic field $`𝐁`$ (and also time independent area elements $`𝐚_m`$), one finds an exact solution for the phase difference $`\phi _0(t)`$:
$`V_0(t)={\displaystyle \frac{\mathrm{}}{2e}}_t\phi _0(t)=I_cR{\displaystyle \frac{J_N^2\left|S_N\left(𝐁\right)\right|^2}{J_N+\left|S_N\left(𝐁\right)\right|\mathrm{sin}\left(\omega _Bt\alpha _B\right)}}`$
For a static magnetic field $`𝐁`$ the voltage response function $`V_{xy}`$ measured between the nodes $`x`$ and $`y`$ of the circuit is equal to the dc-part of the rapidly oscillating voltage signal $`V_0(t)`$. All Josephson junctions in the 1D array oscillate at the same frequency $`\omega _B=2\pi \nu _B`$, which is related to $`V_{xy}`$ by:
$$\frac{h}{2e}\nu _B=V_0=I_cR\sqrt{J_N^2\left|S_N\left(𝐁\right)\right|^2}=V_{xy}$$
(13)
Note that the oscillation frequency $`\nu _B`$ of such a local oscillator is even more sensitive to changes of strength or orientation of the external magnetic field $`𝐁`$ than the structure factor of the array itself, since $`\left|S_N\left(𝐁\right)\right|`$ enters Eq.(13) quadratically.
Consider, as a special case, an ordered array, consisting of $`N1`$ identical SQUID loops, such that $`𝐁,𝐚_n=\mathrm{\Phi }=𝐁,𝐚_L`$, and $`I_{c,n}=I_c`$ independent on the junction index $`n`$. Then the structure factor $`S_N\left(𝐁\right)`$ $`S_N^{\left(\mathrm{\Phi }\right)}`$ becomes a simple geometrical series:
$$S_N^{\left(\mathrm{\Phi }\right)}=\frac{\mathrm{sin}\left(\pi \frac{\mathrm{\Phi }}{\mathrm{\Phi }_0}N\right)}{N\mathrm{sin}\left(\pi \frac{\mathrm{\Phi }}{\mathrm{\Phi }_0}\right)}\mathrm{exp}\left[\pi i\frac{\mathrm{\Phi }}{\mathrm{\Phi }_0}\left(N1\right)\right]$$
(14)
In Fig.(1b) one observes the usual narrowing proportional to $`\frac{1}{N}`$ of the width of the voltage response signal $`V_{xy}`$ around its minima. Note the periodicity property $`\left|S_N^{\left(\mathrm{\Phi }+\mathrm{\Phi }_0\right)}\right|=\left|S_N^{\left(\mathrm{\Phi }\right)}\right|`$ for all $`N2`$. For $`N=2`$ Eq.(13) is the periodic voltage response of a symmetric two junction SQUID in the overdamped junction regime.
A structure factor with a much longer period may be obtained in a parallel junction array where the orientated area elements increase in size according to a linear relation:
$$𝐚_m=(2m1)𝐚_1$$
(15)
For simplicity, identical junction parameters $`R_n`$, $`C_n`$ and $`I_{c,n}`$ are assumed. Then:
$$S_N(𝐁)=\frac{1}{N}\underset{n=0}{\overset{N1}{}}\mathrm{exp}\left[2\pi i\frac{𝐁,𝐚_1}{\mathrm{\Phi }_0}n^2\right]$$
(16)
The total area occupied by such a Gaussian array is $`\left(N1\right)^2𝐚_1`$, where $`𝐚_1`$ is the smallest area element, and $`𝐚_{N1}=\left(2N3\right)𝐚_1`$ is the largest area element. Compare a Gaussian array, with $`N1`$ area elements as described in Eq.(15), with a periodic array, consisting of $`N_P1`$ identical SQUID-loops with size $`\left|𝐚_L\right|`$. For a useful comparison, both arrays should occupy the same total area: $`\left(N1\right)^2𝐚_1=\left(N_P1\right)𝐚_L`$. Also the largest area element in the Gaussian array should coincide with the area element of an optimal single SQUID-loop, i.e. $`𝐚_{N1}=𝐚_L`$. Both requirements together imply for $`N_P1`$ that the Gaussian array has the double number of junctions compared to a corresponding periodic junction array: $`N2N_P`$. To determine the period of the Gaussian array consider a case where the flux threading the area of the smallest element, $`𝐚_1`$, is equal to a rational multiple of half a flux quantum: $`𝐁,𝐚_1=\frac{M}{N}\frac{\mathrm{\Phi }_0}{2}`$. Then the largest area element in the array, $`𝐚_{N1}`$, is threaded by a flux $`\mathrm{\Phi }_M=\left(1\frac{3}{2N}\right)M\mathrm{\Phi }_0`$. In this case the structure factor $`S_N(𝐁)S_N^{\left(\mathrm{\Phi }_M\right)}`$ may be determined using a result of C.F. Gauss:
$$S_N^{\left(\mathrm{\Phi }_M\right)}=\frac{1}{N}\underset{n=0}{\overset{N1}{}}e^{\pi i\frac{M}{N}n^2}=\frac{e^{i\frac{\pi }{4}}}{\sqrt{NM}}\underset{n=0}{\overset{M1}{}}e^{\pi i\frac{N}{M}n^2}$$
(17)
Note the periodicity $`\left|S_N^{\left(\mathrm{\Phi }_M+\mathrm{\Phi }_{2N}\right)}\right|=\left|S_N^{\left(\mathrm{\Phi }_M\right)}\right|`$, with period $`\mathrm{\Phi }_{2N}=\left(2N3\right)\mathrm{\Phi }_0`$. Remarkably, for $`M=2`$, and $`N=N_1N_2`$ being the product of two prime numbers $`N_1`$ and $`N_2`$, there holds the factorization:
$$S_N^{\left(\mathrm{\Phi }_2\right)}=\left(1\right)^{\frac{\left(N_11\right)\left(N_21\right)}{4}}S_{N_1}^{\left(\mathrm{\Phi }_2\right)}S_{N_2}^{\left(\mathrm{\Phi }_2\right)}$$
(18)
Apparently, such Gaussian junction arrays are governed by the laws of number theory (quadratic residues).
The long periodicity of the structure factor vs. flux $`\mathrm{\Phi }^{\left(1\right)}`$ threading the largest area element $`𝐚_L`$ of Gaussian junction arrays is also visible in the calculated voltage response function $`V_{xy}`$, irrespective of the degree of the inductive coupling represented by the parameter $`\beta _L`$. This is illustrated in Fig.(2). Note the asymmetry of $`V_{xy}`$ under $`\mathrm{\Phi }\mathrm{\Phi }`$ (for a constant bias current $`I`$) for finite inductive coupling. As far as disorder is concerned, we also find that $`V_{xy}`$ in Gaussian junction arrays is very responsive to adding small random fluctuations to the size distribution of the area elements, so that $`V_{xy}`$ becomes non periodic with a pronounced antipeak only around $`\mathrm{\Phi }=0`$. As $`\beta _L`$ increases, the difference $`\mathrm{max}V_{xy}\mathrm{min}V_{xy}`$ decreases, and the linewidth of the global minimum $`V_{xy}`$ around $`\mathrm{\Phi }=0`$ ceases in this case to scale proportional to $`\frac{1}{N}`$. Note all this applies in the overdamped junction regime. Our results for weak damping will be published elsewhere .
We hope that an experimental verification of the predicted magnetic field dependence of the voltage response function of disordered $`1D`$ parallel Josephson junction arrays will stimulate the development of new types of robust superconducting quantum interferometers, which would allow (for the first time) a technically rather simple precision measurement of absolute strength of external magnetic fields.
Acknowledgments: We thank R.P. Huebener, R. Kleiner and T. Träuble for useful discussions. Support by ”Forschungsschwerpunktprogramm des Landes Baden-Württemberg” is gratefully acknowledged.
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# 1 Introduction
## 1 Introduction
Since its introduction, string field theory held out the promise for nonperturbative studies of string theory. Recently Sen has argued that open bosonic (string) field theory describes the dynamics of $`D\overline{D}`$ system with the tachyon providing the instability inherent in such pair. Tachyon condensation also describes the decay of a single unstable $`D`$-brane. In addition, according to the kink and lump solutions of such field theory lead to lower dimensional branes. In recent work Sen and Zweibach studied in detail tachyon condensation in 26D open string theory. This follows the earlier, pioneering work of Samuel and Kostelecky who were the first to consider the vacuum structure of string field theory . They used a level truncation scheme to generate an approximation for the tachyon effective potential. Following this scheme, results were found that show great agreement in numerical values with the expected exact results. Similar results were recently obtained for superstring field theory . Impressive high level studies appeared in .
It is equally important to give a construction of non-constant kink and lump-like solutions. One can expect that it is for these that the stringy effects present in string field theory might play the most important role. One of the characteristic features present in construction of Witten’s version of the theory is the appearance in the interaction of terms exponential in derivatives. These terms can be moved from the interaction to give a nontrivial kinetic term. They have frustrated early attempts for construction of nonperturbative soliton (or instanton) solutions of the theory.
In the present work, we consider this problem in numerical terms. Concentrating on the lump of open string theory, we develop methods for its numerical solution. In this we keep the nontrivial exponential terms characteristic of the string vertex interaction and simultaneously perform a level truncation. This, we argue is to provide a very good approximation to the exact result.
The content of the paper is as follows. After a short description of open SFT, we discuss some features relevant to the present work. We explain (based on earlier observations ) how and why the approach of keeping higher derivatives and simultaneous level truncation holds the promise for a good approximation. We then proceed to the numerical work.
While this work was in progress, there appeared the work of Ref. which considers the problem in its field theory limit.
## 2 Open String Field Theory
We begin by describing some features of open string field theory which are of relevance to the investigation that follows. One has the cubic action:
$$S=A|Q|A+\frac{g}{3}V_3||A|A|A$$
(1)
with Q being the first quantized BRST operator. The geometric, three string interaction is realized in the Hilbert space by the vertex
$$V_3|=0|\mathrm{exp}\{\alpha _{n}^{}{}_{}{}^{r}N_{nm}^{rs}\alpha _m^s+\alpha ^{}\mathrm{ln}\gamma \underset{r=1}{\overset{3}{}}_r^2\}$$
(2)
Here the Neumann coefficient $`N^{rs}`$ are determined in terms of appropriate conformal mapping ,their explicit values are determined in .
One of the main properties of the three-string vertex is an explicit appearance of higher derivative terms. They come in exponential form acting on each string field
$$|A\mathrm{exp}\left(\alpha ^{}\mathrm{ln}\gamma _\mu ^\mu \right)|A$$
(3)
with the constant
$$\gamma =\frac{3\sqrt{3}}{4}$$
(4)
The exponential terms are not relevant in studies of vacuum structure but they can have a nontrivial effect in any other nonzero momentum process. Concerning a systematic approximation or expansion scheme one notes the following. The masses of tachyon (and higher mass fields) are proportional to $`1/\alpha ^{}`$. In general $`\alpha ^{}`$ serves as a scale ,it can be scaled out in front of the SFT action. At a nonperturbative level, there is no small free parameter and no systematic expansion. Since one is not able(at present) to solve the theory in exact terms, one relies on seemingly and hoc approximations. Such is the process of level truncation. In order to understand more clearly the procedure involved and the relative relevance of particular terms, let us recall the original argument given for the level truncation in an unpublished work of ref.. Considering an approximate calculation of a nonzero momentum amplitude, for example for four tachyons one starts from:
$$A_s=_0^1𝑑xx^{s/22}V_{34}(\overline{3^{}})|b_0x^R|V_{12}(3^{}),$$
(5)
representing the $`s`$-chanel Feynman diagram.
Using the fact that $`x^R\alpha _nx^R=x^n\alpha _n`$, as well as similar results for $`b^{}s`$ and $`c^{}s`$, we have
$$A_s=_0^1𝑑xx^{s/22}V_{34}(\dot{3})b_0|V_{12}(\dot{3}).$$
(6)
The dot on $`V`$ indicates that $`a_n,b_n,c_n`$ have been replaced by $`x^n\alpha _n,x^nb_n,x^nc_n`$. Expanding in levels corresponds to expanding in powers of $`x=e^r`$. With the use the appropriate Neumann coefficients after a straightforward algebra one has
$$A_s=_0^1𝑑xx^{s/22}\left(1\frac{11^2}{3^6}x^2+\mathrm{}\right)e^{E(x)}.$$
(7)
The term in the exponent has the expansion
$`E(x)={\displaystyle \frac{s}{2}}\mathrm{ln}\gamma `$ $``$ $`\left({\displaystyle \frac{s}{2}}+2\right)\left({\displaystyle \frac{2^3}{3^3}}x+{\displaystyle \frac{2^219}{3^6}}x^2+\mathrm{}\right)`$ (8)
$``$ $`\left({\displaystyle \frac{i}{2}}+2\right)\left({\displaystyle \frac{2^4}{3^3}}x{\displaystyle \frac{2^67^2}{3^{10}}}x^3+\mathrm{}\right)`$
$``$ $`\left({\displaystyle \frac{2^4}{3^3}}x+{\displaystyle \frac{2^3}{3^6}}x^2+\mathrm{}\right)+\left({\displaystyle \frac{265^2}{23^6}}x^2+\mathrm{}\right).`$
This result can easily be rearranged into a form
$$A_s=_0^{z(1)}𝑑zz^{s/22}(1z)^{t/22}$$
(9)
where
$$z(x)=\frac{1}{\gamma ^2}x\left(1\frac{2^3}{3^3}x+\frac{2^2}{3^3}x^2+\mathrm{}\right).$$
(10)
For agreement with the exact result one would need to have z(1)=0.5 (the s and t -chanel diagrams are to cover the full range (0,1).To the present order in the level expansion we have $`z(1)0.50480`$ which is indeed very close to the exact value. A more significant observation is the fact that the main effect is contained in the $`\gamma ^2`$ factor present in the above expression . That term itself gives $`z(1)0.6`$. If we look up the origin of this factor, we see that it comes directly from the exponential higher derivative operator present in the vertex: since the intermediate states in the four point amplitude calculation are not on shell the exponential terms contribute giving the corresponding $`s`$ dependence.
In this nonzero momentum example, we conclude that it is advantageous to keep the higher derivatives exactly and that this followed by a level truncation is likely to give a good overall approximation. Naturaly one still expects this expansion to only be good for certain range of momenta.
Consider then the string field theory with the tachyon (level 0), but with the higher derivative terms kept exactly. The action evaluated originally in reads
$$=\frac{1}{2}\left(_\mu T\right)^2+\frac{1}{2\alpha ^{}}T^2\frac{g}{3}\gamma ^3\stackrel{~}{T}^3\mathrm{with}\stackrel{~}{T}=e^{\mathrm{ln}\gamma ^2}T$$
(11)
The first sign that the presence of the exponential term in the cubic interaction profoundly influences the nature of the problem is seen in attempting to evaluate the asymptotics of a possible static solution. In ordinary field theories, kink and lump solution can be asymptoticaly characterized as
$$\varphi (x)\varphi _0+a_1e^{mx}$$
(12)
where $`\varphi _0`$ is the constant vacuum solution and $`m`$ is the physical mass of the scalar field. Based on this, one can write a systematic expansion for the lump
$$\varphi (x)=\underset{m=0}{\overset{\mathrm{}}{}}a_ne^{nmx}$$
(13)
with the classical equations providing a recursion formula for coefficients $`a_n`$. The exponential decay $`e^{nmx}`$ has a physical meaning, the lump form factor receives a contribution from $`n`$ mesons. In attempting an analogue expansion in the case of string field theory, one meets a surprise. After moving the exponential into the kinetic term or equivalently denoting $`\stackrel{~}{T}=\varphi `$ the string theory tachyon field equation reads
$$\left(\left(_x^2+1\right)e^{c_x^2}2\right)\varphi (x)=\overline{g}\varphi (x)^2$$
(14)
with $`c=2\mathrm{ln}\gamma =\mathrm{ln}\mathrm{\hspace{0.17em}3}^3/4^2`$. The Ansatz
$$\varphi (x)\varphi _0+a,e^{\overline{m}x}$$
(15)
leads to an eigenvalue equation
$$(\overline{m}^2+1)e^{c\overline{m}^2}2=0$$
(16)
for $`\overline{m}`$. This equation turns out to have no real solution. With the exponential present the nature of the lump solution has changed from that of ordinary field theory. In particular the decay at asymtopic infinity in string field theory has to be stronger than a simple exponential. The non-exponential decay of the full lump (kink) solution signals that the form factor will have a more complex physical meaning. This feature also necessitates a purely numerical approach to the problem which we attempt in the next section.
## 3 Calculations and Results
The prescence of the higher derivative terms in $`(14)`$ prevent a numerical analysis directly in $`x`$ space. Upon transforming to momentum space, one finds the following nonlinear integral equation
$$\left(\left(1\stackrel{}{k}^2\right)e^{c\stackrel{}{k}^2}2\right)\varphi (\stackrel{}{k})=\overline{g}d^{25p}q\varphi (\stackrel{}{q})\varphi (\stackrel{}{k}\stackrel{}{q}),$$
(17)
where $`\stackrel{}{k}`$ is a $`25p`$ dimensional vector. To solve an integral equation of this type, one typically discretizes the problem and solves the resulting matrix equation. In the case on hand, the resulting matrix equation is nonlinear and one is forced to search for the solution $`\varphi (\stackrel{}{k})`$ using a numerical minimization. It is computationally difficult and expensive to minimize functions with a large number of variables, and for this reason we have worked directly with a spherically symmetric ansatz for the tachyon field. Stable numerical solutions were obtained using a lattice having $`81`$ points. This implies a nonlinear minimization problem with $`81`$ parameters. The choice of a good objective function to be used in the minimization as well as an accurate initial guess are crucial to obtain a nontrivial solution. Indeed, in practice, we find that the majority of initial guesses lead to the trivial $`\varphi =0`$ solution. Our objective function was constructed as follows: Start by making an initial guess, $`\widehat{G}(k_n)`$ for the right hand side of $`(17)`$. This initial guess is used to compute a value for the tachyon wave function
$$\widehat{\varphi }(k_n)=\frac{\overline{g}\widehat{G}(k_n)}{\left(\left(1k_n^2\right)e^{ck_n^2}2\right)}.$$
(18)
The wave function $`\widehat{\varphi }(k_n)`$ can now be used to compute the value of the left hand side of $`(17)`$. However, this involves a convolution which must be performed with care. Evaluating the convolution directly in momentum space is not optimal: by exploiting the spherical symmetry of the wave function, the convolution can reduced to the integration over a single angle and a radius. The integrand contains the factor $`\widehat{\varphi }(p_n)\widehat{\varphi }(\delta _n)`$ with
$$\delta _n=|\stackrel{}{k}\stackrel{}{p}|=\sqrt{p_n^2+k_m^22k_mp_n\mathrm{cos}(\theta _l)}.$$
(19)
In general, $`\delta _n`$ does not correpsond to a lattice point and one is forced to interpolate from the known points to find $`\widehat{\varphi }(\delta _n)`$. A much more efficient way of performing the convolution is by transforming to $`x`$ space, squaring the function and then transforming back to $`k`$ space. By using the spherical symmetry of wave function, the Fourier transform can be reduced to a single integration. Thus, the previous double integration has been replaced by two single integrations, which is more efficient. In addition, the integrations only require a knowledge of $`\widehat{\varphi }`$ at the lattice points. After convolving $`\widehat{\varphi }`$ with itself to obtain $`G(k_n)`$, the error, $``$ to be minimized is constructed as
$$=\underset{n}{}|G(k_n)\widehat{G}(k_n)|.$$
(20)
The minimization was performed using the fmins.m subroutine of MATLAB, which employs a simplex search method. A suitable initial guess is obtained by taking a function which initially falls off slightly faster than a Gaussian, but reaches zero at some finite momentum. As an example we have shown the wavefunction for the D20 brane below.
> Fig1: The wave function of the D20 brane as a function of the radial coordinate $`r`$. The plot was obtained by taking the Fourier transform of the numerical solution of equation $`(18)`$.
The above form for the wave function is typical. The value of the wave function at the origin in position space decreased from $`0.62`$ for the D24 brane to $`4.46`$ for the D18 brane. The point at which the wave function reaches zero is very nearly constant for all the Dp-branes considered here. In figure 2 we have shown the error in the D24 brane solution. The tensions of these solutions was evaluated directly in momentum space. For example, in the case of the D24 brane, we compute
$`\widehat{T}_{24}=2\pi ^2T_{25}`$ $`{\displaystyle }`$ $`{\displaystyle \frac{dp}{2\pi }}({\displaystyle \frac{1}{2}}\varphi (p)(p^21)e^{cp^2}\varphi (p)+\varphi (p)\varphi (p)`$
$`+`$ $`{\displaystyle \frac{g}{3}}{\displaystyle }{\displaystyle \frac{dk}{2\pi }}\varphi (k)\varphi (p)\varphi (pk))`$
$`={\displaystyle \frac{2\pi ^2T_{25}}{6}}`$ $`{\displaystyle }`$ $`{\displaystyle \frac{dp}{2\pi }}\left(\varphi (p)(p^21)e^{cp^2}\varphi (p)+2\varphi (p)\varphi (p)\right).`$ (21)
> Fig2: The error in the wave function of the D20 brane as a function of the radial coordinate in momentum space $`p`$. The plot shows $`/(_n\widehat{G}(k_n)`$ as a function of $`k_n`$.
In our conventions, a Dp-brane has tension
$$T_p=(2\pi )^{25p}T_{25}.$$
(22)
In the table below we compare our numerically evaluated tensions, $`\widehat{T}_p`$ with the above values.
| $`p`$ | $`\mathrm{\Delta }`$ $`T_p`$ | $`\widehat{T}_p/T_p`$ | $`\varphi (0)`$ |
| --- | --- | --- | --- |
| 24 | -29.4 | 0.706 | -0.63 |
| 23 | -27.5 | 0.725 | -0.86 |
| 22 | -14.3 | 0.857 | -1.19 |
| 21 | -8.3 | 0.917 | -1.63 |
| 20 | 6.4 | 1.064 | -2.26 |
| 19 | 25.3 | 1.253 | -3.15 |
| 18 | 64.0 | 1.640 | -4.46 |
Table I
> Values of the Dp brane tensions computed using the numerical tachyon lump solution. The parameter $`\mathrm{\Delta }T_p`$ is defined as $`\mathrm{\Delta }T_p=\frac{\widehat{T}_pT_p}{T_p}\times 100.`$ $`\varphi (0)`$ is the value of the wavefunction at the origin in position space.
Clearly there is no obstacle for the existence of lower p-Brane solutions. The trend is that the static lump gives systematicaly a growing tension. This is actually opposite to what one finds in the extreme field theory limit. Concerning further improvement of present results one expects (esspecialy for lower p) the relevance of higher massive levels. It is next important to study the direction of their contributions. It is also relevant to perform a similar study in superstring theory.
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# Resistive axisymmetric equilibria with arbitrary flow
## Abstract
An analysis of axisymmetric equilibria with arbitrary incompressible flow and finite resistivity is presented. It is shown that with large aspect ratio approximation or vanishing poloidal current, a uniform conductivity profile is consistent with equilibrium flows. Also a comment made on coexistence of both toroidal and poloidal flows in an axisymmetric field-reversed configuration.
Calculating the equilibrium is one of the fundamental problems of magnetically confined plasmas. Most studies are directed toward finding an ideal (i.e. infinitely conducting) magnetohydrodynamic (MHD) equilibria in an axisymmetric plasma. The earliest calculations are those of Grad and Shafranov , leading to the famous Grad-Shafranov equation . The ideal and static Grad-Shafranov equation is an elliptic differential equation in the magnetic flux function $`\psi `$ with two arbitrary surface quantities as the pressure $`p(\psi )`$ and the poloidal current $`I(\psi )`$.
Consequently, there have been attempts to include various effects e.g. mass flow into the equilibrium equations. An equivalent of Grad-Shafranov equation in an axisymmetric ideal plasma with arbitrary flow has been given by Hameiri . In some recent works, Steinhauer deals with a generalization of Grad-Shafranov equilibria in a multi-fluid with flow and Throumoulopoulos and Tasso consider a helically symmetric equilibria with flow. The situation with flow becomes more realistic when one realizes the existence of equilibrium flows both in toroidal and poloidal directions in tokamaks following momentum deposition through heating by neutral beam injection . With equilibrium flows, the resultant governing differential equation does not remain always elliptic . The investigation of a general MHD equilibrium becomes much more complicated when one tries to include the effects of other important factors, say of viscous stress tensor. Recently Ren et al. have studied the deformation of magnetic island by including the effect of sheared flow and viscosity into an ideal two-dimensional MHD equilibrium configuration. However, there is an element of inconsistency, whether an ideal equilibrium is realistic . Heuristically, one ignores the resistivity in the Ohm’s law while calculating the equilibrium, but then a resistive stability analysis based on a stationary equilibrium remains questionable as long as the field diffusion is not taken into account . Montgomery et al. have investigated the problem on non-ideal static axisymmetric equilibrium. There have also been attempts to calculate resistive axisymmetric equilibrium with only toroidal flow . It has been further argued that tokamak equilibrium flow is either purely toroidal or the poloidal component is small and quickly damped by magnetic pumping . So there is a natural tendency to exclude the poloidal flow while calculating an equilibrium. But when one considers finite conductivity with purely equilibrium toroidal flow, the conductivity (hence the resistivity) becomes a function of space. In general, the resistivity is not a flux function irrespective of equilibrium flow. In this report, we ask the very pertinent question, whether the situation changes in presence of poloidal flow. As we show that a uniform resistivity profile is consistent in presence of poloidal flow, whereas it has been shown that a scalar pressure equilibrium can not have uniform resistivity . Further we show that in a field-reversed (FRC) axisymmetric configuration with no toroidal magnetic field, both toroidal and poloidal equilibrium flows can coexist with finite resistivity, which is not found to be the case with ideal equilibrium .
We consider the equilibrium resistive MHD equations with plasma flow. The equations are
$`(\rho 𝐯)`$ $`=`$ $`0,`$ (1)
$`𝐁=0,\times 𝐄`$ $`=`$ $`0,\times 𝐁=0,`$ (2)
$`𝐄+𝐯\times 𝐁`$ $`=`$ $`𝐣/\sigma ,`$ (3)
$`\rho (𝐯)𝐯`$ $`=`$ $`𝐣\times 𝐁p,`$ (4)
where the symbols have their usual meanings. We use a right handed cylindrical system $`(r,\theta ,z)`$ with $`z`$ as the axis of symmetry, $`\theta `$ as the toroidal angle, and $`r`$ along the major radius of an axisymmetric device. We assume the plasma flow to be arbitrary (toroidal and poloidal) and axisymmetry is assumed i.e. $`/\theta =0`$. The plasma resistivity $`\eta =\sigma ^1`$ is assumed to be an unspecified function of $`r`$ and $`z`$. We have further assumed here that the equilibrium is maintained in a steady-state through resistive diffusion. The magnetic induction equation allows us to write the magnetic field as
$$𝐁=\frac{1}{r}\psi \times \widehat{𝐞}_\theta +\frac{I}{r}\widehat{𝐞}_\theta ,$$
(5)
where $`\psi `$ is the magnetic flux function which is the azimuthal component of the vector potential $`𝐀`$ and $`I`$ is the current function. Similarly, following the continuity equation, Eq.(1), we can express the plasma equilibrium velocity as
$$𝐯=\frac{1}{\rho r}\phi \times \widehat{𝐞}_\theta +\omega r\widehat{𝐞}_\theta ,$$
(6)
where $`\phi `$ is the velocity stream function and $`\omega =v_\theta /r`$ is the toroidal angular velocity. We also assume that the flow is incompressible i.e. $`𝐯=0`$.
Because the flow is now in both toroidal and poloidal direction, the poloidal component of current, $`𝐣_p`$ need not vanish. In general the current can be expressed as
$$𝐣=\frac{1}{r}\mathrm{\Delta }^{}\psi \widehat{𝐞}_\theta +\frac{1}{r}I\times \widehat{𝐞}_\theta ,$$
(7)
where $`\mathrm{\Delta }^{}`$ is the elliptic operator defined by $`\mathrm{\Delta }^{}\psi =r^2\left(\frac{1}{r^2}\psi \right).`$ Taking curl of the Ohm’s law Eq.(3), we have
$$\times (𝐯\times 𝐁)=\times (𝐣/\sigma )$$
(8)
with
$$𝐯\times 𝐁=\frac{1}{\rho r^2}\phi \times \psi \frac{I}{\rho r^2}\phi +\omega \psi .$$
(9)
The $`\widehat{𝐞}_\theta `$ component of Eq.(8) can be now written as
$`\widehat{𝐞}_\theta \omega \times \psi \widehat{𝐞}_\theta \left({\displaystyle \frac{I}{\rho r^2}}\right)\times \phi `$ (10)
$`={\displaystyle \frac{1}{\sigma r}}\left({\displaystyle \frac{1}{\sigma }}\sigma I+{\displaystyle \frac{2}{r}}{\displaystyle \frac{I}{r}}^2I\right).`$ (11)
We now invoke the large aspect ratio expansion and assume that toroidal magnetic field, to the first approximation can be written $`B_\theta B_0r_0/r`$. Here, $`B_0`$ is the value of the toroidal magnetic field at center of the cylindrical cross section of the torus at distance $`r_0`$ from the axis of symmetry. To this effect we have the current function $`IB_0r_0=\mathrm{const}`$. Under these assumption, the above equation can be written as
$`\widehat{𝐞}_\theta \omega \times \psi `$ $`=`$ $`{\displaystyle \frac{I}{\rho ^2r^2}}\widehat{𝐞}_\theta \rho \times \phi `$ (13)
$`{\displaystyle \frac{2I}{\rho r^3}}\widehat{𝐞}_\theta r\times \phi .`$
The first term on the right hand side of the above equation vanishes by virtue of the incompressibility condition and the continuity equation. We neglect the second term because of its $`1/r^3`$ dependence and find that the toroidal angular velocity $`v_\theta /r=\omega (\psi )`$ becomes a surface quantity. This also further means that $`𝐣=j_\theta \widehat{𝐞}_\theta `$ with $`𝐣_p=0`$. We, however, note that the condition $`\omega \omega (\psi )`$ is identically satisfied in a field-reversed configuration (FRC) where $`I=0`$ and large aspect ratio approximation is not required. It can be noted here that without any approximation, $`\omega `$ becomes a flux function when one considers ideal equilibrium or resistive equilibrium with only toroidal flow .
Now, we consider the momentum equation Eq.(4) and its $`\widehat{𝐞}_\theta `$ component. With the above approximations, we can write Eq.(4) as
$`j_\theta \widehat{𝐞}_\theta \times 𝐁_p`$ $`=`$ $`P+\rho \left[{\displaystyle \frac{1}{2\rho ^2r^2}}(\phi )^2\right]`$ (16)
$`\left({\displaystyle \frac{1}{\rho r^2}}\phi \right)\phi \omega ^{}(\phi \times \psi )`$
$`{\displaystyle \frac{2\omega }{r}}{\displaystyle \frac{\phi }{z}}\widehat{𝐞}_\theta \omega ^2rr`$
In the above equation $`𝐁_p`$ is the poloidal component of the magnetic field and $`(^{})`$ denotes derivative with respect to $`\psi `$. Taking the $`\widehat{𝐞}_\theta `$ component of the above equation, we have,
$$\widehat{𝐞}_\theta \phi \times (\omega r^2)=0,$$
(17)
which means that $`\phi \phi (\omega r^2)`$. We take the simplest situation of $`\phi \omega r^2`$ which yields another surface quantity, $`\phi /r^2=\zeta (\psi )`$. However, it is important to note that, $`\zeta (\psi )`$ is not an arbitrary function in the sense that it is proportional to the toroidal velocity $`\omega (\psi )`$ i.e. the toroidal and poloidal flows are no longer independent. Physically, one can understand this by noting that finite resistivity allows plasma motion across the the flux surfaces. Equivalently, toroidal flow, in a resistive axisymmetric plasma, is always associated with poloidal flow.
Because of equilibrium flow, however, plasma pressure $`p`$ is no longer a flux function now. Taking the $`𝐁_p`$ component of the momentum equation (Eq.(16)), we have
$`𝐁_p\left[{\displaystyle \frac{p}{\rho }}+\left\{{\displaystyle \frac{1}{2\rho ^2r^2}}(\phi )^2{\displaystyle \frac{1}{2}}\omega ^2r^2\right\}\right]`$ (18)
$`={\displaystyle \frac{1}{\rho }}\left({\displaystyle \frac{1}{\rho r^2}}(\zeta r^2)\right)𝐁_p(\zeta r^2).`$ (19)
Depending upon the equation of state, now, several options are possible. However, we note that density, in general, is not a flux function in presence of arbitrary plasma flow. this can be easily seen from the equation of continuity Eq.(1), after applying the incompressibility condition,
$$\widehat{𝐞}_\theta \phi \times \rho =0.$$
(20)
It can be seen from the above expression that $`\rho `$ is not a surface quantity. We note here that axisymmetric equilibria with incompressible equilibrium flows are generally associated with constant density magnetic surfaces . One is also free to choose density as a flux function in case of resistive axisymmetric equilibrium with only incompressible toroidal flow .
Taking the $`\widehat{𝐞}_\theta `$ component of Ohm’s law Eq.(3) along with Eq.(9), we have an expression for plasma conductivity,
$$\sigma \left(E_0r_0+\frac{2}{\rho }r\zeta B_r\right)+\mathrm{\Delta }^{}\psi =0,$$
(21)
where $`E_0`$ is the longitudinal externally applied electric field at major radius $`r=r_0`$. We immediately see from the above expression that conductivity, in general, is a space dependent quantity.
In what follows, we shall consider two cases with (i) uniform and constant density and (ii) a nonuniform density. In the second case, we consider isentropic magnetic surfaces. We now assume that plasma density is uniform and constant i.e. $`\rho =\mathrm{const}.`$ and normalize our equations to $`\rho =1`$. We can now write Eq.(19) as,
$`𝐁_p\left[p+{\displaystyle \frac{1}{2r^2}}(\phi )^2{\displaystyle \frac{1}{2}}\omega ^2r^2\right]`$ (22)
$`={\displaystyle \frac{1}{r^2}}\mathrm{\Delta }^{}(\zeta r^2)𝐁_p(\zeta r^2).`$ (23)
Integration of the above equation yields the equivalent Bernoulli’s equation,
$`p+{\displaystyle \frac{(\phi )^2}{2r^2}}{\displaystyle \frac{\omega ^2r^2}{2}}`$ $`=`$ $`{\displaystyle \frac{dl}{B_p}\frac{1}{r^2}\mathrm{\Delta }^{}(\zeta r^2)𝐁_p(\zeta r^2)}`$ (25)
$`+\chi (\psi ),`$
where the integration is along a magnetic field line and $`\chi (\psi )`$ is an arbitrary surface quantity. The solubility condition further requires that
$$\frac{dl}{B_p}\frac{1}{r^2}\mathrm{\Delta }^{}(\zeta r^2)𝐁_p(\zeta r^2)=0.$$
(26)
We further assume that part of the pressure gradient that varies within a magnetic flux tube has no $`\psi `$ component i.e.
$$\psi \frac{dl}{B_p}\frac{1}{r^2}\mathrm{\Delta }^{}(\zeta r^2)𝐁_p(\zeta r^2)=0.$$
(27)
Together with Eq.(25) and the above assumption, the $`\psi `$ component of the momentum equation yields the equivalent Grad-Shafranov equation,
$$\mathrm{\Delta }^{}\psi +r^2(\chi ^{}+\omega ^{}\omega ^2r)+\frac{\mathrm{\Delta }^{}(\zeta r^2)\psi (\zeta r^2)}{|\psi |^2}=0,$$
(28)
with two arbitrary flux functions $`\chi (\psi )`$ and $`\omega (\psi )`$. The primes refer derivative with respect to $`\psi `$.
We now consider the second case where we consider a nonuniform density. With incompressible flow, magnetic surfaces with constant entropy is quite a reasonable approximation in ideal MHD . However, considering long resistive diffusion time, the right hand side of Eq.(3) can be neglected and we can continue to proceed with isentropic magnetic surfaces . The equation of state can now be written as, $`p=S\rho ^\gamma `$, where, $`S(\psi )`$ is the entropy which is a flux function and $`\gamma `$ is the ratio of specific heats. We now write $`𝐁_pp/\rho `$ as $`𝐁_p[\gamma S\rho ^{\gamma 1}/(\gamma 1)]`$, so that equivalent Bernoulli’s equation can be written as,
$`\mathrm{\Theta }(\psi )+{\displaystyle \frac{dl}{B_p}\frac{1}{\rho }\left(\frac{1}{\rho r^2}(\zeta r^2)\right)𝐁_p(\zeta r^2)}`$ (29)
$`={\displaystyle \frac{\gamma }{\gamma 1}}S\rho ^{\gamma 1}+{\displaystyle \frac{1}{2\rho ^2r^2}}(\phi )^2{\displaystyle \frac{1}{2}}\omega ^2r^2,`$ (30)
where $`\mathrm{\Theta }(\psi )`$ is arbitrary. As we have assumed previously, it requires a solubility condition and the equivalent to the assumption (27). We can then continue to write the equivalent Grad-Shafranov equation by taking the $`\psi `$ component of the momentum equation and applying the Bernoulli’s law Eq.(30),
$`\mathrm{\Delta }^{}\psi +r^2\left(\mathrm{\Theta }^{}+\omega ^{}\omega ^2rS^{}{\displaystyle \frac{\rho ^{\gamma 1}}{\gamma 1}}\right)`$ (31)
$`={\displaystyle \frac{r^2}{\rho |\psi |^2}}\left[{\displaystyle \frac{1}{\rho r^2}}(\zeta r^2)\right]\psi (\zeta r^2).`$ (32)
In the above equation we have four arbitrary surface quantities i.e. $`\mathrm{\Theta }(\psi )`$, $`\omega (\psi )`$, and $`S(\psi )`$ and the primes denote derivative with respect to $`\psi `$.
We have derived the differential equations, equivalent to the Grad-Shafranov equation, for resistive axisymmetric plasma with arbitrary equilibrium flows. These equilibrium equations Eqs.(28, 32) have to be solved subject to conductivity constraint Eq.(21). Further, in a field-reversed configuration (FRC) with no toroidal magnetic field, it can be seen from Eq.(17) that both poloidal and toroidal flow can coexist.
We now show that a uniform conductivity profile is consistent with resistive axisymmetric equilibria with arbitrary flow. A simple examination of Eq.(4), though reveals that uniform conductivity may be possible with scalar pressure equilibrium in presence of flow, it however provides no easier way of proving it. We note that the usual procedure for solving Eqs.(28) and (32) requires specifying a priori dependence of the respective arbitrary functions on $`\psi `$. However, in the presence of finite resistivity, the resistivity constraint Eq.(21) can be used to solve for $`\psi `$, which is uniquely determined if the right hand side of Eq.(21) is specified . It should be noted here with caution whether the resultant solution for $`\psi `$ corresponds to realistic profiles for other physical quantities such as pressure, density, velocity etc. However, our sole aim, here, is to demonstrate the existence of a solution consistent with uniform resistivity in presence of flows.
From Eq.(20) we know that $`\rho \rho (\phi )`$, and assume that $`\rho \phi `$. We now assume that conductivity is uniform in space so that the resulting Eq.(21) can be written as,
$$\mathrm{\Delta }^{}\psi +\frac{\alpha }{r^2}\frac{\psi }{z}=\beta ,$$
(33)
where $`\alpha `$ and $`\beta `$ are arbitrary constants. It is worthwhile mentioning at this point that Eq.(33) can not be used in case of very small resistivity. In the limit of vanishing resistivity (large $`\beta `$ in the above equation), the solution of Eq.(33) contains short scale spatial dependence (boundary layers), not present in case of ideal equilibrium and may lead to unphysical results.
Note that the above equation is a elliptic equation and can be treated as boundary value problem. Following Zheng et al. , we assume a solution of the form
$$\psi _h(r,z)=\underset{n=0,1,2,\mathrm{}}{}f_n(r)z^n$$
(34)
for the homogeneous part of Eq.(33). We however retain the odd terms in the summation to take care of the asymmetric-term in Eq.(33). For simplicity we assume that $`f_n(r)=0`$ for $`n3`$, which, however, can be extended up to any number of terms if required, about which we shall make a comment later. Substituting Eq.(34) in the equivalent homogeneous equation for Eq.(33) we can solve for the functions $`f_{0,1,2}(r)`$. The homogeneous solution of Eq.(33) is then given by,
$`\psi _h(r,z)`$ $`=`$ $`a_1r^2\{4r^2+16z^2\alpha ^2[4(\mathrm{ln}r)^24\mathrm{ln}r+2`$ (38)
$`r^2]+8\alpha z(r^22\mathrm{ln}r)\}+a_2[2r^2(2\mathrm{ln}r`$
$`\mathrm{\hspace{0.17em}1})\alpha ^2\mathrm{ln}r(\mathrm{ln}r+1)+4\alpha z\mathrm{ln}r+4z^2]`$
$`+a_3r^2[\alpha (2\mathrm{ln}r1)+4z]+a_4r^2,`$
where $`a_i`$s are arbitrary constants to be determined from the boundary conditions. A particular solution of Eq.(33) is $`\psi _p=\beta r^2(2\mathrm{ln}r1)/4`$. So the complete solution of Eq.(33) is
$$\psi =\psi _h+\psi _p,$$
(39)
which can be verified by direct substitution. For a conducting circular boundary of a toroidal axisymmetric device, the constant flux ($`\psi `$) contours are shown in Fig.1 (a) which shows a scaler pressure equilibrium. The solution for $`\sigma r^2`$ is shown in Fig.1 (b). Note that $`\sigma r^2`$ is the only possible solution for resistive axisymmetric equilibrium without flow .
In principle the expansion in Eq.(38) should be retained with a large number of terms which will result a equally large number of arbitrary constants for the solution in $`\psi `$. These constants can then be used to shape any arbitrarily shaped plasma boundary.
In passing, we would like to note that resistive field diffusion ($`B/t0`$) is intrinsically involved with non-stationary equilibria ($`v0`$). However, a series of ideal quasi-stationary equilibrium states can be built up with $`B/t=0`$ in which, the effect of finite resistivity is only to slowly evolve the equilibrium in a diffusive time scale .
It is a pleasure to thank A. Sen for the kind hospitality at IPR where part of this work has been completed.
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# The Interstellar Medium of Young Stellar Clusters from the Mid-Infrared Point of View
## 1. Introduction
The subject of stellar cluster formation history and environment has made great headway lately, with the high resolution and sensitivity currently available at optical and near infrared (NIR) wavelengths. In principle, the mid-infrared (MIR) wavelength regime should provide numerous advantages for such studies, since this wavelength range is relatively extinction free (A<sub>15μ</sub><sub>m</sub> $``$ 5% A<sub>J</sub>) and contains diagnostic ionic lines to probe HII regions. In addition, hot dust emission provides us with another link to the ultraviolet starlight that has been absorbed and reemitted by the nearby grains. Our knowledge of the MIR wavelength window has been limited by the low spatial and spectral resolution provided by the IRAS satellite, and has remained rather sketchy when it comes to detailed studies of the ISM of individual galaxies. The Infrared Space Observatory (ISO; Kessler et al. 1996) has been a recent turning point in this effort, providing high spectral and spatial resolution and unprecedented sensitivity in the MIR through the far infrared (FIR). We have incorporated these MIR and FIR observations in a study of the energy redistribution in starburst galaxies, with the aim of understanding the impact of the star formation on the surrounding gas and dust.
The main limitation in MIR star cluster studies remains the spatial resolution, despite the great improvement over previous instrumentation provided by ISOCAM ($``$ 6<sup>′′</sup> at 15$`\mu `$m; Cesarsky, C.J. et al. 1996). ISOCAM resolves about 600 pc at 20 Mpc, the distance of the closest massive merging system, the Antennae. However, even with this limitation, we are able to draw noteworthy conclusions from the MIR, from unique MIR diagnostics. Here, I concentrate primarily on results of the nearer dwarf galaxies, since impacts of the massive clusters on the global dust and gas environment are very pronounced in these relatively small objects.
## 2. What Do MIR Wavelengths Trace?
Figure 1 shows ISOCAM 5-17 $`\mu `$m spectra for the three dwarf galaxies II Zw40, NGC 1140 and NGC 1569, with metallicities of $`{}_{}{}^{1}/_{7}^{}`$ to $`{}_{}{}^{1}/_{3}^{}`$ Z (Madden et al. 2000), along with the spectrum of the notoriously metal poor SBS 0335-052 ($`{}_{}{}^{1}/_{40}^{}`$ Z; Thuan, Sauvage, & Madden 1999). All these spectra show obvious MIR signatures of massive stars.
As often seen in starburst galaxies, the MIR spectra are dominated by steeply rising continua longward of $``$ 10 $`\mu `$m. Thermal emission from hot small grains with mean temperatures of the order of hundreds of Kelvin are responsible for the MIR continuum emission.
The unidentified infrared bands (UIBs) at 6.2, 7.7, 8.6, 11.3 and 12.6 $`\mu `$m, have been attributed to aromatic hydrocarbon particles undergoing stochastic temperature fluctuations (i.e, PAHs: Léger & Puget 1984; Allamendola, Tielens, & Barker 1989; coal grains: Papoular, Reynaud, & Nenner 1991). They are observed to peak in the photodissociation (PDR) zones around HII regions, but are destroyed deep within the HII regions themselves (Verstraete et al. 1996; Cesarsky, D. et al. 1996; Tran 1998). While the UIBs are not obvious in the spectra of II Zw40 and SBS 0335-052, and are only very weakly present NGC 1569, they can be distinguished in the spectrum of NGC 1140. Several ground state fine structure nebular lines are present also in 3 of the spectra, the most prominent being 15.6 $`\mu `$m \[NeIII\] (ionisation potential $``$ 41 eV) and 10.5 $`\mu `$m \[SIV\] ($``$ 35 eV). Weaker, lower energy lines may also present, such as the 8.9 $`\mu `$m \[ArIII\] line and the \[NeII\] 12.8 $`\mu `$m line, which can be blended with the 12.6 $`\mu `$m UIB.
While all of these spectra are very different from one another, all differ significantly from those of normal metallicity starburst galaxies. Normal starburst galaxies show prominent UIBs, in contrast to AGNs, which are devoid of UIBs (e.g. Roche et al. 1991; Dudley 1999, Laurent et al. 2000; Sturm et al. 2000). When compared to spectra characteristic of PDRs and HII regions (e.g. M17, Cesarsky, D. et al. 1996; Verstraete et al. 1996), II Zw40 is remarkably similar to that of an HII region. In contrast, NGC 1140, which has a very flat continuum yet a very strong \[NeIII\] line, does have a more obvious contribution from PDR regions in its spectra. Note that the MIR spectrum of N66, the most prominent HII region in the SMC, also shows a scarcity of UIBs in the vicinity of the most massive central cluster (Contursi et al. 2000), as does the low metallicity source NGC 5253 (Crowther et al. 1999).
In some starburst galaxies, amorphous silicate is seen in absorption centered at 9 and 18 $`\mu `$m (Roche et al. 1991; Dudley 1999; Laurent et al. 2000). We can fit the MIR region of the II Zw40 spectrum with a blackbody at 193 K and an absorption equivalent to A<sub>v</sub> $``$ 4. Dust temperatures derived assuming blackbodies should be interpreted with care, since the dust emitting in the MIR is expected to be undergoing stochastic processes rather than being in thermal equilibrium with the radiation field. The amount of absorption in II Zw40 (A<sub>v</sub> $``$ 4) has yet to be confirmed. In SBS 0335-052, A<sub>v</sub> $``$ 20 has been deduced from the absorption in the ISOCAM MIR spectra (Fig. 1). The presence of a significant amount of dust at a metallicity as low as Z/40 is surprising, especially since star formation in SBS 0335-052 began as recently as 100 Myr ago (Papaderos et al. 1998; Thuan, Izotov, & Foltz 1999). Such high extinction implies that the current star formation rate, hidden by dust, might be underestimated by at least 50% (Thuan, Sauvage & Madden 1999).
## 3. Effects of the Massive Star Formation on the Gas
As a consequence of the smaller dust abundance of most dwarf galaxies, the ISM throughout these galaxies is affected by the hard radiation field of the massive stellar clusters. All star forming dwarf galaxies contain evidence for Wolf-Rayet stars (Schaerer, Contini, & Pindao 1999) and super star clusters have been detected in NGC 1140 (Hunter, O’Connell, & Gallagher 1994), NGC 1569 (O’Connell, Gallagher, & Hunter 1994) and SBS 0335-052 (Thuan, Izotov, & Lipovetsky 1997). Their harsh radiation fields, which more readily permeate the ISM compared those in solar metallicity environments, are capable of destroying the UIB carriers, for example, over very extensive spatial areas. The effect of the pervasive radiation field can be witnessed in NGC 1569 (Fig. 2), where photodissociation occurs on global scales. Violent activity is revealed by the H$`\alpha `$ distribution (Waller 1991; Martin 1998) and the 15.8 $`\mu `$m \[NeIII\] emission, with giant streamers suspected to originate from the energetic winds of the super star clusters A & B, (black stars in Fig. 2). The UIB, \[SIV\] and \[NeIII\] emission seems to avoid the super star clusters, which blow out much of the gas and dust on relatively short time scales. This effect is also seen in the CO (Taylor et al. 1999), HI (Israel & van Driel 1990) and H$`\alpha `$ (Waller 1991) distributions. Likewise we see the destruction of the UIBs in the beam-averaged spectrum of the entire galaxies II Zw40 and SBS 0335-052 (the available spatial resolution prevents us from seeing more detail within these galaxies in the MIR).
## 4. Modelling the Spectral Energy Distribution
We have compiled broad band data from the literature for II Zw40, NGC 1569 and NGC 1140 and have combined these with our MIR data to construct appropriate stellar spectral energy distributions (SEDs). In so doing, we fit the observed optical and NIR data with population synthesis models of PEGASE (Fioc & Rocca-Volmerange 1997), taking into account the constraints of the MIR line emission by modelling the corresponding photoionisation with CLOUDY (Ferland 1996). After briefly describing the results of this process, we discuss the results of the use of the reconstructed stellar SEDs as input to our dust model.
### 4.1. Combined stellar evolution and photoionisation model results
When assuming instantaneous star formation, a metallicity Z/5 and a Salpeter IMF (with upper and lower mass cut-offs of 0.1 and 120 solar masses), we find solutions to the observed broad band colours for a variety of ages and ionisation parameters. The ISOCAM MIR observations provide the diagnostic lines of neon, sulphur and argon, that have been recently addressed e.g. by Lutz et al. (1998), Crowther et al. (1999), Schaerer & Stasińska (1999) and Genzel et al. (1998).
For example, the \[NeIII\]/\[NeII\] ratio is a measure of T<sub>eff</sub>, the hardness of the radiation field, and therefore traces the massive stellar population. For the dwarf galaxies, we find \[NeIII\]/\[NeII\] ratios in the range of 5 to 10 - much higher values than those of normal metallicity galaxies ($``$1; Thornley et al. 2000). The extreme values of the \[NeIII\]/\[NeII\] ratios are related to the low metallicities of the systems: the T<sub>eff</sub> of the stars increases as the metallicity decreases for a specific stellar age. High ratios of \[NeIII\]/\[NeII\] and the prominent \[SIV\] in these spectra limit the age of the present star formation to $`<`$ 5 Myr. Beyond this age, the massive stars have died and the \[NeIII\]/\[NeII\] ratio drops dramatically. The high excitation 24.9 $`\mu `$m \[OIV\] line, covered by the ISO SWS data, is observed in some dwarf galaxies (Lutz et al. 1998) and has been attributed to the presence of Wolf-Rayet stars (Schaerer & Stasińska 1999).
For NGC 1569, NGC 1140 and II Zw40, we construct composite stellar SEDs that require 70 to 95 % of the stellar mass to be provided by an ’older’ population with ages between about 10 and 30 Myr, with the remaining 5% to 30% corresponding to a very young population ($`<`$ 5 Myr). Observational evidence for the presence of Wolf-Rayet stars corroborates the existence of this very young stellar population (Vacca & Conti 1992). The broad band optical and NIR data alone reveal predominantly the older population in our apertures. Fig. 3 shows an example of the resultant composite SED for II Zw 40, including the extreme ultraviolet (EUV) radiation that the young, massive stellar population traces.
### 4.2. Effects of the Massive Star Formation on the Dust
We use the modelled stellar spectra of II Zw40, NGC 1569 and NGC 1140 as input to a dust model to study the effects of this radiation field on the dust properties. This is an important step since dust plays a major role in influencing the chemical and physical state of the ISM. We use the Désert, Boulanger, & Puget (1990) model to fit the various dust components emitting in the MIR and the FIR. This model calculates the IR emission from large silicate grains (BGs), very small amorphous carbon grains (VSGs), and stochastically heated polycyclic aromatic hydrocarbons (PAHs), for various grain size distributions.
In these three galaxies the MIR spectrum is clearly dominated by emission from VSGs with very little PAH emission. The BG component dominates the overall dust emission with mass fractions ranging from 93% to 99%, while the PAH mass fraction is relatively insignificant — 5 orders of magnitude lower. The model gives a PAH/VSG mass ratio of $`2\mathrm{3\hspace{0.17em}10}^4`$ for NGC 1569 and II Zw 40 and 10 times this value for NGC 1140. For comparison, the Désert et al. model applied to the Galactic cirrus gives a PAH/VSG mass ratio $`1`$. Thus, even compared to the VSG population, we find an insignificant mass fraction of PAHs, reflecting the fact that the PAHs are destroyed throughout the entire galaxies, as a result of the hard radiation fields originating from the few massive stellar clusters. This in an important result, since PAHs are thought to be the primary particles responsible for the photoelectric heating process (Bakes & Tielens 1994) and are incorporated in PDR models (Kaufman et al. 1999). Our preliminary results, while not statistically robust at this stage, suggest that even in the absence of PAHs, the photoelectric effect is efficient, as both II Zw40 and NGC 1569 are relatively prominent \[CII\] sources among the galaxies surveyed (Jones et al. 1997). On the contrary, in NGC 1140, where PAHs are more obvious in the MIR spectra (Fig. 1), we do not detect \[CII\]. VSGs with derived sizes of $``$40 to 300 Å, which are very abundant relative to the PAHs in NGC 1569 and II Zw40 but less so in NGC 1140, may therefore be the more efficient sources of photoelectric gas heating in these environments, rather than PAHs.
This scenario is in contrast to normal metallicity galaxies, where the PAHs are prominently observed, and the effects of the numberous massive stellar activity are much more local rather than global.
## 5. Summary
MIR ISOCAM spectroscopy provides details of ionic lines, UIBs and the distribution of small hot grain emission in dwarf galaxies. The strong MIR \[NeIII\]/\[NeII\] ratios are signatures of the hard radiation fields and indicate the presence of clusters of young massive stars in dwarf galaxies. Because of the increase in T<sub>eff</sub> in low metallicity environments, this ratio is enhanced in dwarf galaxies to at least 5 to 10 times that observed in normal metallicity galaxies. The penetrating radiation field also affects the dust components, destroying the UIBs in some dwarf galaxies on global scales, as is evident in the MIR spectra and in the dust modeling. This dramatic global effect of the massive stellar population in dwarf galaxies, due to the decrease in attenuation of the UV flux, is not apparant in normal metallicity galaxies, where these effects are experienced much more locally.
#### Acknowledgments.
It is a pleasure to acknowledge my collaborators D. Ragaigne and A. Jones. I wish to also thank W. Waller for his H$`\alpha `$ image of NGC 1569.
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# Is the Multichannel Kondo Model Appropriate to Describe the Single Electron Transistor?
## I Introduction
The single electron box (SEB) and the single electron transistor (SET) are basic elements of mesoscopic devices and have been studied extensively. Both consist of a single small metallic or semiconducting box connected to one (single electron box, or SEB) or two (single electron transistor, or SET) leads. Additionally, in the SET a gate electrode is attached to the box to control the actual charge on the dot (see Fig. 1).
The electrostatic energy of the box is well-described by the classical expression
$$H_C=\frac{e^2}{2C}\left(n_{\mathrm{box}}\frac{V_gC_g}{e}\right)^2,$$
(1)
where $`C`$ denotes the capacitance of the island, $`C_g`$ is the gate capacitance, $`e`$ is the electric charge, $`V_g`$ stands for the gate voltage, and $`n_{\mathrm{box}}`$ is the number of extra electrons on the island. For box sizes in the $`0.1\mu \mathrm{m}`$ range the capacitance $`C`$ of the box can be small enough so that the charging energy $`E_C=e^2/2C`$ associated with putting an extra electron on the island can safely be around $`\mathrm{mV}`$ range. Therefore, unless $`N_g=\frac{V_gC_g}{e}`$ is a half-integer, it costs a finite energy to charge the island, and at low enough temperatures the number of electrons on the island becomes quantized and a Coulomb blockade develops provided the quantum fluctuations are not too strong. In this Coulomb blockade regime the transport through the island is suppressed.
The situation is dramatically different for dimensionless gate voltages $`N_g=\frac{V_gC_g}{e}n+1/2`$. For $`N_g=n+1/2`$ the two states $`n_{\mathrm{box}}=n`$ and $`n_{\mathrm{box}}=n+1`$ become degenerate, and quantum fluctuations between the island and the leads become important. Assuming that the mean level spacing on the island is smaller than any energy scale (temperature, $`E_C`$, etc.) two scenarios have been suggested: (a) It has been proved by Matveev that if the leads are connected to the island via a single conduction mode then — close to the degeneracy point and at low enough temperature — the physics of the SET (SEB) becomes identical to that of the two-channel Kondo model. Indeed, the fingerprints of the two-channel Kondo behavior have been observed recently on semiconducting single electron transistors. (b) In the opposite limit one assumes that the tunneling to the island happens through an infinite number of identical modes. This model has been applied very successfully for the description of metallic islands. The predictions of this infinite channel model are, however, very different from those of the two-channel Kondo model: the conductance of the SET, for instance, scales to zero as $`T`$ in the two-channel Kondo picture, while it is proportional to $`1/\mathrm{ln}(E_C/T)`$ in the infinite channel scenario.
The purpose of the present paper is to treat the general case of finite conduction modes in the lead and to reconcile the apparent contradiction between the two pictures above. We show that both models capture the physical properties of the SEB (SET), however, they are appropriate in very different regimes. Carrying out a renormalization group analysis we show that there exists a crossover energy $`T^{}`$. Above $`T^{}`$, even for $`N20`$ tunneling modes the system is well characterized by the conductance of the tunnel junctions and the SEB (SET) is satisfactorily described by the $`2N`$-channel model of Ref. . Nevertheless, for small mode numbers pronounced deviations occur, and similar deviations appear for larger values of $`N`$ in the presence of pinholes in the junction, which offer a plausible explanation to the deviations observed in Ref. .
Below $`T^{}`$, on the other hand, the detailed structure of the tunneling matrix becomes important: At very low $`T`$ only a single conductance mode dominates the physics and a two-channel Kondo effect develops. Unfortunately, in most situations $`T^{}`$ (and thus the Kondo temperature $`T_KT^{}`$) turns out to be extremely small, and the two-channel Kondo physics cannot be observed. In fact, very special experimental setups are needed to observe the two-channel Kondo behavior, as we shall discuss it in detail in our concluding section.
The paper is organized as follows: In Sec. II we describe the models applied. Secs. III and IV are devoted to the analysis of the single electron box and the single electron transistor, respectively. Finally, in Sec. V we discuss the possibility of experimental observations of the low-energy Kondo-like behavior and summarize our conclusions.
## II The Models
### A Hamiltonian of the SEB
For the sake of simplicity, let us first concentrate on the single electron box and generalize our results to the SET later on. Usually, the lead is described by means of $`N`$ independent non-interacting one-dimensional electron modes:
$$H_{\mathrm{lead}}=\underset{n=1}{\overset{N}{}}\underset{ϵ,\sigma }{}ϵc_{ϵn\sigma ;\mathrm{lead}}^{}c_{ϵn\sigma ;\mathrm{lead}},$$
(2)
where $`c_{ϵn\sigma ;\mathrm{lead}}^{}`$ crates an electron on the box with spin $`\sigma `$, mode index $`m`$ and energy $`ϵ`$. (Note that to avoid confusion, we do not follow the usual terminology and use deliberately the expression conduction mode instead of the wording ’conductance channels’, more frequent in the literature.)
In the present work we assume that the level spacing $`\mathrm{\Delta }`$ at the island is much smaller than any energy scale in the problem, and therefore these discrete levels may be represented as a single particle continuum on the island. This assumption is crucial to obtain the Kondo physics discussed in this paper, since the level spacing provides an infrared cut-off which ultimately kills the logarithmic singularities and the Kondo effect. Based on these assumptions we express the Hamiltonian of the island as
$`H_{\mathrm{box}}`$ $`=`$ $`{\displaystyle \underset{m=1}{\overset{M}{}}}{\displaystyle \underset{ϵ}{}}ϵc_{ϵm\sigma ;\mathrm{box}}^{}c_{ϵm\sigma ;\mathrm{box}}`$ (3)
$`+`$ $`E_C\left(n_{\mathrm{box}}V_gC_g/e\right)^2.`$ (4)
In Eq. (4) we assumed $`M`$ independent modes on the box, and defined the number $`n_{\mathrm{box}}`$ of electrons on the island as
$$n_{\mathrm{box}}=\underset{m=1}{\overset{M}{}}\underset{ϵ,\sigma }{}:c_{ϵm\sigma ;\mathrm{box}}^{}c_{ϵm\sigma ;\mathrm{box}}:,$$
(5)
where the symbol $`:\mathrm{}:`$ denotes normal ordering.
As usually, in Eq. (4) we implicitly used the assumption that the collective charge excitations decouple from the single particle excitations and that the electron-electron interaction can be fully taken into account by the classical Coulomb interaction term. The validity of this approximation relies heavily on the fact the collective charge excitations relax extremely fast compared to all other time scales involved.
The coupling of the box to the lead is described by a standard tunneling Hamiltonian:
$$H_{\mathrm{tun}}=\underset{n=1}{\overset{N}{}}\underset{m=1}{\overset{M}{}}\underset{ϵ,ϵ^{},\sigma }{}\left(T_{mn}c_{ϵm\sigma ;\mathrm{box}}^{}c_{ϵ^{}n\sigma ;\mathrm{lead}}+\text{h.c.}\right),$$
(6)
where we neglected the energy dependence of the elements of the $`M\times N`$ tunneling matrix $`T_{mn}`$
It is very important that the tunneling is diagonal in the spin indices, however, it is generally non-diagonal in the mode indices. As we shall see later, the twofold spin degeneracy is the basic origin of the very low-temperature two-channel Kondo effect. Once magnetic field is applied the symmetry between spin up and spin down conduction electrons is broken, and the system flows to the single channel Kondo fixed point.
The tunneling Hamiltonian of the SET differs only slightly from that of the SEB. In this case there are two leads that are connected to the island. However, as first shown by Averin and Nazarov, at temperatures larger than the level spacing coherent processes connecting the two leads are strongly suppressed. Therefore, one can formally separate from each-other those single particle states on the island which participate in the tunneling from the first and the second lead, respectively. These tunneling processes are then only correlated by the very fast Coulomb interaction which allows for the presence of only one excess electron on the island.
Thus for the SET the effective Hamiltonian of the island becomes:
$`H_{\mathrm{box}}`$ $`=`$ $`{\displaystyle \underset{m=1}{\overset{M}{}}}{\displaystyle \underset{ϵ}{}}ϵ(c_{}^{(1)}{}_{ϵm\sigma ;\mathrm{box}}{}^{}c_{}^{(1)}{}_{ϵm\sigma ;\mathrm{box}}{}^{}`$ (7)
$`+`$ $`c_{}^{(2)}{}_{ϵm\sigma ;\mathrm{box}}{}^{}c_{}^{(2)}{}_{ϵm\sigma ;\mathrm{box}}{}^{})`$ (8)
$`+`$ $`E_C\left(n_{\mathrm{box}}^{(1)}+n_{\mathrm{box}}^{(2)}V_gC_g/e\right)^2,`$ (9)
where the indices $`(1)`$ and $`(2)`$ refer to single particle states participating in the tunneling from the first and second leads, respectively, and the number operators $`n_{\mathrm{box}}^{(1)}`$ and $`n_{\mathrm{box}}^{(2)}`$ are defined similarly to Eq. (5).
The tunneling Hamiltonian of the SET reads:
$$H_{\mathrm{tun}}=\underset{\genfrac{}{}{0pt}{}{f=1,2}{n,m,ϵ,ϵ^{},\sigma }}{}\left(T_{mn}^{(f)}c_{}^{(f)}{}_{ϵm\sigma ;\mathrm{box}}{}^{}c_{}^{(f)}{}_{ϵ^{}n\sigma ;\mathrm{lead}}{}^{}+\text{h.c.}\right),$$
(10)
where the index $`f`$ refers to the two junctions. For the sake of simplicity, we assumed that the number of modes in the two leads ($`N^{(1)}`$ and $`N^{(2)}`$) and the number of tunneling modes on the island ($`M^{(1)}`$ and $`M^{(2)}`$ ) is identical for both junctions($`N^{(1)}=N^{(2)}=N`$ and $`M^{(1)}=M^{(2)}=M`$ ). This simplification does not modify our results because the two tunneling matrices $`T_{mn}^{(1)}`$ and $`T_{mn}^{(2)}`$ are assumed to be completely uncorrelated.
In the following we focus to the vicinity of the degeneracy points, $`VC_G/e=n+1/2`$. As already mentioned in the introduction at these gate voltages the charge states $`n_{\mathrm{box}}=n`$ and $`n_{\mathrm{box}}=n+1`$ become degenerate and quantum fluctuations dominate. For temperatures (energy scales) below the charging energy $`E_C=e^2/2C`$ one can safely project out all the other charging states, represent the two states $`n_{\mathrm{box}}=n+1/2\pm 1/2`$ as two states of a pseudospin $`S_z=\pm 1/2`$, and rewrite the tunneling part of SEB Hamiltonian in the following form:
$`H`$ $`=`$ $`{\displaystyle \underset{ϵ,\sigma ,\alpha }{}}{\displaystyle \underset{n=1}{\overset{N_\alpha }{}}}ϵc_{ϵn\sigma ;\alpha }^{}c_{ϵn\sigma ;\alpha }hS_z`$ (11)
$`+`$ $`{\displaystyle \underset{\genfrac{}{}{0pt}{}{ϵ,ϵ^{},\sigma }{\alpha ,\beta ,n_\alpha ,n_\beta }}{}}(c_{ϵn_\alpha \sigma ;\alpha }^{}T_{n_\alpha ,m_\beta }\sigma _{\alpha \beta }^+S^{}c_{ϵ^{}m_\beta \sigma ;\beta }+\mathrm{h}.\mathrm{c}.),`$ (12)
where the ’orbital pseudospins’ $`\alpha ,\beta =({}_{}{}^{}\mathrm{box}_{}^{},{}_{}{}^{}\mathrm{lead}_{}^{})`$ indicate the position of an electron and couple to $`S`$, $`\sigma ^\pm =\sigma _x\pm i\sigma _y`$ denote Pauli matrices, and the index $`n_\alpha `$ takes values $`n_\alpha =1,..,N_\alpha `$ ($`N_{\mathrm{lead}}=N`$ and $`N_{\mathrm{box}}=M`$). The effective field $`h=e\delta VC_G/C`$ measures the distance from the degeneracy point, with $`\delta V=Ve(n+1/2)/C_G`$.
For the SET Eq. (11) gets modified in that the two leads provide two conduction electron ’channels’ ($`f=1,2`$) coupled to the charge pseudospin:
$$H_{\mathrm{tun}}=\underset{\genfrac{}{}{0pt}{}{f,ϵ,ϵ^{},\sigma }{\alpha ,\beta ,n,m}}{}(c_{}^{(f)}{}_{ϵn\sigma ;\alpha }{}^{}T_{n,m}^{(f)}\sigma _{\alpha \beta }^+S^{}c_{ϵ^{}m\sigma ;\beta }^{(f)}+\mathrm{h}.\mathrm{c}.).$$
(13)
Note that the ’channel’ label $`f`$ has a role essentially different from that of the real spin of the electrons: While there is a full SU(2) symmetry associated to the latter, the former is merely a conserved quantity (corresponding to a $`U(1)`$ symmetry only).
In the formulation above the case $`M=N=1`$ corresponds to Matveev’s two-channel Kondo model, while in the limit $`M=N\mathrm{}`$ and $`T_{nm}=\delta _{nm}T1/\sqrt{N}`$ we recover the infinite channel model mentioned in the Introduction. Obviously, both limits are somewhat specific: In many realistic systems $`M,N>1`$ and the first approximation seems to be inadequate. The second approximation, on the other hand, contains an artificial $`SU(2N)`$ symmetry: There is no reason for the tunneling matrix element to be diagonal in the mode indices at all, and even more to have identical matrix element in each tunneling mode. In fact, any defect, roughness, etc. present in a real junction will produce cross-channel tunneling, and even the simplest models of a perfect tunnel junction with $`N=M`$ give different tunneling eigenvalues for the different tunneling modes. The philosophy behind this second approach is that the only physically relevant parameter is the conductance of the junction and therefore the artificial symmetry introduced has no effect. As we shall see, this philosophy is only partially justified: cross-mode tunneling — breaking this artificial $`SU(2N)`$ symmetry — is in reality a relevant perturbation and leads the system ultimately to the two-channel Kondo fixed point.
## III Perturbative scaling analysis of the SEB
It has been shown longtime ago that the Hamiltonian of Eq. (11) generates logarithmic singularities when perturbation theory in the tunneling amplitude is developed. To deal with these logarithmic singularities one has to sum up the perturbation series up to infinite order. The easiest way to do this is by constructing the renormalization group (RG) equations. Fortunately, this straightforward but rather tedious calculation can be avoided by rewriting the Hamiltonian (11) in the following form:
$$H_{\mathrm{int}}=\underset{i=1}{\overset{3}{}}\underset{\sigma ,ϵ,ϵ^{},r,r^{}}{}V_{rr^{}}^i\tau ^ic_{ϵr\sigma }^+c_{ϵ^{}r^{}\sigma },$$
(14)
and observing that Eq. (14) is formally identical to the Hamiltonian of a non-commutative two-level system. The indices $`r`$ and $`r^{}`$ in Eq. (14) take the values $`r,r^{}=(1,..,M+N)`$, the $`\tau ^i`$’s denote Pauli matrices ($`\tau ^i=2S^i`$), and the $`V_{rr^{}}^i`$’s can be written in a block matrix notation as
$`𝐕^x={\displaystyle \frac{1}{2}}\left(\begin{array}{cc}0& 𝐓\\ 𝐓^{}& 0\end{array}\right),`$ (17)
$`𝐕^y={\displaystyle \frac{1}{2i}}\left(\begin{array}{cc}0& 𝐓\\ 𝐓^+& 0\end{array}\right),`$ (20)
$`𝐕^z=\left(\begin{array}{cc}𝐐& 0\\ 0& \stackrel{~}{𝐐}\end{array}\right),`$ (23)
where we introduced the tensor notation $`T_{mn}𝐓`$. The $`M\times M`$ and $`N\times N`$ Hermitian matrices $`𝐐`$ and $`\stackrel{~}{𝐐}`$ vanish in the bare Hamiltonian, but they are dynamically generated under scaling. They correspond to charging state dependent back scattering off the tunnel junction:
$`H_{\mathrm{back}}`$ $`=`$ $`{\displaystyle \underset{ϵ,ϵ^{},\sigma ,\sigma ^{}}{}}2S_z[{\displaystyle \underset{n,n^{}=1}{\overset{N}{}}}\stackrel{~}{Q}_{nn^{}}c_{ϵn\sigma ;\mathrm{lead}}^{}c_{ϵ^{}n^{}\sigma ^{};\mathrm{lead}}`$ (24)
$`+`$ $`{\displaystyle \underset{m,m^{}=1}{\overset{M}{}}}Q_{mm^{}}c_{ϵm\sigma ;\mathrm{box}}^{}c_{ϵ^{}m^{}\sigma ^{};\mathrm{box}}].`$ (25)
The scaling equations of the two-level system have been first derived in Ref. and its possible fixed points and their stability have been carefully analyzed in Ref. . Using this mapping we can easily construct the RG equations for the SEB:
$`{\displaystyle \frac{d𝐭}{dx}}`$ $`=`$ $`2(\mathrm{𝐭𝐪}\stackrel{~}{𝐪}𝐭)`$ (26)
$``$ $`2𝐭\left[\mathrm{Tr}\{\mathrm{𝐭𝐭}^{}\}+2\mathrm{T}\mathrm{r}\{\mathrm{𝐪𝐪}\}+2\mathrm{T}\mathrm{r}\{\stackrel{~}{𝐪}\stackrel{~}{𝐪}\}\right],`$ (27)
$`{\displaystyle \frac{d\stackrel{~}{𝐪}}{dx}}`$ $`=`$ $`𝐭^{}𝐭4\stackrel{~}{𝐪}\mathrm{Tr}\{\mathrm{𝐭𝐭}^{}\},`$ (28)
$`{\displaystyle \frac{d𝐪}{dx}}`$ $`=`$ $`\mathrm{𝐭𝐭}^{}4𝐪\mathrm{Tr}\{\mathrm{𝐭𝐭}^{}\}.`$ (29)
Here we introduced the scaling variable $`x=\mathrm{ln}(E_c/\mathrm{max}\{T,\omega ,..\})`$ and defined dimensionless couplings as $`t_{mn}(\varrho _m^{\mathrm{box}}\varrho _n^{\mathrm{lead}})^{1/2}T_{mn}`$ (and similarly, $`q_{mm^{}}(\varrho _m^{\mathrm{box}}\varrho _m^{}^{\mathrm{box}})^{1/2}Q_{mm^{}}`$, and $`\stackrel{~}{q}_{nn^{}}(\varrho _n^{\mathrm{lead}}\varrho _n^{}^{\mathrm{lead}})^{1/2}\stackrel{~}{Q}_{nn^{}}`$) with $`\varrho _m^{\mathrm{box}}`$ ($`\varrho _n^{\mathrm{lead}}`$) the density of states at the Fermi energy in mode $`m`$ (mode $`n`$) of the box (lead).
The scaling equation for the effective field $`h`$ can be obtained from that of the splitting in the two-level system problem:
$$\frac{dh}{dx}=4h\mathrm{Tr}\{\mathrm{𝐭𝐭}^{}\}.$$
(30)
These scaling equations are appropriate provided $`g\mathrm{Tr}\{\mathrm{𝐭𝐭}^{}\}<1`$, otherwise the perturbative RG breaks down. They must be solved with the initial condition $`𝐪(x=0)=\stackrel{~}{𝐪}(x=0)=0`$ and $`𝐭(x=0)=𝐭^{\mathrm{bare}}`$. Away from the degeneracy point the “magnetic field”, i.e. the deviation from the degeneracy point is a relevant perturbation, and the scaling must be cut off at an energy scale $`h^{}`$ determined selfconsistently from the condition $`h(x=\mathrm{ln}(E_C/h^{}))=h^{}`$.
Up to logarithmic accuracy, the conductivity can then be expressed in terms of the scaled dimensionless tunneling rate $`𝐭(x=\mathrm{ln}[E_C/\mathrm{max}\{T,h^{}\}])`$ as
$$G(T)=G_0\mathrm{Tr}\{𝐭(x)𝐭^{}(x)\}$$
(31)
where $`G_0=8\pi ^2e^2/h`$ denotes a universal conductance unit.
The analogues of Eqs. (27-29) have been analyzed very carefully in Refs. , where it was shown by means of a systematic $`1/N_s`$ expansion ($`N_s`$ being the spin of the electrons) that they have a unique stable fixed point, identical to the two-channel Kondo fixed point. At this latter two orbital quantum numbers prevail, and the others become irrelevant. In the present context this statement means that there will be a unique ’effective tunneling mode’ in the lead (it is some combination of the original tunneling modes in the lead), and another unique ’tunneling mode’ in the box (also a linear combination of the original modes in the box): The $`T=0`$ effective Hamiltonian of the model at the degeneracy point corresponds to tunneling between these two modes only, and all the other modes can be neglected. This theorem therefore justifies the use of Matveev’s effective model at very low temperatures even if the number of modes in the lead or the box is larger than one.
However, as we show now, the temperature $`T^{}`$, below which the two channel Kondo behavior appears turns out to be extremely small in most cases. To see this let us consider a tunneling junction with a rough tunneling surface, and a dimensionless conductance $`g=G/G_0`$. Obviously, in this case the matrix elements $`t_{mn}`$ scale as $`t_{mn}\sqrt{g}/N`$ so that $`\mathrm{Tr}\{𝐭^{}𝐭\}g`$, and they have random sign or phase relative to each-other. The occupation dependent back scatterings $`𝐪`$ and $`\stackrel{~}{𝐪}`$ are generated by Eqs. (28) and (29), and their typical matrix elements can easily be estimated to be of order $`1/N^{3/2}`$. Therefore, the first two terms in Eq. (27) can be estimated to be smaller than $`g/N^2`$, while the matrix elements of the last term are dominated by the term proportional to $`\mathrm{Tr}\{𝐭^{}𝐭\}`$ and are of order $`g^{3/2}/N`$. Thus for large enough $`NM`$ one can simply substitute $`𝐪=\stackrel{~}{𝐪}=0`$ and the scaling equations reduce to:
$$\frac{d𝐭}{dx}=2𝐭\mathrm{Tr}\{\mathrm{𝐭𝐭}^{}\}$$
(32)
Multiplying this equation by $`𝐭^{}`$ and taking its trace we finally arrive at the scaling equation:
$$\frac{dg}{dx}=4g^2,$$
(33)
which is identical to the one obtained in Ref. using the $`2N`$-channel model. The scaling equation for the effective field can also be expressed in terms of the dimensionless conductance:
$$\frac{dh}{dx}=4gh.$$
(34)
The above two scaling equations can readily be integrated to obtain:
$`g(x)={\displaystyle \frac{g_0}{1+4g_0x}},`$ (35)
$`h(x)={\displaystyle \frac{h_0}{1+4g_0x}}.`$ (36)
Obviously, in this approximation the physics of the SEB is completely characterized by the effective field $`h`$ and the conductance of the junction, and the details about the specific structure of the junction or the tunneling amplitudes are unimportant.
It is easy to estimate the crossover scale $`T^{}`$ below which this approximation breaks down. The scaling towards the two-channel Kondo fixed point is generated by the second order ’coherence’ terms in Eqs. (27-29), while the third order terms tend to reduce all couplings to zero. Therefore the scale $`T^{}`$ is determined by the condition that the second and third order terms in Eqs. (27-29) be of the same order of magnitude, giving $`g1/N`$. Replacing Eq. (35) by its asymptotic form $`g(x)1/4x`$ we finally obtain
$$T^{}E_C\mathrm{exp}\{CN\}(N1),$$
(37)
where $`C`$ is a constant of the order of unity. In view of the experimental values of $`E_C`$, this scale is extremely small for $`N>10`$, which justifies the use of the $`2N`$-channel model in many experimental setups.
In Figs. 2 and 3 we show the typical scaling of $`g(x)`$ for various $`N`$ and $`M`$ values, obtained by solving Eqs. (27-29) numerically. While in the Figures the initial hopping amplitudes have been generated completely randomly, similar results have been obtained when we used simple model estimates for the $`t_{mn}`$’s. Eq. (35) approximates very nicely the conductance above $`T^{}`$ even for rather small channel numbers. Below $`T^{}`$, however, the conductance starts to increase until it reaches its fixed point value $`g_{\mathrm{fp}}1`$. In the inset we show that the scale $`T^{}`$ decreases exponentially with the number of modes in agreement with Eq. (37).
## IV Discussion of the SET
To derive the scaling equations for the SET we repeated the derivation of the scaling equations of the two-level system with a generalized version of the interaction Hamiltonian (14). Here we only quote the results.
The scaling equations become:
$`{\displaystyle \frac{d𝐭^{(\mathrm{f})}}{dx}}`$ $`=`$ $`2(𝐭^{(\mathrm{f})}𝐪^{(\mathrm{f})}\stackrel{~}{𝐪}^{(\mathrm{f})}𝐭^{(\mathrm{f})})`$ (38)
$``$ $`2𝐭^{(\mathrm{f})}{\displaystyle \underset{f^{}}{}}[\mathrm{Tr}\{𝐭^{(\mathrm{f}^{})}𝐭_{}^{(\mathrm{f}^{})}{}_{}{}^{}\}+2\mathrm{T}\mathrm{r}\{𝐪^{(\mathrm{f}^{})}𝐪^{(\mathrm{f}^{})}\}`$ (39)
$`+`$ $`2\mathrm{T}\mathrm{r}\{\stackrel{~}{𝐪}^{(\mathrm{f}^{})}\stackrel{~}{𝐪}^{(\mathrm{f}^{})}\}],`$ (40)
$`{\displaystyle \frac{d\stackrel{~}{𝐪}^{(\mathrm{f})}}{dx}}`$ $`=`$ $`𝐭_{}^{(\mathrm{f})}{}_{}{}^{}𝐭^{(\mathrm{f})}4\stackrel{~}{𝐪}^{(\mathrm{f})}{\displaystyle \underset{f^{}}{}}\mathrm{Tr}\{𝐭^{(\mathrm{f}^{})}𝐭_{}^{(\mathrm{f}^{})}{}_{}{}^{}\},`$ (41)
$`{\displaystyle \frac{d𝐪^{(\mathrm{f})}}{dx}}`$ $`=`$ $`𝐭^{(\mathrm{f})}𝐭_{}^{(\mathrm{f})}{}_{}{}^{}4𝐪^{(\mathrm{f})}{\displaystyle \underset{f^{}}{}}\mathrm{Tr}\{𝐭^{(\mathrm{f}^{})}𝐭_{}^{(\mathrm{f}^{})}{}_{}{}^{}\}.`$ (42)
Similarly to the SEB, the ’incoherent’ scaling equations can be obtained by dropping all second order ’coherent’ terms in the equations above. In this way we obtain for the dimensionless conductances $`g^{(1)}`$ and $`g^{(2)}`$ of the two junctions and the effective field $`h`$ the following scaling equations:
$`{\displaystyle \frac{dg^{(1)}}{dx}}`$ $`=`$ $`4g^{(1)}(g^{(1)}+g^{(2)}),`$ (43)
$`{\displaystyle \frac{dg^{(2)}}{dx}}`$ $`=`$ $`4g^{(2)}(g^{(1)}+g^{(2)}),`$ (44)
$`{\displaystyle \frac{dh}{dx}}`$ $`=`$ $`4h(g^{(1)}+g^{(2)}).`$ (45)
These equations can be readily solved to give:
$`g^{(1)}(x)={\displaystyle \frac{g_0^{(1)}}{1+4(g_0^{(1)}+g_0^{(2)})x}},`$ (46)
$`g^{(2)}(x)={\displaystyle \frac{g_0^{(2)}}{1+4(g_0^{(1)}+g_0^{(2)})x}},`$ (47)
$`h(x)={\displaystyle \frac{h_0}{1+4(g_0^{(1)}+g_0^{(2)})x}}.`$ (48)
Thus in this approximation the only effect of the presence of several leads is to replace the dimensionless conductance in the denominator of Eqs. (35) and (36) by the parallel conductance of all junctions attached to the island.
In Fig. 4 we show the conductance calculated from the solution of the full scaling equations Eq. (40-42) for $`N=M=15`$ conduction modes. Initially, both conductances follow Eqs. (47). However, at a temperature $`T^{}10^3E_C`$ a single conduction mode prevails in one of the junctions and starts to induce the two-channel Kondo effect. The conductance of this junction approaches a universal value characteristic to the two-channel Kondo fixed point while the resistivity of the other junction is suppressed to zero below this temperature. Since the tunneling between the island and this latter lead is an irrelevant operator of dimension $`1/2`$ and therefore its conductance scales as $`t^2T`$, the total conductance of the device at the degeneracy point scales as $`GT`$, in agreement with Matveev’s result. Note that the conductance of the other junction is universal, i.e. independent of the number of modes in the junction.
## V Conclusions
In the present work we studied in detail the physics of the single electron box and the single electron transistor close to their degeneracy points using renormalization group methods. In particular, we investigated the effect of cross-mode scattering and showed that this is a relevant perturbation, and drives the system towards a stable two-channel Kondo fixed point, a prototype of non-Fermi liquid fixed points. At very low temperatures we recover Matveev’s mapping to the two-channel Kondo model: In this case at very low temperatures the system dynamically selects a single mode on the box and another one on one of the leads, and all the other modes become irrelevant. This fixed point has an $`SU(2)\times SU(2)\times U(1)`$ symmetry, where the first symmetry is generated dynamically and is connected to the structure of the effective tunneling Hamiltonian, while the second $`SU(2)`$ symmetry is associated to the real spin of the electrons and is responsible for the non-Fermi liquid behavior. The $`U(1)`$ symmetry is simply due to charge conservation.
Due to this two-channel Kondo fixed point the linear coefficient of the specific heat and the capacitance of the SET diverge logarithmically at the degeneracy point, $`c/TC\mathrm{ln}(T)`$, while the resistivity of the device diverges as $`1/T`$.
However, as our detailed analysis demonstrates, if the number of tunneling modes is larger than one, then a new small energy scale $`T^{}E_C\mathrm{exp}(CN)`$ enters into the problem, with $`N`$ the total number of tunneling modes. Above this scale coherent processes leading to the Kondo physics can be neglected, and the properties of the system are very well described solely in terms of the conductances of the various junctions. Neglecting the aforementioned coherent terms we were able to re-derive the equations of Ref. , and generalize them for the case of SET.
At this point we have to make an important remark. Our results for the two-channel Kondo fixed point rely heavily on the fact that the electrons at temperatures $`T`$ larger than the level spacing on it can hardly travel coherently from one junction to another. Neglecting this coherent process then one of the SET conductances scales to zero and so does the total conductance of the SET at the degeneracy point.
Apart from inelastic scattering, the main origin of this loss of coherence is in the random scattering from the wall of the island and the impurities on it. Any model assuming the existence of consecutive coherent tunneling events between different leads gives an essentially different, and very likely in most situations unphysical result at very low temperatures. Indeed, repeating our analysis for the model used in Ref. we find that at $`T=0`$ both tunnel junctions of the SET have a finite conductance at the degeneracy point, and thus the total conductance of the SET also remains finite even at zero temperature. This result is essentially different from the one obtained by Matveev or the calculations presented here, where coherent tunneling processes between different leads are excluded.
The difference between these two approximations becomes even more striking for the case of an $`M`$-fold degenerate multi-degeneracy points of a multi-dot system. In this case, neglecting the above-mentioned coherent tunneling between different leads, we find that the low-temperature physics is again described by the two-channel Kondo model. In the opposite case, however, the system would scale to another non-Fermi liquid fixed point, the $`SU(M)\times SU(2)`$ Coqblin-Schrieffer fixed point.
How and at what energy scale the crossover between these two behaviors happens, seems to be presently an open question, which can only be answered by somehow incorporating more details about the scattering on the impurities and the energy and spatial dependence of the tunneling matrix elements.
Because of the exponential dependence of $`T^{}`$, even a relatively small number of tunneling modes leads to an extremely small $`T^{}`$ and renders the non-Fermi liquid behavior in most cases inaccessible for the experimentalists. To observe the 2-channel Kondo behavior one can use semiconducting or metallic quantum boxes.
Semiconducting devices have the advantage that using suitable gate electrodes one can realize the idealistic case of a single mode contact, and therefore $`T^{}E_C`$ can be achieved. There is a serious difficulty, however, since one should keep $`E_C`$ large enough in order to have a measurable Kondo temperature while having a small level spacing on the island, the latter playing the role of an infrared cutoff for the Kondo physics. The former requirement immediately sets an upper limit on the size of the box $`d<e^2/E_C1000\AA `$, and therefore a lower limit on the mean level spacing $`\mathrm{\Delta }1/(m^{}d^2)0.11K`$, where we assumed a two dimensional electron gas with effective mass $`m^{}`$. This means that even if one manages to build a semiconducting device with a Kondo temperature in the measurable range, the level spacing will be too large to observe the two-channel Kondo behavior in detail. Indeed, even in the very recent experiments of Ref. the ratio $`E_C/\mathrm{\Delta }`$ was in the range $`100`$, and consequently only some fingerprints of the two-channel Kondo behavior could have been observed. For smaller islands $`E_C`$ would become larger, however, the ratio $`E_C/\mathrm{\Delta }`$ would get even worse. Therefore there seems to be no hope to investigate the two-channel Kondo behavior with semiconducting devices more in detail.
The other possibility is to prepare metallic boxes. Since these metallic boxes are three dimensional objects and $`m^{}m_e`$ (unlike semiconducting devices with $`m^{}m_e`$), metallic grains of the size of only $`d100\AA `$ may have quite large $`E_C100500K`$, and very small mean level spacing on the other hand. Indeed, using STM devices to tunnel into metallic drops it was possible to observe the Coulomb blockade even at room temperature.
The difficulty in this case is connected to the preparation of the junctions. As emphasized earlier, one needs practically single mode or at most few mode junctions in order to have $`T^{}`$ large enough, which requires atomic size contacts/junctions. To establish such a junction we propose the experimental setup illustrated in Fig. 5. In the suggested experiment a metallic electrode is covered by a thin insulating layer, and a metallic drop is deposited on the top of it. The atomic size contacts can be formed by plugging an STM needle into one of the drops and then gently pulling it out of it. An additional gate electrode should be built in the vicinity of the drop to control the charging of the island. The two-channel Kondo effect would then show up in the gate voltage dependence of the conductance through the island.
The biggest difficulty in the experiment above is to establish a stable contact, so that one has enough time to tune the drop to its degeneracy point and carry out the measurement. It would be much more advantageous to use nanotechnology to build atomic size contacts instead of using an STM, a possibility which may be not too far away any more.
The authors are grateful to G. Schön, and M. Devoret for valuable discussions. This research has been supported by the U.S - Hungarian Joint Fund No. 587, grant No. DE-FG03-97ER45640 of the U.S DOE Office of Science, Division of Materials Research, and Hungarian Grant Nr. OTKA T026327, OTKA F29236, and OTKA T029813.
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# Quantum Phase Transitions of a Square-Lattice Heisenberg Antiferromagnet with Two Kinds of Nearest-Neighbour Bonds: A High-Order Coupled-Cluster Treatment
## I Introduction
The subject of quantum spin-half antiferromagnetism in low-dimensional systems has attracted a great deal of interest in recent times in connection with the magnetic properties of the cuprate high-temperature superconductors. However, low-dimensional quantum spin systems are of interest in their own right as examples of strongly interacting quantum many-body systems. Although we know from the Mermin-Wagner theorem that thermal fluctuations are strong enough to destroy magnetic long-range order at any finite temperature, the role of quantum fluctuations is less understood. As a result of intensive work in the late eighties, it is now well-established that the ground-state of the Heisenberg antiferromagnet on the square lattice with nearest-neighbour interactions is long-range ordered (see for example the review in Ref. ). However, Anderson’s and Fazekas’ investigations of the triangular lattice led to a conjecture that quantum fluctuations plus frustration may be sufficient to destroy the Néel-like long-range order in two dimensions. Another specific area of recent research is the spin-half $`J_1`$$`J_2`$ antiferromagnet on the square-lattice where the frustrating diagonal $`J_2`$ bonds plus quantum fluctuations are able to realize a second-order transition from Néel ordering to a disordered quantum spin liquid (see for example Refs. and references therein). On the other hand, there are cases in which frustration causes a first-order transition in quantum spin systems in contrast to a second-order transition in the corresponding classical model (see for example Refs. ).
Besides frustration, there is another mechanism to realize the “melting” of Néel ordering in the ground states of unfrustrated Heisenberg antiferromagnets, namely the formation of local singlet pairs of two coupled spins. This mechanism may be relevant for the quantum disordered state in bilayer systems as well as in CaV<sub>4</sub>O<sub>9</sub> (see for example Refs. and references therein). The formation of local singlets is connected with a gap in the excitation spectrum. By contrast, the opening of a gap in the excitation spectrum of frustrated systems seems to be less clear and might be dependent on details of the exchange interactions.
In the present paper, we study a model which contains both mechanisms, frustration and singlet formation, in different parameter regions. We mainly use in this article the coupled cluster method, which has become widely recognized as one of the most powerful and most universal techniques in quantum many-body theory. In recent years there has been increasing success in applying the CCM to quantum spin systems, especially with the advent of high-order approximations which utilize computer algebra. Subsequently, high-order CCM approximations have been applied to the $`XXZ`$ model, the anisotropic $`XY`$ model and the $`J_1`$$`J_2`$ model. In addition to the CCM results we also present variational, spin-wave theory (SWT) and exact diagonalization (ED) results for the sake of comparison.
## II The model
We consider a spin-half Heisenberg model on a square lattice with nearest neighbor bonds $`J`$ and $`J^{}`$ in a regular zigzag pattern as shown in Fig. 1. The Hamiltonian is given by
$`H`$ $`=`$ $`J{\displaystyle \underset{<ij>_1}{}}\text{s}_i\text{s}_j+J^{}{\displaystyle \underset{<ij>_2}{}}\text{s}_i\text{s}_j`$ (1)
$`\stackrel{J=1}{=}`$ $`{\displaystyle \underset{iA}{}}{\displaystyle \underset{p}{}}(1+\delta _{p,p_J^{}}(J^{}1))\text{s}_i\text{s}_{i+p}.`$ (2)
The sums over $`<ij>_1`$, and $`<ij>_2`$ represent sums over the nearest-neighbour bonds shown in Fig. 1 by dashed and solid lines, respectively. Throughout the paper we fix the $`J`$ bond to be antiferromagnetic ($`J>0`$) and henceforth scale it to the value $`J=1`$, and consider $`J^{}`$ as the free parameter of the model. We also split the square lattice into the equivalent $`A`$ and $`B`$ sublattices shown in Fig. 1. In Eq. (1) the sum over $`i`$ runs over the sites of the sublattice $`A`$, with vectors $`p=\{(0,\pm 1),(\pm 1,0)\}`$ connecting nearest neighbours. In particular, $`p_J^{}=(1,0)`$ pertains to the coupling with $`J^{}`$ bonds.
Each square lattice plaquette consists of three $`J=1`$ bonds and one $`J^{}`$ bond. In the case of ferromagnetic $`J^{}`$ bonds (i.e., $`J^{}<0`$), the plaquettes are frustrated. Conversely, for antiferromagnetic $`J^{}`$ bonds (i.e., $`J^{}>0`$) there is no frustration in the system, although the difference of the coupling strengths $`J`$ and $`J^{}`$ leads to quantum competition. This model has been treated previously using perturbation theory, renormalized spin wave theory (RSWT) and exact diagonalization (ED). It allows us to study the influence of local singlet formation ($`J^{}>1`$) and frustration ($`J^{}<0`$) on the stability of the Néel order within a single model.
Ferromagnetic bonds in an antiferromagnetic matrix have been discussed in recent times in connection with the proposal by Aharony and coworkers to model localized oxygen holes in the Cu-O-planes by local ferromagnetic bonds between the copper spins. It was argued that random ferromagnetic bonds may influence the antiferromagnetic order drastically and may support the realization of a quantum spin-liquid state.
On the other hand, the case of antiferromagnetic $`J^{}`$ bonds with $`J^{}>1`$ resembles the situation in bilayer systems and in the depleted square-lattice antiferromagnet CaV<sub>4</sub>O<sub>9</sub>, in which the competition between two different antiferromagnetic bonds leads to a phase transition from antiferromagnetic long-range order to quantum disorder with a finite gap. It is seen in this article that the transition point obtained for the model of Eq. (1) is quite close to that obtained for the bilayer model.
There are some special cases of the model Hamiltonian of Eq. (1):
(i) $`J^{}=1`$: square-lattice antiferromagnet, for which the ground state is long-range ordered;
(ii) $`J^{}=0`$: honeycomb-lattice antiferromagnet, for which the ground state is long-range ordered;
(iii) $`J^{}=\mathrm{}`$: spin-$`1`$ triangular lattice, for which the ground state is long-range ordered; and
(iv) $`J^{}=+\mathrm{}`$: valence-bond solid, for which the ground state is a rotationally invariant quantum dimer state with an excitation gap.
Classical ground state. For $`J^{}>1/3`$ the Néel state is the classical ground state of the Hamiltonian of Eq. (1). At $`J_c^{}=1/3`$ there is classically a second-order phase transition to a ground state of helical nature (see Fig. 1), with a characteristic pitch angle $`\mathrm{\Phi }=\pm |\mathrm{\Phi }_{\mathrm{cl}}|`$ given by
$$|\mathrm{\Phi }_{\mathrm{cl}}|=\{\begin{array}{cc}0\hfill & J^{}>\frac{1}{3}\hfill \\ \mathrm{arccos}\left(\frac{1}{2}\sqrt{1\frac{1}{J^{}}}\right)\hfill & J^{}\frac{1}{3}\hfill \end{array}$$
(3)
where the different signs correspond to the two chiralities of this helical state. Note that for $`\mathrm{\Phi }=0`$ this is just the Néel state. More generally, the pitch angle varies with $`J^{}`$ from $`|\mathrm{\Phi }_{\mathrm{cl}}|=0`$ for $`J^{}>1/3`$ to $`|\mathrm{\Phi }_{\mathrm{cl}}|=\pi /3`$ for $`J^{}=\mathrm{}`$. Note that $`|\mathrm{\Phi }_{\mathrm{cl}}|=\pi /3`$ (realized at $`J^{}=\mathrm{}`$) corresponds to the ground state of the spin-1 triangular lattice. We describe the directions of the spins $`\text{s}_A`$ and $`\text{s}_B`$, belonging to the $`A`$ and $`B`$ sublattices respectively, for the classical helical state with a characteristic angle $`\mathrm{\Phi }`$ as follows (and see Fig. (1)),
$`\text{s}_A(𝐑)`$ $`=`$ $`\widehat{𝐮}\mathrm{cos}𝐐𝐑+\widehat{𝐯}\mathrm{sin}𝐐𝐑,`$ (4)
$`\text{s}_B(𝐑+\widehat{x})`$ $`=`$ $`\widehat{𝐮}\mathrm{cos}(𝐐𝐑+\pi +3\mathrm{\Phi })+\widehat{𝐯}\mathrm{sin}(𝐐𝐑+\pi +3\mathrm{\Phi }),`$ (5)
where $`\widehat{𝐮}`$ and $`\widehat{𝐯}`$ are perpendicular unit vectors in the spin space, $`𝐑`$ runs over the sites of the sublattice $`A`$, and we have $`𝐐=(2\mathrm{\Phi },0)`$ for the pitch vector $`𝐐`$. We note that this general helical state does not have a periodicity in the $`x`$-direction because $`\mathrm{\Phi }`$ is in general not of the form $`m\pi /n`$ with $`m`$ and $`n`$ integral. We also note that we have only three different angles between nearest-neighbour spins, namely $`\pm (\pi \mathrm{\Phi })`$ for the $`J=1`$ couplings and $`\pi 3\mathrm{\Phi }`$ for the coupling with $`J^{}`$.
The maximum frustration is in the region around $`J^{}1`$. Bearing in mind the situation for the $`J_1`$$`J_2`$ model, one might expect that for the extreme quantum case (spin-half) quantum fluctuations might be able to open the window to a spin-liquid phase for a finite range of parameters around this region of maximum frustration. On the other hand, for strong antiferromagnetic bonds ($`J^{}1`$) there is, of course, no indication in the classical model for the breakdown of the Néel order.
Simple variational ansatz for the quantum ground-state. In the quantum case, the region of strong antiferromagnetic $`J^{}`$ bonds ($`J^{}1`$) is characterized by a tendency to singlet pairing of the two spins corresponding to a $`J^{}`$ bond. Using a high-order series expansion the Néel order was found to be stable up to a critical value $`J_s^{}2.56`$.
A comparable value can be obtained using a simple variational wave function similar to that used for bilayer systems, namely
$$|\mathrm{\Psi }_{\mathrm{var}}=\underset{iA}{}\frac{1}{\sqrt{1+t^2}}[|_i_{i+\widehat{x}}t|_i_{i+\widehat{x}}],$$
(6)
where the lattice sites $`i`$ and $`i+\widehat{x}`$ correspond to a $`J^{}`$ bond, and where the product in Eq. (6) is thus effectively taken over the $`J^{}`$ bonds of the lattice of Eq. (1). The trial function depends on the variational parameter $`t`$ and interpolates between a valence-bond state realized for $`t=1`$ and the Néel state for $`t=0`$. For $`t=1`$, the singlet pairing is complete and $`|\mathrm{\Psi }_{\mathrm{var}}`$ represents an eigenstate of the model of Eq. (1) in the limit $`J^{}\mathrm{}`$ (dimer state). By minimizing $`\mathrm{\Psi }_{\mathrm{var}}|H|\mathrm{\Psi }_{\mathrm{var}}`$ with respect to the variational parameter $`t`$ we get an upper bound for the ground-state energy per spin of the model of Eq. (1),
$$E_{\mathrm{var}}/N=\{\begin{array}{cc}(J^2+3J^{}+9)/24\hfill & J^{}3\hfill \\ 3J^{}/8\hfill & J^{}>3\hfill \end{array}.$$
(7)
The relevant order parameter describing the Néel order is
$$M_{\mathrm{var}}=\mathrm{\Psi }_{\mathrm{var}}|s_i^z|\mathrm{\Psi }_{\mathrm{var}}=\{\begin{array}{cc}1/2\sqrt{1J^2/9}\hfill & J^{}3\hfill \\ 0\hfill & J^{}>3\hfill \end{array},$$
(8)
showing a breakdown of the Néel order at a critical value $`J_s^{}=3`$.
## III Coupled cluster calculations
### A The ground-state formalism
The starting point for any CCM calculation (see overwiev in Ref. ) is the choice of a normalized model or reference state $`|\mathrm{\Phi }`$, together with a set of mutually commuting multispin creation operators $`C_I^+`$ which are defined over a complete set of many-body configurations $`I`$. The operators $`C_I`$ are the multispin destruction operators and are defined to be the Hermitian adjoints of the $`C_I^+`$. We choose $`\{|\mathrm{\Phi };C_I^+\}`$ in such a way that we have $`\mathrm{\Phi }|C_I^+=0=C_I|\mathrm{\Phi }`$, $`I0`$, where, by definiton, $`C_0^+=𝟙`$, the identity operator.
For spin systems, an appropriate choice for the CCM model state $`|\mathrm{\Phi }`$ is often a classical spin state, in which the most general situation is that each spin can point in an arbitrary direction. For the case of the Hamiltonian of Eq. (1), we choose the helical state illustrated in Fig. 1 to be our model state. Although the classical ground state of Eq. (1) is precisely of this form, we do not choose the classical result for the pitch angle $`\mathrm{\Phi }`$ but we consider it rather as a free parameter in the CCM calculation.
In order to perform a CCM calculation, we would like to treat each site equivalently and we do this by performing a rotation of the local spin axes at each site about the $`y`$-axis such that all spins in the model state align in the same direction, say down (along the negative $`z`$-axis). After this transformation we have
$$|\mathrm{\Phi }=|\mathrm{}\mathrm{};C_I^+=s_r^+,s_r^+s_r^{}^+,s_r^+s_r^{}^+s_{r^{\prime \prime }}^+,\mathrm{},$$
(9)
(where the indices $`r,r^{},r^{\prime \prime },\mathrm{}`$ denote any lattice site) respectively, for the model state and the multispin creation operators, which now consist of spin-raising operators only.
In order to make the spin $`\text{s}_i`$ point down let us suppose we need to perform such a rotation of the spin axes by an angle $`\delta _i`$. This is equivalent to the transformation
$`s_i^x`$ $``$ $`\mathrm{cos}\delta _is_i^x+\mathrm{sin}\delta _is_i^z`$ (10)
$`s_i^y`$ $``$ $`s_i^y`$ (11)
$`s_i^z`$ $``$ $`\mathrm{sin}\delta _is_i^x+\mathrm{cos}\delta _is_i^z.`$ (12)
A similar rotation about the $`y`$-axis by an angle $`\delta _j`$ is performed for the spin $`\text{s}_j`$. Thus we get for the transformation of the scalar product of the two spins, $`\text{s}_i\text{s}_j(\text{s}_i\text{s}_j)_\phi `$, where
$`(\text{s}_i\text{s}_j)_\phi `$ $``$ $`\mathrm{sin}\phi [s_i^xs_j^zs_i^zs_j^x]+\mathrm{cos}\phi [s_i^xs_j^x+s_i^zs_j^z]+s_i^ys_j^y`$ (13)
$`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{sin}\phi [s_i^+s_j^zs_i^zs_j^++s_i^{}s_j^zs_i^zs_j^{}]+\mathrm{cos}\phi s_i^zs_j^z`$ (16)
$`+{\displaystyle \frac{1}{4}}(\mathrm{cos}\phi +1)[s_i^+s_j^{}+s_i^{}s_j^+]`$
$`+{\displaystyle \frac{1}{4}}(\mathrm{cos}\phi 1)[s_i^+s_j^++s_i^{}s_j^{}].`$
The angle $`\phi \delta _j\delta _i`$ is the angle between the two spins, and $`s^\pm s^x\pm \mathrm{i}s^y`$ are the spin-raising and spin-lowering operators. Note that this product of two spins after the rotation depends not only on the angle between them, but also on the sign of this angle. In case of the Néel model state ($`\mathrm{\Phi }=0`$), the angle between any neighbouring spins is $`\pi `$, and hence Eq. (13) becomes $`\text{s}_i\text{s}_j\frac{1}{2}[s_i^+s_j^++s_i^{}s_j^{}]s_i^zs_j^z`$.
Using the helical state of Eq. (4) with the characteristic angle $`\mathrm{\Phi }`$, the Hamiltonian of Eq. (1) is now rewritten in the local coordinates as,
$$H=\underset{iA}{}\underset{p}{}(1+\delta _{p,p_J^{}}(J^{}1))(\text{s}_i\text{s}_{i+p})_{\phi _p},$$
(17)
where the angles between neighbouring spins are $`\phi _{\pm \widehat{y}}=\pi +\mathrm{\Phi }`$, $`\phi _{\widehat{x}}=\pi \mathrm{\Phi }`$ and $`\phi _{\widehat{x}}=\pi +3\mathrm{\Phi }`$. While the general helical state (see Fig. 1) does not have translational symmetry in the $`x`$-direction, the transformed Hamiltonian of Eq. (17) does have this symmetry since it depends only on the angles between neighbouring spins.
Having defined an appropriate model state $`|\mathrm{\Phi }`$ with creation operators $`C_I^+`$, the CCM parameterizations of the ket and bra ground states are given by
$$|\mathrm{\Psi }=e^S|\mathrm{\Phi },S=\underset{I0}{}𝒮_IC_I^+,$$
(18)
$$\stackrel{~}{\mathrm{\Psi }}|=\mathrm{\Phi }|\stackrel{~}{S}e^S,\stackrel{~}{S}=1+\underset{I0}{}\stackrel{~}{𝒮}_IC_I.$$
(19)
The correlation operator $`S`$ is expressed in terms of the creation operators $`C_I^+`$ and the ket-state correlation coefficients $`𝒮_I`$. We can now write the ground-state energy as,
$$E=\mathrm{\Phi }|e^SHe^S|\mathrm{\Phi }.$$
(20)
To describe the magnetic order of the system, we use a simple order parameter which is expressed in terms of the local, rotated spin axes, and which is given by
$$M\stackrel{~}{\mathrm{\Psi }}|s_i^z|\mathrm{\Psi },$$
(21)
such that the order parameter represents the on-site magnetization. Note that $`M`$ is the usual sublattice magnetization for the case of the Néel state as the CCM model state.
To find the ket-state and bra-state correlation coefficients we have to require that the expectation value $`\overline{H}=\stackrel{~}{\mathrm{\Psi }}|H|\mathrm{\Psi }`$ is a minimum with respect to the bra-state and ket-state correlation coefficients. This formalism is exact if we include all possible multispin configurations in the correlation operators $`S`$ and $`\stackrel{~}{S}`$, which is usually impossible in a practical situation. We use the LSUB$`n`$ approximation scheme to truncate the expansion of $`S`$ and $`\stackrel{~}{S}`$ in the Eqs. (18) and (19).
Using the lattice symmetries, we have now to find all different possible configurations with respect to the point and space group symmetries of both the lattice and Hamiltonian with up to $`n`$ spins spanning a range of no more than $`n`$ adjacent lattice sites (LSUB$`n`$ approximation) and these are referred to as the fundamental configurations.
The Hamiltonian of Eq. (1) has four lattice point-group symmetries namely two rotational operations ($`0^{}`$, $`180^{}`$) and two reflections (along the $`x`$\- and $`y`$-axes), defined by:
$`xx,yy;`$ $`x(x+1),yy,`$ (22)
$`xx,yy;`$ $`x(x+1),yy.`$ (23)
The rotation of $`180^{}`$ and the reflection along the $`y`$-axis are connected by a shift of $`\widehat{x}=(1,0)`$. The translational operator $`T`$ is defined by
$$T=(n+m)\widehat{x}+(mn)\widehat{y},n\text{}m\text{ integral},$$
(24)
such that translational symmetry is preserved.
The Néel model state also contains these symmetries, and so for this model state we can directly apply all these symmetries in finding the fundamental configurations. On the other hand the general helical model state ($`\mathrm{\Phi }0`$) has only two of the above four lattice point-group symmetries, namely $`xx,yy`$ and $`xx,yy`$, and so this reduced symmetry yields a larger number of fundamental configurations.
In the case of the Néel model state ($`\mathrm{\Phi }=0`$), the number of fundamental configurations can further be reduced by explicitly conserving the total uniform magnetization $`s_T^z_ks_k^z`$ (the sum on $`k`$ runs over all lattice sites) because the ground state is known to lie in the $`s_T^z=0`$ subspace. This means that we exclude configurations with an odd number of spins, and therefore we do not use LSUB3, LSUB5, etc. approximations. The helical state is not an eigenstate of $`s_T^z`$ and we cannot apply this property when using the helical model state. The fundamental configurations can now be calculated computationally, and the resulting numbers of LSUB$`n`$ configurations for $`n8`$ are given in Table I.
The Ket-state and bra-state equations are calculated computationally. For the Néel model state, we are able to carry out the CCM up to the LSUB8 level (where we need to solve 4986 coupled equations), whereas for the helical state we could do this only up to the LSUB6 level (where we need to solve 1638 coupled equations).
### B The excited state formalism
We use the excited-state formalism of Emrich to approximate the excited-state wavefunctions. We apply an excitation operator $`X^e`$ linearly to the ket state wavefunction (18), such that
$$|\mathrm{\Psi }_e=X^ee^S|\mathrm{\Phi };X^e=\underset{I0}{}𝒳_I^eC_I^+.$$
(25)
Using the Schrödinger equation, $`H|\mathrm{\Psi }_e=E_e|\mathrm{\Psi }_e`$, we find that
$$ϵ_eX^e|\mathrm{\Phi }=e^S[H,X^e]_{}e^S|\mathrm{\Phi },$$
(26)
where $`ϵ_e`$ ($`E_eE`$) is the difference between the excited-state energy ($`E_e`$) and the ground-state energy ($`E`$). Applying $`\mathrm{\Phi }|C_I`$ to Eq. (26) we find that,
$$ϵ_e𝒳_I^e=\mathrm{\Phi }|C_Ie^S[H,X^e]_{}e^S|\mathrm{\Phi },$$
(27)
which is an eigenvalue equation with eigenvalues $`ϵ_e`$ and corresponding eigenvectors $`𝒳_I^e`$.
As for the ground state, we must use an approximation scheme for $`X^e`$ in Eq. (25). Although it is not necessary to use the same approximation for the excited state as for the ground state, we in fact do so to keep the CCM calculations as systematic and self-consistent as possible. We define the fundamental configurations for LSUB$`n`$ (for the Néel state) as previously, but we now restrict the choice of configurations to contain only those which produce a change of $`s_T^z`$ of $`\pm 1`$ with respect to the model state. Since we are only interested in the lowest-lying excitations, the restriction to these single-magnon spin-wave-like excitations is the correct choice. The number of fundamental excited-state configurations for LSUB$`n`$ is given in Table I.
To calculate the terms of the right hand side of Eq. (27) we use the same computational algorithm as for the calculation of the ground-state ket equation. The terms contain the ground ket-state correlation coefficients $`𝒮_I`$, so once these coefficients have been determined the eigenvalue equation (27) can be solved (numerically). We furthermore choose the lowest energy eigenvalue of Eq. (27) in order to calculate the excitation energy gap, $`\mathrm{\Delta }`$. We note that the eigenvalues of Eq. (27) are not guaranteed to be real, since as a generalized eigenvalue equation it is not symmetric. However, over the entire regime of interest, the values of $`\mathrm{\Delta }`$ so obtained are found to be real. We have performed these calculations for the excited state up to the LSUB6 level of approximation.
### C Extrapolation of the CCM-LSUB$`n`$ results
Although no scaling theory for results of LSUB$`n`$ approximations has yet been proven, there are empirical indications of scaling laws for the energy, the magnetization, and the excited-state energy gap for various spin models.
These scaling laws can be justified by the observations that they fit the results well (i.e., with low mean-square deviation), and that the extrapolated results are in good agreement with results of other methods (e.g., Green function Monte Carlo or series expansion for the 2D $`XXZ`$ model) or with exact results (e.g., 1D $`XY`$ model). In accordance with those previous results we use the following scaling laws: for the ground-state energy,
$$E=a_0+a_1(1/n^2)+a_2(1/n^2)^2;$$
(28)
for the ground-state magnetization,
$$M=b_0+b_1(1/n)+b_2(1/n)^2;$$
(29)
and for the gap of the lowest-lying excitations,
$$\mathrm{\Delta }=c_0+c_1(1/n)+c_2(1/n)^2;$$
(30)
where $`n`$ is the LSUB$`n`$ approximation level.
### D Choice of the CCM model state
As stated previously, we use the helical state of Eq. (4) with the characteristic angle $`\mathrm{\Phi }`$, illustrated in Fig. 1, as the model state for the CCM. We must therefore make a selection of an appropriate value for $`\mathrm{\Phi }`$. A possible choice would be the classical ground state of the Hamiltonian of (1) (i.e., $`\mathrm{\Phi }=\mathrm{\Phi }_{\mathrm{cl}}`$ as given by Eq. (3)).
Another possibility is to perform a CCM-LSUB$`n`$ approximation calculation and then to minimize the corresponding LSUB$`n`$ approximation to the energy with respect to $`\mathrm{\Phi }`$,
$$E_{\mathrm{LSUB}n}(\mathrm{\Phi })\mathrm{min}\mathrm{\Phi }=\mathrm{\Phi }_{\mathrm{LSUB}n}.$$
(31)
The results for $`\mathrm{\Phi }_{\mathrm{LSUB}n}`$ will be given later (Fig. 7). However, we note now that although the CCM does not yield a strict upper bound for the ground-state energy, using $`\mathrm{\Phi }=\mathrm{\Phi }_{\mathrm{LSUB}n}`$ (i.e., using the CCM with a variational parameter) has been found to be a reasonable assumption. There are several additional arguments to suggest that $`\mathrm{\Phi }=\mathrm{\Phi }_{\mathrm{LSUB}n}`$ is indeed a better choice than $`\mathrm{\Phi }=\mathrm{\Phi }_{\mathrm{cl}}`$, as indicated below.
In the first place we note that we cannot find solutions for the LSUB6 equations using $`\mathrm{\Phi }=\mathrm{\Phi }_{\mathrm{cl}}(J^{})`$ in the region $`0.7J^{}0.47`$ insofar as the Newton method used to solve these equations does not converge in that region. This is a clear indication that this model state is not a good one. By contrast, such behaviour is not found for $`\mathrm{\Phi }=\mathrm{\Phi }_{\mathrm{LSUB}n}`$.
Secondly, it is generally known that quantum fluctuations tend to prefer collinear order (e.g., Néel order). We will indeed find (and see Fig. 7 below) that the Néel ordering ($`\mathrm{\Phi }=0`$) seems to survive for some $`J^{}<1/3`$, in which region it has already broken down in the classical case. This is also in agreement with results of exact diagonalizations for our model (and see Fig. 5 below).
Thirdly we find better agreement of the CCM results for the energy compared to exact diagonalization results by using the helical state as the model state with the value $`\mathrm{\Phi }=\mathrm{\Phi }_{\mathrm{LSUB}n}`$ rather than with the classical value $`\mathrm{\Phi }=\mathrm{\Phi }_{\mathrm{cl}}`$. We find that CCM results for the ground-state energy usually agree well with the corresponding ED results (and with results of other methods), provided that a good CCM model state is chosen.
We therefore use the helical state with $`\mathrm{\Phi }=\mathrm{\Phi }_{\mathrm{LSUB}n}`$ as the CCM model state throughout this paper. Note that for $`J^{}1/3`$ this model state is identical to the classical ground state of Eq. (1) but that for $`J^{}<1/3`$ it is not.
## IV Results
Using the CCM scheme described above, we calculate the approximate ground state and the low-lying excitations of the Hamiltonian of Eq. (1). For comparison we also exactly diagonalize finite sized lattices of square shape. We use periodic boundary conditions with $`N=16,18,20,26`$ and $`32`$ spins, and we extrapolate to the infinite system using standard finite-size scaling laws. We present results for the ground-state energy, the order parameter and the excitation gap. We examine the formation of local singlets (for $`J^{}>1`$), the effects of frustration (for $`J^{}<0`$), and the special case of the honeycomb lattice ($`J^{}=0`$).
$`J^{}>1`$: Formation of local singlets. Using the CCM we obtain clear indications of a second-order phase transition to a disordered dimer-like phase at a certain critical value of $`J^{}`$, namely $`J_s^{}`$. For $`J^{}>J_s^{}`$, the Néel-like long range order melts, (i.e., the sublattice magnetization, $`M`$, given by Eq. (21) becomes zero). Our estimate for $`J_s^{}`$ using the four extrapolated LSUB$`n`$ results for $`M`$ with $`n=2,4,6,8`$ (see Fig. 2) is $`J_s^{}3.41`$. However, using only the three CCM LSUB$`n`$ approximations with $`n=4,6,8`$ for the extrapolation, we obtain a value $`J_s^{}3.16`$, which indicates that the true value could be even somewhat smaller. This is in agreement with our corresponding result using exact diagonalizations of small systems. By using the extrapolation scheme of Ref. , we find a critical value $`J_s^{}2.45`$ for the magnetization. Note however that better accuracy requires larger systems because of the exact diagonalization (ED) extrapolation ansatz for $`M`$ (i.e., $`M=M_{\mathrm{}}+\mathrm{const}\times N^{1/2}`$). Therefore, we cannot consider the ED results for the magnetization (and see Fig. 2) as quantitatively correct. Our two results for the critical $`J_s^{}`$ also agree with the estimate $`J_s^{}2.56`$ from series expansion, and even the result $`J_s^{}=3`$ from the simple variational ansatz of Eq. (6) agrees surprisingly well with these values. By contrast, the second-order renormalized spin wave theory (RSWT) gives the larger result $`J_s^{}5.0`$, indicating that the standard spin-wave approach is insufficient to describe this type of transition.
Another indication of a dimerized phase is the appearance of a gap $`\mathrm{\Delta }`$ between the ground state and the lowest-lying excited state. We clearly expect a spectrum with gapless Goldstone modes if the ground state is Néel long-range ordered, whereas for a disordered singlet ground-state the formation of triplet excitations may cost a finite amount of energy. This behaviour is reflected by our results using both CCM and exact diagonalization (see Fig. 3), which agree well with each other. For $`J^{}J_s^{}`$, there is a gap proportional to $`J^{}`$, corresponding to the dimer-like nature of the ground state. The gap obviously opens in the range $`2.5J_s^{}3.0`$ in both the ED and CCM calculations. This is in good agreement with the corresponding estimates for the critical point using the order parameter. Note, that the standard linear spin wave theory fails in calculating the gap (i.e., gives gapless modes for all values of $`J^{}`$ with $`J^{}>0`$).
By comparing the results for the ground-state energy, we find excellent agreement between the CCM results and the results from exact diagonalization (see Fig. 4) for $`J^{}>0`$. By contrast, spin-wave theory (SWT) calculations show a significant deviation from these results for larger $`J^{}1`$. These spin-wave results are obviously poor since the simple upper bound for the energy given by Eq. (7) (e.g., $`E_0=1.5`$ for $`J^{}=4`$) is smaller than the corresponding SWT results (e.g., $`E_0=1.42`$ from 2nd-order RSWT). By contrast, CCM and ED results (both are about $`E_0=1.54`$ for $`J^{}=4`$) are slightly smaller than the variational result. While both CCM and SWT calculations have the Néel state as starting point, we find the CCM is much better able than SWT to describe the transition to the rotationally invariant disordered state and to the completely dimerized state (represented by the variational function of Eq. (6) with $`t=1`$). Note that even the simplest CCM approximation (LSUB2) gives the correct asymptotic result for the energy (i.e., Eq. (7)) for very large values of $`J^{}`$, whereas SWT does not.
For the case of the pure square-lattice Heisenberg antiferromagnet (i.e., $`J^{}=1`$), we reproduce the CCM results of Refs. , which have already been demonstrated to agree well with those from other methods.
$`J^{}=0`$: honeycomb lattice. For the special case of $`J^{}=0`$ (which is equivalent to the honeycomb lattice), we find that the CCM and the ED results are in good agreement (see Table II). However, the magnetization $`M`$ for ED is found to be smaller then the CCM result. We note, however, that the CCM result for $`M`$ at this point agrees with the result of high-order SWT ($`M=0.28`$) as well as with the result of series expansion ($`M=0.26`$), although it does not agree so well with the result of Monte Carlo calculations ($`M=0.22`$). We note too that our CCM results here agree perfectly with previous lower-order CCM calculations.
$`J^{}<0`$: Frustration. For $`J^{}2`$, we find that the extrapolated ED results for the energy lie appreciably above the CCM (and SWT) results (see Fig. 4). This is because the energies for the small lattices considered do not fit well to the finite-size scaling law ($`E/N=E_{\mathrm{}}/N+\mathrm{const}\times N^{3/2}`$) in this region. The finite-size effects for systems with an incommensurate helical structure are found to be larger than for systems with, for example, Néel order or with dimerized spin pairs. However, we find that our best ED result (with 32 spins) shows only very small deviations from the CCM result, even in the frustrated region.
While classically we have a second-order phase transition (from Néel order to helical order) at $`J_c^{}=1/3`$, using the CCM we find indications for a shift of this critical point to a value $`J_c^{}1.35`$ (see Fig. 7) in the quantum case. The ED data of the structure factors (see Fig. 5) also show a shift of the transition to stronger ferromagnetic $`J^{}`$ bonds. Both correspond to the general picture that quantum fluctuations prefer a collinear ordering (such as Néel order). Hence this ordered state can survive for the quantum case into a region where classically it is already unstable. The Néel model state ($`\mathrm{\Phi }=0`$) gives the minimum ground-state energy for all values $`J^{}>J_c^{}`$, where $`J_c^{}`$ is also dependent on the level of LSUB$`n`$ approximation level. For $`J^{}<J_c^{}`$ another minimum in the energy for $`\mathrm{\Phi }0`$ is found to lie lower than the minimum at $`\mathrm{\Phi }=0`$ (see Fig. 6). The state for $`\mathrm{\Phi }0`$ is believed to be a quantum analogue of the classical ground state in two dimensions. Furthermore, the crossover from one minimum solution to the other is not smooth but is abrupt at this point (see Fig. 6 and Fig. 7). This behaviour is assumed here to be an indication of a phase transition. Furthermore, it might also indicate that this is a first-order phase transition and, consequently, that due to quantum fluctuations the nature of this phase transition is changed from the classical second-order type to a first-order type.
The behaviour of the order parameter (see Fig. 2) in the region around $`J_c^{}1.35`$ where we expect the above-mentioned quantum phase transition is also quite marked. We cannot extrapolate the LSUB$`n`$ results directly, because the phase transition points shift with the order of the LSUB$`n`$ approximations (Fig. 7). We thus find a large statistical deviation of the extrapolated results in the region $`1.4J^{}1.0`$. Hence, we use the minima for $`M`$ in that region to extrapolate an estimation of the order parameter. We find that minimum ($`M0.05`$) to be at a value of $`J^{}1.2`$. The extrapolated ED results do not agree very well with this result, since these give $`M`$ to be zero at $`J^{}0.8`$. However, these results are also very poor in that region because of the strong influence of the boundary conditions and large statistical errors. As a result of these difficulties we are not able to decide whether or not quantum fluctuations and frustration are able to form a disordered quantum spin liquid phase (i.e., with $`M=0`$) between the Néel state and the helical state for some finite frustrating $`J^{}<0`$ regime. However, the CCM results suggest that there is either no quantum spin liquid phase or that, if it does exist, it does so only in a very small region.
## V Summary
Using the CCM we have studied the influence of quantum spin fluctuations on both the ground-state phase diagram and the excited states of a spin-half square-lattice Heisenberg antiferromagnet with two kinds of nearest-neighbour exchange bonds. The phase diagram is found to contain a quantum helical phase, a Néel ordered phase, and a finite-gap quantum disordered phase. While we have clearly a second-order transition from the Néel phase to the finite-gap quantum disordered phase, we also found indications of a quantum-induced first-order transition from the Néel phase to the helical phase, for which classically we have a second-order transition. While our CCM results were in general in good agreement with the ED data, we found the CCM particularly good at describing the dimerized phase. By contrast, spin wave theory fails in that region due to enhanced longitudinal spin fluctuations. Accurate high-order CCM results for the antiferromagnet on the honeycomb lattice were also presented.
###### Acknowledgements.
We would like to thank N.B. Ivanov for his stimulating discussions. This work has also been supported in part by the Deutsche Forschungsgemeinschaft (GRK 14, Graduiertenkolleg on “Classification of Phase Transitions in Crystalline Materials”), and by a research grant (GR/M45429) from the Engineering and Physical Sciences Research Council (EPSRC) of Great Britain.
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# Mu and Tau Neutrino Thermalization and Production in Supernovae: Processes and Timescales
## I Introduction
The cores of protoneutron stars and core-collapse supernovae are characterized by mass densities of order $`10^{10}10^{14}`$ g cm<sup>-3</sup> and temperatures that range from $`1`$ to $`50`$ MeV. The matter is composed predominantly of nucleons, electrons, positrons, and neutrinos of all species. For $`\nu _\mu `$ and $`\nu _\tau `$ types (collectively ‘$`\nu _\mu `$s’), which carry away 50$``$60% of the $`23\times 10^{53}`$ ergs liberated during collapse and explosion, the prevailing opacity and production processes are $`\nu _\mu `$-electron scattering, $`\nu _\mu `$-nucleon scattering, electron-positron annihilation ($`e^+e^{}\nu _\mu \overline{\nu }_\mu `$), and nucleon-nucleon bremsstrahlung. While all of these processes contribute for the electron types ($`\nu _e`$s and $`\overline{\nu }_e`$s), the charged-current absorption processes $`\nu _enpe^{}`$ and $`\overline{\nu }_epne^+`$ dominate their opacity so completely that in this paper we address only $`\nu _\mu `$ production and thermalization.
Supernova theorists had long held that $`\nu _\mu `$-nucleon scattering was unimportant as a mechanism for neutrino equilibration. While this process was included as a source of opacity , it served only to redistribute the neutrinos in space, not in energy. In contrast, $`\nu _\mu `$-electron scattering was thought to dominate $`\nu _\mu `$ neutrino thermalization. In addition, the only $`\nu _\mu \overline{\nu }_\mu `$ pair production mechanisms employed in full supernova calculations were $`e^+e^{}\nu _\mu \overline{\nu }_\mu `$ and plasmon decay ($`\gamma _{pl}\nu _\mu \overline{\nu }_\mu `$) ; nucleon-nucleon bremsstrahlung was neglected as a source. Recent developments, however, call both these practices into question and motivate a re-evaluation of these opacities in the supernova context. Analytic formulae have recently been derived which include the full nucleon kinematics and Pauli blocking in the final state at arbitrary nucleon degeneracy. These efforts reveal that the average rate of energy transfer in $`\nu _\mu `$-nucleon scattering may surpass previous estimates by an order of magnitude . Hence, this process may compete with $`\nu _\mu `$-electron scattering as an equilibration mechanism. Similarly, estimates for the total nucleon-nucleon bremsstrahlung rate have been obtained which indicate that this process might compete with $`e^+e^{}`$ annihilation.
These results suggest that the time is ripe for a technical study of the relative importance of each process for production or thermalization. To conduct such a study, we consider $`\nu _\mu `$ neutrinos in an isotropic homogeneous thermal bath of scatterers and absorbers. In this system, the full transport problem is reduced to an evolution of the neutrino distribution function ($`_\nu `$) in energy space alone. Although this is a simplification of the true problem, it provides a theoretical laboratory in which to analyze the rates both for equilibration of an initial neutrino distribution function with dense nuclear matter and for production of the neutrinos themselves. From these rates we determine the importance and particular character of each process, and discover in which energy, temperature, or density regime each dominates. We employ a general prescription for solving the Boltzmann equation in this system with the full energy redistribution collision term. We compare quantitatively, via direct numerical evolution of an arbitrary neutrino distribution function, the rates for thermalization and production by each process, at all neutrino energies. Furthermore, we present the total nucleon-nucleon bremsstrahlung rate for arbitrary nucleon degeneracy and derive the single $`\nu _\mu `$ and $`\overline{\nu }_\mu `$ production spectra. This facilitates a more comprehensive evaluation of its relative importance in neutrino production than has previously been possible.
In §II, we discuss the general form of the Boltzmann equation and our use of it to study $`\nu _\mu `$ equilibration and production rates. In §III, we provide formulae for each of the four processes we consider: $`\nu _\mu `$-nucleon scattering, $`\nu _\mu `$-electron scattering, and $`\nu _\mu \overline{\nu }_\mu `$ pair production via both nucleon-nucleon bremsstrahlung and $`e^+e^{}`$ annihilation. In §IV, we present the results of our equilibration calculations, showing the time evolution of $`\nu _\mu `$ distribution functions as influenced by each of these processes individually. We include plots of thermalization and production rates for each process as a function of neutrino energy and time. For the scattering interactions we include figures of the time evolution of the net energy transfer to the medium as a function of incident neutrino energy. We repeat this analysis at points in temperature, density, and composition space relevant to supernovae and protoneutron stars, taken from snapshots of a stellar profile during a realistic collapse calculation . Using these results, we discuss the relative importance of each process in shaping the emergent $`\nu _\mu `$ spectrum. In §V, we recapitulate our findings and conclusions.
## II The Boltzmann Equation
The static (velocity$`=\mathrm{\hspace{0.17em}0}`$) Boltzmann equation for the evolution of the neutrino distribution function ($`_\nu `$), including Pauli blocking in the final state, and for a spherical geometry, is
$$\left(\frac{1}{c}\frac{}{t}+\mu \frac{}{r}+\frac{1\mu ^2}{r}\frac{}{\mu }\right)_\nu =(1_\nu )j_\nu _\nu \chi _\nu ,$$
(1)
where $`t`$ is the time, $`r`$ is the radial coordinate, and $`\mu (=\mathrm{cos}\theta )`$ is the cosine of the zenith angle. $`j_\nu `$ and $`\chi _\nu `$ are the total $`source`$ and $`sink`$, respectively. For emission and absorption, $`j_\nu `$ is the emissivity and $`\chi _\nu `$ is the extinction coefficient. For scattering, both $`j_\nu `$ and $`\chi _\nu `$ are energy redistribution integrals which couple one neutrino energy bin with all the others. The matrix element and associated phase-space integrations which comprise $`j_\nu `$ and $`\chi _\nu `$ for electron and nucleon scattering yield the probability that a given collision will scatter a particle into any angle or energy bin. A full transport calculation couples energy and angular bins to each other through the right hand side of eq. (1).
In a homogeneous, isotropic thermal bath of scatterers and absorbers no spatial or angular gradients exist. Consequently, the Boltzmann equation becomes
$$\frac{1}{c}\frac{_\nu }{t}=(1_\nu )j_\nu _\nu \chi _\nu .$$
(2)
By dealing with this system, the transport problem reduces to an evolution of $`_\nu `$ in just energy and time. Note that for scattering processes, both $`j_\nu `$ and $`\chi _\nu `$ require an integral over the scattered neutrino distribution function $`_\nu ^{}`$. Similarly, in evolving $`_\nu `$ via the production and absorption processes, $`j_\nu `$ and $`\chi _\nu `$ involve an integration over the anti-neutrino distribution function $`_{\overline{\nu }}`$. Therefore, $`_{\overline{\nu }}`$ must be evolved simultaneously with $`_\nu `$. While $`j_\nu `$ and $`\chi _\nu `$ may be fairly complicated integrals over phase-space, the numerical solution of eq. (2) is straightforward.
Given an arbitrary initial $`_\nu `$, we divide the relevant energy range into $`n`$ energy bins. We then solve eq. (2) for each bin individually and explicitly. Angular integrals over scattering cosines which appear in the $`\nu _\mu `$-nucleon and $`\nu _\mu `$-electron scattering formalism, as well as the electron energy integration needed for $`e^+e^{}`$ annihilation, are carried out with a 4-point Gauss-Legendre integration scheme. The double integral over dimensionless nucleon momentum variables needed to obtain the contribution from nucleon-nucleon bremsstrahlung is computed using nested 16-point Gauss-Laguerre quadratures.
### A Rates for $`_\nu `$ Evolution and Energy Transfer
Scattering, emission, and absorption processes, at a given neutrino energy ($`\epsilon _\nu `$), produce and remove neutrinos from the phase-space density at that energy. The former achieves this by transferring energy to the matter during scattering, the latter two by emitting or absorbing directly from that bin. The Boltzmann equation can then be written in terms of an in and an out channel, the former a source and the latter a sink:
$$\frac{_\nu }{t}=\frac{_\nu }{t}|_{in}\frac{_\nu }{t}|_{out}.$$
(3)
Consequently, for any interaction, there are two rates to consider: the rate for scattering or production into a given energy bin ($`\mathrm{\Gamma }_{in}`$) and the inverse rate for scattering or absorption out of that bin ($`\mathrm{\Gamma }_{out}`$). The rates $`cj_\nu `$ and $`c\chi _\nu `$ yield timescales for an interaction to occur, but fail, in the case of the former, to fold in Pauli blocking in the final state. Equation (3) includes these effects and provides natural timescales for $`_\nu `$ evolution:
$$\mathrm{\Gamma }_{in}=\frac{1}{_\nu }\frac{_\nu }{t}|_{in}=\frac{(1_\nu )}{_\nu }cj_\nu $$
(4)
and
$$\mathrm{\Gamma }_{out}=\frac{1}{_\nu }\frac{_\nu }{t}|_{out}=c\chi _\nu .$$
(5)
Note that although eq. (5) does not explicitly contain a Pauli blocking term, $`\chi _\nu `$ contains an integral over $`(1_\nu ^{})`$, in the case of scattering, and an appropriate final-state blocking term, in the case of absorption. At a given $`\epsilon _\nu `$, then, $`\mathrm{\Gamma }_{in}`$ incorporates information about the $`\nu _\mu `$ phase-space density at that energy. Conversely, at that same $`\epsilon _\nu `$, $`\mathrm{\Gamma }_{out}`$ contains information about the phase-space at all other energies. Regardless of the initial distribution, $`_\nu /t=0`$ in equilibrium. This implies $`\mathrm{\Gamma }_{in}=\mathrm{\Gamma }_{out}`$ in equilibrium and, hence, we build in a test for the degree to which the system has thermalized.
Just as there are distinct rates for the $`in`$ and $`out`$ channels of the Boltzmann equation during equilibration, so too are there distinct scattering energy transfers. For $`\nu _\mu `$ scattering with a scatterer $`s`$ (electron or nucleon), at a specific $`\epsilon _\nu `$, two thermal average energy transfers can be defined;
$$\omega _{in}=d^3p_\nu ^{}\omega _\nu ^{}^{in}\left[\nu _\mu s\nu _\mu ^{}s^{}\right]/d^3p_\nu ^{}_\nu ^{}^{in}\left[\nu _\mu s\nu _\mu ^{}s^{}\right]$$
(6)
and
$$\omega _{out}=d^3p_\nu ^{}\omega (1_\nu ^{})^{out}\left[\nu _\mu s\nu _\mu ^{}s^{}\right]/d^3p_\nu ^{}(1_\nu ^{})^{out}\left[\nu _\mu s\nu _\mu ^{}s^{}\right],$$
(7)
where primes denote the scattered neutrino, $`\omega (=\epsilon _\nu \epsilon _\nu ^{})`$ is the energy transfer, and $`^{in}`$ and $`^{out}`$ are the kernels for scattering into and out of a given energy bin, respectively. As a consequence of detailed balance between the in and out channels of the Boltzmann equation, $`^{in}=e^{\beta \omega }^{out}`$, where $`\beta =1/k_BT`$ and $`T`$ is the matter temperature. (The scattering kernels are discussed in detail in §III for both scattering processes.) Note that the denominators in eqs. (6) and (7), up to constants which divide out in the definitions of $`\omega _{in}`$ and $`\omega _{out}`$, are just $`j_\nu `$ and $`\chi _\nu `$, respectively.
In an effort to provide more than one measure of the timescale for $`_\nu `$ equilibration due to scattering and to make contact with previous neutrino thermalization studies we also define a set of timescales in terms of $`\omega _{out}`$ and the higher $`\omega `$-moment, $`\omega ^2_{out}`$;
$$\mathrm{\Gamma }_D=c\chi _\nu \left|\frac{\omega _{out}}{\epsilon _\nu }\right|$$
(8)
and
$$\mathrm{\Gamma }_E=c\chi _\nu \frac{\omega ^2_{out}}{\epsilon _\nu ^2}.$$
(9)
$`\mathrm{\Gamma }_D`$ is the rate for shifting the centroid of a given distribution and $`\mathrm{\Gamma }_E`$ is the rate for spreading an initial distribution . In contrast with the work of , we include the full effects of Pauli blocking in the final state, allowing us to deal consistently with cases in which the $`\nu _\mu `$s are partially degenerate.
## III Individual Interactions
This section details the source and sink terms necessary to solve the Boltzmann equation for the time-evolution of $`_\nu `$. Sections §III A and §III B are dedicated to the presentation and discussion of the collision terms for $`\nu _\mu `$-nucleon and $`\nu _\mu `$-electron scattering, respectively. Section §III C describes the Legendre series expansion approximation and the use of it to compute the contribution to the Boltzmann equation, the pair emissivity, and the single $`\nu _\mu `$ spectrum due to $`e^+e^{}\nu _\mu \overline{\nu }_\mu `$. Our derivations of $`j_\nu `$ and $`\chi _\nu `$, as well as the single and pair spectra from nucleon-nucleon bremsstrahlung at arbitrary nucleon degeneracy and in the non-degenerate limit, are presented in §III D. In what follows, we take $`G^21.55\times 10^{33}`$ cm<sup>3</sup> MeV<sup>-2</sup> s<sup>-1</sup>, $`\mathrm{sin}^2\theta _W0.231`$, and employ natural units in which $`\mathrm{}=c=k_B=1`$.
### A Nucleon Scattering: $`\nu _\mu n\nu _\mu n`$ and $`\nu _\mu p\nu _\mu p`$
Researchers working on supernova and protoneutron star evolution have recently re-evaluated the issue of energy transfer via $`\nu _\mu `$-nucleon scattering . Originally, the assumption was made that the nucleons were stationary . If a neutron of mass $`m_n`$ is at rest with respect to an incoming neutrino of energy $`\epsilon _\nu `$, one finds that the energy transfer ($`\overline{\omega }`$) is $`\epsilon _\nu ^2/m_n`$. For $`\epsilon _\nu =\mathrm{\hspace{0.17em}10}`$ MeV, $`\overline{\omega }0.1`$ MeV, a fractional energy lost of 1%. Using these simple kinematic arguments and disregarding neutrino and nucleon Pauli blocking, one finds that the thermalization rate for $`\nu _\mu `$-electron scattering should be approximately a factor of 20 larger than that for $`\nu _\mu `$-nucleon scattering. In the context of interest, however, at temperatures of order 10 MeV and mass densities of order $`10^{13}`$ g cm<sup>-3</sup>, free nucleons are not stationary, but have thermal velocities. The fractional energy exchange per collision, in the case of $`\nu _\mu `$-neutron scattering, is then $`p_n/m_nc`$ . For $`T10`$ MeV this gives a $``$10$``$20% change in $`\epsilon _\nu `$ per collision. This calls the naive estimate of the relative importance of $`\nu _\mu `$-nucleon scattering as a thermalization process into question and a more complete exploration of the relative importance of the two scattering processes is necessary.
Recently, analytic formulae have been derived which include the full kinematics of $`\nu _\mu `$-nucleon scattering at arbitrary nucleon degeneracy . At the temperatures and densities encountered in the supernova context non-interacting nucleons are not relativistic. Due to nucleon-nucleon interactions, however, at and around nuclear density ($`2.68\times 10^{14}`$ g cm<sup>-3</sup>), the nucleon’s effective mass drops and is expected to be comparable with its Fermi momentum . In such a circumstance, a relativistic description of the $`\nu _\mu `$-nucleon scattering interaction is warranted. In addition, spin and density correlation effects engendered by these nucleon-nucleon interactions have been found to suppress the $`\nu _\mu `$-nucleon interaction rate by as much as a factor of $`23`$ .
In this study, we focus on $`\nu _\mu `$ equilibration rates at densities $`1\times 10^{14}`$ g cm<sup>-3</sup> where it is still unclear if nucleon-nucleon interactions will play an important role. This ambiguity is due in part to uncertainties both in the nuclear equation of state and the nucleon-nucleon interaction itself. For this reason we choose to treat the nucleons as non-relativistic and non-interacting, thereby ignoring collective effects which might enhance or reduce the $`\nu _\mu `$-nucleon scattering rate. Making these assumptions, we find that $`j_\nu `$ and $`\chi _\nu `$ in eq. (2) are given by
$$j_\nu =\frac{G^2}{(2\pi )^3}d^3\stackrel{}{p}_\nu ^{}_{\mathrm{NC}}_\nu ^{}e^{\beta \omega }$$
(10)
and
$$\chi _\nu =\frac{G^2}{(2\pi )^3}d^3\stackrel{}{p}_\nu ^{}_{\mathrm{NC}}[1_\nu ^{}],$$
(11)
where $`\beta =1/T`$, $`\stackrel{}{p}_\nu ^{}`$ is the final state neutrino momentum, and $`\omega `$ is the energy transfer. In eqs. (10) and (11), the neutral-current scattering kernel is given by
$$_{\mathrm{NC}}=S(q,\omega )[(1+\mu )V^2+(3\mu )A^2],$$
(12)
where $`\mu (=\mathrm{cos}\theta )`$ is the cosine of the scattering angle between incident and final state neutrinos and $`S(q,\omega )`$ is the dynamic structure function. In eq. (12), $`V`$ and $`A`$ are the appropriate vector and axial-vector coupling constants; for $`\nu _\mu `$-neutron scattering, $`V=1/2`$ and $`A=1.26/2`$. The dynamic structure function is
$`S(q,\omega )`$ $`=`$ $`2{\displaystyle \frac{d^3\stackrel{}{p}}{(2\pi )^3}(1^{})(2\pi )\delta (\omega +\epsilon \epsilon ^{})}`$ (13)
$`=`$ $`2\mathrm{Im}\mathrm{\Pi }^{(0)}(q,\omega )(1e^{\beta \omega })^1`$ (14)
where $`q=|p_\nu p_\nu ^{}|=[\epsilon _\nu ^2+\epsilon _\nu ^{\mathrm{\hspace{0.17em}2}}2\epsilon _\nu \epsilon _\nu ^{}\mu ]^{1/2}`$ is the magnitude of the momentum transfer, and $``$ and $`^{}`$ are the incident and scattered nucleon distribution functions, respectively. In eq. (14), $`\stackrel{}{p}`$ is the incident nucleon momentum, $`\epsilon `$ is the incident nucleon energy, and $`\epsilon ^{}`$ is the scattered nucleon energy. The imaginary part of the free polarization is given by
$$\mathrm{Im}\mathrm{\Pi }^{(0)}(q,\omega )=\frac{m^2}{2\pi \beta q}\mathrm{ln}\left[\frac{1+e^{Q^2+\eta }}{1+e^{Q^2+\eta \beta \omega }}\right],$$
(15)
where
$$Q=\left(\frac{m\beta }{2}\right)^{1/2}\left(\frac{\omega }{q}+\frac{q}{2m}\right),$$
(16)
$`\eta `$ is the nucleon degeneracy ($`\mu /T`$), and $`m`$ is the nucleon mass. The factor $`e^{\beta \omega }`$ which appears in eq. (10) is a consequence of the fact that $`S(q,\omega )=e^{\beta \omega }S(q,\omega )`$, itself a consequence of detailed balance between the $`in`$ and $`out`$ channels of the Boltzmann equation. The dynamic structure function can be thought of as a correlation function which connects $`\epsilon _\nu `$ and $`\epsilon _\nu ^{}`$.
The $`\varphi `$ angular integrations implicit in eqs. (10) and (11) can be computed trivially assuming the isotropy of $`_\nu `$. Furthermore, defining the coordinate system with the momentum vector of the incident neutrino, the scattering angle and the direction cosine are equivalent. Combining these two equations in the Boltzmann equation for the evolution of $`_\nu `$ due to neutral-current $`\nu _\mu `$-nucleon scattering, we obtain
$$\frac{_\nu }{t}=\frac{G^2}{(2\pi )^2}_0^{\mathrm{}}𝑑\epsilon _\nu ^{}\epsilon _\nu ^{\mathrm{\hspace{0.17em}2}}_1^1𝑑\mu _{\mathrm{NC}}\left\{[1_\nu ]_\nu ^{}e^{\beta \omega }_\nu [1_\nu ^{}]\right\}.$$
(17)
### B Electron Scattering: $`\nu _\mu e^{}\nu _\mu e^{}`$
At the temperatures and densities encountered in supernovae and protoneutron stars, electrons are highly relativistic. A formalism analogous to that used for $`\nu _\mu `$-nucleon scattering is desired in order to include the full electron kinematics at arbitrary electron degeneracy. Reddy et al. have developed a relativistic generalization of the structure function formalism described in §III A. They obtain a set of polarization functions which characterize the relativistic medium’s response to a neutrino probe in terms of polylogarithmic functions. In analogy with eq. (17), we can write the Boltzmann equation for the evolution of $`_\nu `$ due to $`\nu _\mu `$-electron scattering, as
$$\frac{_\nu }{t}=\frac{G^2}{(4\pi )^3}d^3p_\nu ^{}_{\mathrm{NC}}^r\left\{[1_\nu ]_\nu ^{}e^{\beta \omega }_\nu [1_\nu ^{}]\right\},$$
(18)
where $`_{\mathrm{NC}}^r`$ is the relativistic neutral-current scattering kernel for $`\nu _\mu `$s, analogous to $`_{\mathrm{NC}}`$ in eq. (12). All the physics of the interaction is contained in $`_{NC}^r`$, which can be written as
$$_{\mathrm{NC}}^r=\mathrm{Im}\{\mathrm{\Lambda }^{\alpha \beta }\mathrm{\Pi }_{\alpha \beta }^R\}(1e^{\beta \omega })^1.$$
(19)
As in the non-relativistic case, $`_{\mathrm{NC}}^r`$ is composed of the lepton tensor,
$$\mathrm{\Lambda }^{\alpha \beta }=8[2k^\alpha k^\beta +(kq)g^{\alpha \beta }(k^\alpha q^\beta +q^\alpha k^\beta )iϵ^{\alpha \beta \mu \nu }k^\mu q^\nu ],$$
(20)
which is just the squared and spin-summed matrix element for the scattering process written in terms of $`k_\alpha `$, the incident $`\nu _\mu `$ four-momentum, and $`q_\alpha \left(=(\omega ,\stackrel{}{q})\right)`$, the four-momentum transfer. The scattering kernel also contains the retarded polarization tensor, $`\mathrm{\Pi }_{\alpha \beta }^R`$, which is directly analogous to the free polarization in the non-relativistic case given in eq. (14). The retarded polarization tensor is related to the causal polarization by
$$\mathrm{Im}\mathrm{\Pi }_{\alpha \beta }^R=\mathrm{tanh}\left(\frac{1}{2}\beta \omega \right)\mathrm{Im}\mathrm{\Pi }_{\alpha \beta }$$
(21)
and
$$\mathrm{\Pi }_{\alpha \beta }=i\frac{d^4p}{(2\pi )^4}\mathrm{Tr}[G_e(p)J_\alpha G_e^{}(p+q)J_\beta ].$$
(22)
In eq. (22), $`p_\alpha `$ is the electron four-momentum and $`J_\alpha `$ is the current operator. The electron Green’s functions ($`G_e`$ and $`G_e^{}`$), explicit in the free polarization, connect points in electron energy space and characterize the effect of the interaction on relativistic electrons. The polarization tensor can be written in terms of a vector part, an axial-vector part, and a mixed part, so that
$$\mathrm{\Pi }_{\alpha \beta }=V^2\mathrm{\Pi }_{\alpha \beta }^V+A^2\mathrm{\Pi }_{\alpha \beta }^A2VA\mathrm{\Pi }_{\alpha \beta }^{VA}.$$
(23)
In turn, the vector part of the polarization tensor can be written in terms of two independent components, $`\mathrm{\Pi }_T`$ and $`\mathrm{\Pi }_L`$. In contrast with eq. (12), since $`v/c1`$ for the electrons, the angular terms which were dropped from the matrix element in the non-relativistic case, leading to a single structure function, must now be retained. $`_{\mathrm{NC}}^r`$ can then be written as a set of three structure functions :
$$_{\mathrm{NC}}^r=8[A𝒮_1(q,\omega )+𝒮_2(q,\omega )+B𝒮_3(q,\omega )](1e^{\beta \omega })^1,$$
(24)
where $`A=(4\epsilon _\nu \epsilon _\nu ^{}+q_\alpha ^2)/2q^2`$ and $`B=\epsilon _\nu +\epsilon _\nu ^{}`$. These structure functions can be written in terms of the vector parts of the retarded polarization tensor ($`\mathrm{\Pi }_T^R`$ and $`\mathrm{\Pi }_L^R`$), the axial part ($`\mathrm{\Pi }_A^R`$), and the mixed part ($`\mathrm{\Pi }_{VA}^R`$):
$$𝒮_1(q,\omega )=(V^2+A^2)\left[\mathrm{Im}\mathrm{\Pi }_L^R(q,\omega )+\mathrm{Im}\mathrm{\Pi }_T^R(q,\omega )\right],$$
(25)
$$𝒮_2(q,\omega )=(V^2+A^2)\mathrm{Im}\mathrm{\Pi }_T^R(q,\omega )A^2\mathrm{Im}\mathrm{\Pi }_A^R(q,\omega ),$$
(26)
and
$$𝒮_3(q,\omega )=\mathrm{\hspace{0.17em}2}VA\mathrm{Im}\mathrm{\Pi }_{VA}^R(q,\omega ).$$
(27)
The retarded polarization functions, in terms of differences between polylogarithmic integrals, can be found in Appendix A.
### C Electron-Positron Annihilation: $`e^+e^{}\nu _\mu \overline{\nu }_\mu `$
Fermi’s Golden Rule for the total volumetric emission rate for the production of $`\nu _\mu `$s via electron-positron annihilation can be written as
$$Q=\frac{d^3\stackrel{}{p}}{(2\pi )^32\epsilon }\frac{d^3\stackrel{}{p}^{}}{(2\pi )^32\epsilon ^{}}\frac{d^3\stackrel{}{q}_\nu }{(2\pi )^32\epsilon _\nu }\frac{d^3\stackrel{}{q}_{\overline{\nu }}}{(2\pi )^32\epsilon _{\overline{\nu }}}\epsilon _\nu \left(\frac{1}{4}\underset{s}{}|^2|\right)(2\pi )^4\delta ^4(\mathrm{P})\mathrm{\Xi }[],$$
(28)
where
$$\mathrm{\Xi }[]=(1_\nu )(1_{\overline{\nu }})_e^{}_{e^+},$$
(29)
and $`\delta ^4(\mathrm{P})`$ conserves four-momentum. In eq. (28), $`p_\alpha (=(\epsilon ,\stackrel{}{p}))`$ and $`p_\alpha ^{}(=(\epsilon ^{},\stackrel{}{p}^{}))`$ are the four-momenta of the electron and positron, respectively, and $`q_\nu ^\alpha (=(\epsilon _\nu ,\stackrel{}{q_\nu }))`$ and $`q_{\overline{\nu }}^\alpha (=(\epsilon _{\overline{\nu }},\stackrel{}{q}_{\overline{\nu }}))`$ are the four-momenta of the $`\nu _\mu `$ and $`\overline{\nu }_\mu `$, respectively. The process of electron-positron annihilation into a neutrino/anti-neutrino pairs is related to neutrino-electron scattering considered in §III B via a crossing symmetry. In order to make the problem tractable, we follow the standard procedure of expanding the production kernel in a Legendre series in the scattering angle to first order (see Appendix B). Near the neutrinospheres, at densities which render neutrino transport diffusive this approximation holds. In a full neutrino transport algorithm, however, which must handle both the diffusion and free-streaming limits, the second-order term, with proper closure relations, must be used in the semi-transparent regime between the neutrinospheres and the shock . Having made this approximation, including only the zeroth- and first-order terms, the single $`\nu _\mu `$ spectrum is
$$\frac{dQ}{d\epsilon _\nu }=(1_\nu )\frac{\epsilon _\nu ^3}{8\pi ^4}_0^{\mathrm{}}𝑑\epsilon _{\overline{\nu }}\epsilon _{\overline{\nu }}^2\mathrm{\Phi }_0^p(\epsilon _\nu ,\epsilon _{\overline{\nu }})(1_{\overline{\nu }}),$$
(30)
where $`\mathrm{\Phi }_0^p(\epsilon _\nu ,\epsilon _{\overline{\nu }})`$ is the zeroth-order production kernel expansion coefficient, an integral over the electron energy (see Appendix B) . With the differential spectrum or emissivity ($`dQ/d\epsilon _\nu `$) in hand, it is a simple matter to extract the contribution to the Boltzmann equation due to $`e^+e^{}`$ annihilation. As eq. (30) already contains the $`\nu _\mu `$ blocking factor, the contribution to the Boltzmann equation, the $`in`$ channel explicit in eq. (3), can be written as
$$\frac{_\nu }{t}|_{in}=\frac{1}{4\pi }\frac{(2\pi )^3}{\epsilon _\nu ^3}\frac{dQ}{d\epsilon _\nu }.$$
(31)
In order to obtain the $`out`$ channel for absorption due to $`e^+e^{}`$ annihilation, we need only replace $`_e^{}_{e^+}`$ in eq. (29) with an electron/positron blocking term, $`(1_e^{})(1_{e^+})`$, and replace the $`\nu _\mu `$ and $`\overline{\nu }_\mu `$ blocking terms in eq. (30) with $`_\nu _{\overline{\nu }}`$. Finally, the Boltzmann equation for the evolution of $`_\nu `$ in time due to $`e^+e^{}\nu _\mu \overline{\nu }_\mu `$ can be written as
$$\frac{_\nu }{t}=\frac{2G^2}{(2\pi )^3}_0^{\mathrm{}}𝑑\epsilon _{\overline{\nu }}\epsilon _{\overline{\nu }}^2_0^ϵ𝑑\epsilon H_o(\epsilon _\nu ,\epsilon _{\overline{\nu }},\epsilon )\left\{(1_\nu )(1_{\overline{\nu }})_e^{}_{e^+}_\nu _{\overline{\nu }}(1_e^{})(1_{e^+})\right\},$$
(32)
where $`ϵ=\epsilon _\nu +\epsilon _{\overline{\nu }}`$ and $`H_o(\epsilon _\nu ,\epsilon _{\overline{\nu }},\epsilon )`$ is given in eq. (B5). In solving eq. (32), $`_{\overline{\nu }}`$ must be evolved simultaneously with $`_\nu `$. To do so, in addition to making the appropriate changes to the vector and axial-vector coupling constants, $`V`$ and $`A`$, one needs to integrate over $`\epsilon _\nu `$ instead of $`\epsilon _{\overline{\nu }}`$. Note that the electron and positron distribution functions appear explicitly in eq. (32). We take these distributions to be Fermi-Dirac at temperature $`T`$ and with $`\eta _e`$ determined by $`T`$, $`\rho `$, and $`Y_e`$.
Equation (28) may also be used to find the total volumetric $`\nu _\mu \overline{\nu }_\mu `$ pair spectrum by replacing $`\epsilon _\nu `$ in the numerator with $`ϵ`$. Ignoring neutrino blocking in the final state one can show that
$$Q_{\nu _\mu \overline{\nu }_\mu }\mathrm{\hspace{0.17em}2.09}\times 10^{24}\left(\frac{T}{\mathrm{MeV}}\right)^9f(\eta _e)\mathrm{ergs}\mathrm{cm}^3\mathrm{s}^1,$$
(33)
where
$$f(\eta _e)=\frac{F_4(\eta _e)F_3(\eta _e)+F_4(\eta _e)F_3(\eta _e)}{2F_4(0)F_3(0)},$$
(34)
and
$$F_n(y)=_0^{\mathrm{}}\frac{x^n}{e^{xy}+1}𝑑x$$
(35)
are the Fermi integrals.
### D Nucleon-Nucleon Bremsstrahlung
The importance of nucleon-nucleon bremsstrahlung in late-time neutron star cooling has been acknowledged for some time . Recently, however, this process has received more attention as a contributor of $`\nu _\mu \overline{\nu }_\mu `$ pairs and as an energy transport mechanism in both core-collapse supernova and nascent neutron star evolution . The contribution from nucleon-nucleon bremsstrahlung is a composite of neutron-neutron ($`nn`$), proton-proton ($`pp`$), and neutron-proton ($`np`$) bremsstrahlung. Fermi’s Golden Rule for the total volumetric emissivity of single $`\nu _\mu `$s due to $`nn`$, $`pp`$, or $`np`$ bremsstrahlung, including $`\nu _\mu `$ and $`\overline{\nu }_\mu `$ blocking in the final state, is given by
$$Q=\left[\underset{i=1}{\overset{4}{}}\frac{d^3\stackrel{}{p}_i}{(2\pi )^3}\right]\frac{d^3\stackrel{}{q}_\nu }{(2\pi )^32\epsilon _\nu }\frac{d^3\stackrel{}{q}_{\overline{\nu }}}{(2\pi )^32\epsilon _{\overline{\nu }}}\epsilon _\nu \left(s||^2\right)(2\pi )^4\delta ^4(\mathrm{P})\mathrm{\Xi }[]$$
(36)
where
$$\mathrm{\Xi }[]=_1_2(1_3)(1_4)(1_\nu )(1_{\overline{\nu }}).$$
(37)
The product of differential phase space factors in eq. (36) includes a term for each of the four nucleons involved in the process; 1 and 2 denote initial-state nucleons whereas 3 and 4 denote final-state nucleons. In eq. (36), $`s`$ is a symmetry factor for identical initial-state fermions, $`\stackrel{}{q}_\nu `$ is the neutrino three-momentum, $`\epsilon _\nu `$ is the neutrino energy, and the four-momentum conserving delta function is explicit. In a one-pion exchange model for the nucleon-nucleon interaction, the spin-summed matrix element can be approximated by
$$\underset{s}{}||^264G^2g_A^2\left(\frac{f}{m_\pi }\right)^4\left[\left(\frac{k^2}{k^2+m_\pi ^2}\right)^2+\mathrm{}\right]ϵ^2(\epsilon _\nu \epsilon _{\overline{\nu }}\stackrel{}{q}_\nu \widehat{k}\stackrel{}{q}_{\overline{\nu }}\widehat{k})$$
(38)
where $`ϵ=\epsilon _\nu +\epsilon _{\overline{\nu }}`$, $`k`$ is the magnitude of the nucleon momentum transfer, $`g_A1.26`$, $`f1`$ is the pion-nucleon coupling, and $`m_\pi `$ is the mass of the pion. In order to make the 18-dimensional phase-space integration in eq. (36) tractable we assume the quantity in square brackets to be of order unity, but possibly as low as 0.1 . To acknowledge our ignorance, we introduce the factor, $`\zeta `$, and assume these momentum terms are constant. Furthermore, we neglect the momentum of the neutrinos relative to the momentum of the nucleons. We are left with a simple, but general, form for the bremsstrahlung matrix element:
$$||^2A\zeta \frac{\epsilon _\nu \epsilon _{\overline{\nu }}}{ϵ^2},$$
(39)
where $`A=64G^2g_A^2f^4/m_\pi ^4`$. In the case of $`nn`$ or $`pp`$ bremsstrahlung, as appropriate for identical particles in the initial state, the symmetry factor ($`s`$) in eq. (36) is $`1/4`$. Such a symmetry factor does not enter for the mixed-nucleon process, $`np`$, which is still further enhanced by the fact that a charged pion mediates the nucleon exchange . This increases the matrix element in eq. (39) by a factor of $`7/3`$ in the degenerate nucleon limit and $`5/2`$ in the non-degenerate limit . Considering the already substantial simplifications made by choosing not to handle the momentum terms directly, we will adopt the more conservative $`4\times (7/3)`$ enhancement for the $`np`$ matrix element. The total volumetric emission rate combining all processes is just $`Q_{tot}=Q_{nn}+Q_{pp}+Q_{np}`$. What remains is to reduce eq. (36) to a useful expression in asymmetric matter and at arbitrary neutron and proton degeneracy.
Following ref. , we define new momenta, $`p_\pm =(p_1\pm p_2)/2`$ and $`p_{3c,4c}=p_{3,4}p_+`$, new direction cosines, $`\gamma _1=p_+p_{}/|p_+||p_{}|`$ and $`\gamma _c=p_+p_{3c}/|p_+||p_{3c}|`$, and let $`u_i=p_i^2/2mT`$. Furthermore, we note that $`d^3p_1d^3p_2=8d^3p_+d^3p_{}`$. Using the three-momentum conserving delta function, we can do the $`d^3\stackrel{}{p}_4`$ integral trivially. Rewriting eq. (36) with these definitions, we find that
$$Q=2As\zeta (2mT)^{9/2}(2\pi )^9𝑑\epsilon _\nu \epsilon _\nu ^3𝑑\epsilon _{\overline{\nu }}𝑑u_{}𝑑u_+𝑑u_{3c}𝑑\gamma _1𝑑\gamma _c(\epsilon _{\overline{\nu }}/ϵ)^2(u_{}u_+u_{3c})^{1/2}\delta (E)\mathrm{\Xi }[],$$
(40)
where
$$\delta (E)=\delta (\underset{i=1}{\overset{4}{}}\epsilon _iϵ)=\delta (2T(u_{}u_{3c}ϵ/2T)).$$
(41)
The nucleon distribution functions in the term $`\mathrm{\Xi }[]`$ in eq. (40) have been rewritten in terms of the new direction cosines, the dimensionless momenta ($`u_i`$), and the initial-state nucleon degeneracy factors $`\eta _{1,2}=\mu _{1,2}/T`$:
$$_1=\frac{e^{(a_1^{}+b^{}\gamma _1)}}{2\mathrm{cosh}(a_1^{}+b^{}\gamma _1)}\mathrm{and}_2=\frac{e^{(a_2^{}b^{}\gamma _1)}}{2\mathrm{cosh}(a_2^{}b^{}\gamma _1)},$$
(42)
where $`a_{1,2}^{}=a_{1,2}/2=\frac{1}{2}(u_++u_{}\eta _{1,2})`$ and $`b^{}=b/2=(u_+u_{})^{1/2}`$. Furthermore,
$$(1_3)=\frac{e^{(c_1^{}+d^{}\gamma _c)}}{2\mathrm{cosh}(c_1^{}+d^{}\gamma _c)}\mathrm{and}(1_4)=\frac{e^{(c_2^{}d^{}\gamma _c)}}{2\mathrm{cosh}(c_2^{}d^{}\gamma _c)},$$
(43)
where $`c_{1,2}^{}=c_{1,2}/2=\frac{1}{2}(u_++u_{3c}\eta _{1,2})`$ and $`d^{}=d/2=(u_+u_{3c})^{1/2}`$. $`_\nu `$ and $`_{\overline{\nu }}`$, in contrast with the nucleon distribution functions, are independent of angle; for a given set of thermodynamic conditions, they remain functions of energy alone. While non-trivial, the integrations over $`\gamma _1`$ and $`\gamma _c`$ can be performed. For example, the result for the $`\gamma _1`$ integration is of the form
$$\frac{1}{2\sqrt{B(B+1)}}\mathrm{ln}[(B(1+2B)\xi ^2+2\xi \sqrt{B}(B+1)(\xi ^21)],$$
(44)
where $`B=\mathrm{sinh}^2a^{}`$ and $`\xi =\mathrm{cosh}b^{}\gamma _1`$. With a proper evaluation of the integration limits and some algebra one can rewrite this result as
$$\frac{1}{2\mathrm{sinh}a^{}\mathrm{cosh}a^{}}\mathrm{ln}\left[\frac{(1+\mathrm{cosh}a\mathrm{cosh}b+\mathrm{sinh}a\mathrm{sinh}b)}{(1+\mathrm{cosh}a\mathrm{cosh}b\mathrm{sinh}a\mathrm{sinh}b)}\right].$$
(45)
Similar operations yield a result for the $`\gamma _c`$ integral in terms of $`c`$ and $`d`$. In addition, eq. (41) can be used to eliminate the integral over $`u_{}`$. Collectively, these manipulations reveal that the differential $`\nu _\mu `$ bremsstrahlung emissivity at arbitrary neutron and proton degeneracy is simply a three-dimensional integral over $`u_+`$, $`u_{3c}`$, and $`\epsilon _{\overline{\nu }}`$:
$$\frac{dQ}{d\epsilon _\nu }=Ks\zeta (1_\nu )\epsilon _\nu ^3𝑑\epsilon _{\overline{\nu }}𝑑u_+𝑑u_{3c}(\epsilon _{\overline{\nu }}/ϵ)^2u_+^{1/2}e^{\beta ϵ/2}\mathrm{\Phi }(ϵ,u_+,u_{3c})(1_{\overline{\nu }}),$$
(46)
where
$$K=2G^2\left(\frac{m}{2\pi ^2}\right)^{9/2}\left(\frac{f}{m_\pi }\right)^4g_A^2T^{7/2},$$
(47)
$`\mathrm{\Phi }(ϵ,u_+,u_{3c})`$ $`=`$ $`\mathrm{sinh}^1(f)\mathrm{ln}\left[\left({\displaystyle \frac{1+\mathrm{cosh}(e_+)}{1+\mathrm{cosh}(e_{})}}\right)\left({\displaystyle \frac{\mathrm{cosh}(f)+\mathrm{cosh}(g_+)}{\mathrm{cosh}(f)+\mathrm{cosh}(g_{})}}\right)\right]`$ (48)
$`\times `$ $`\mathrm{sinh}^1(j)\mathrm{ln}\left[\left({\displaystyle \frac{1+\mathrm{cosh}(h_+)}{1+\mathrm{cosh}(h_{})}}\right)\left({\displaystyle \frac{\mathrm{cosh}(j)+\mathrm{cosh}(k_+)}{\mathrm{cosh}(j)+\mathrm{cosh}(k_{})}}\right)\right],`$ (49)
and
$`e_\pm `$ $`=`$ $`(u_+^{1/2}\pm u_{}^{1/2})^2\eta _2`$ (50)
$`f`$ $`=`$ $`u_++u_{}\eta _1/2\eta _2/2`$ (51)
$`g_\pm `$ $`=`$ $`\pm 2(u_+u_{})^{1/2}\eta _1/2+\eta _2/2`$ (52)
$`h_\pm `$ $`=`$ $`(u_+^{1/2}\pm u_{3c}^{1/2})^2\eta _2`$ (53)
$`j`$ $`=`$ $`u_++u_{3c}\eta _1/2\eta _2/2`$ (54)
$`k_\pm `$ $`=`$ $`\pm 2(u_+u_{3c})^{1/2}\eta _1/2+\eta _2/2.`$ (55)
Though $`u_{}`$ has been integrated out via the energy-conserving delta function, it appears here in an attempt to make this expression more compact and should be read as $`u_{}=u_{3c}+ϵ/2T`$. Importantly, if $`\eta _1=\eta _2`$ the right-hand term within both logarithmic terms in $`\mathrm{\Phi }(ϵ,u_+,u_{3c})`$ becomes unity.
Using eq. (31), we can easily obtain the contribution to the Boltzmann equation due to nucleon-nucleon bremsstrahlung for arbitrary nucleon degeneracy, in asymmetric matter, and including the full nucleon and neutrino Pauli blocking terms. We find that
$$j_\nu =K^{}s\zeta 𝑑\epsilon _{\overline{\nu }}𝑑u_+𝑑u_{3c}(\epsilon _{\overline{\nu }}/ϵ)^2u_+^{1/2}e^{\beta ϵ/2}\mathrm{\Phi }(ϵ,u_+,u_{3c})(1_{\overline{\nu }})$$
(56)
where $`K^{}=[(2\pi )^3/4\pi ]K`$. The nucleon phase-space integrations above are identical in form for the $`\nu _\mu \overline{\nu }_\mu `$ absorption process, $`\nu _\mu \overline{\nu }_\mu nnnn`$. In this case, then, the primed energies are now associated with nucleons 1 and 2 in the above manipulations and the incident nucleons (3 and 4) have unprimed energies. If we take the form derived above for the nucleon phase-space terms, the absorption channel ($`\chi _\nu `$) must then contain a factor of $`e^{\beta ϵ}`$. In addition, the blocking term, $`(1_{\overline{\nu }})`$, becomes $`_{\overline{\nu }}`$. The Boltzmann equation for the evolution of $`_\nu `$ in time is then,
$$\frac{1}{c}\frac{_\nu }{t}=K^{}s\zeta 𝑑\epsilon _{\overline{\nu }}𝑑u_+𝑑u_{3c}(\epsilon _{\overline{\nu }}/ϵ)^2u_+^{1/2}e^{\beta ϵ/2}\mathrm{\Phi }(ϵ,u_+,u_{3c})\left\{(1_\nu )(1_{\overline{\nu }})_\nu _{\overline{\nu }}e^{\beta ϵ}\right\}.$$
(57)
For the neutron-neutron ($`nn`$) or proton-proton ($`pp`$) bremsstrahlung contribution, we simply set $`s=1/4`$ in eq. (57) and use $`\eta _1=\eta _2=\eta _n`$ or $`\eta _1=\eta _2=\eta _p`$, respectively. For the mixed nucleon ($`np`$) bremsstrahlung we set $`s=1`$, multiply eq. (57) by $`7/3`$, and set $`\eta _1=\eta _n`$ and $`\eta _2=\eta _p`$. While eqs. (46) and (57) may not appear symmetric in $`\eta _1`$ and $`\eta _2`$ the logarithmic terms conspire to ensure that the rates for both $`np`$ and $`pn`$ bremsstrahlung are identical, as they should be. That is, it makes no difference whether we set $`\eta _n`$ or $`\eta _p`$ equal to $`\eta _1`$ or $`\eta _2`$.
Just as in §III C, in considering $`e^+e^{}\nu _\mu \overline{\nu }_\mu `$, $`_{\overline{\nu }}`$ must be evolved simultaneously with $`_\nu `$. In this case, however, the situation is simpler. Suppose we wish to compare electron-positron annihilation with nucleon-nucleon bremsstrahlung by starting at $`t=0`$ with $`_{\overline{\nu }}=_\nu =0`$ over all energies. We then solve eq. (32) and its $`_{\overline{\nu }}`$ counterpart at each timestep and at each energy. For $`e^+e^{}`$ annihilation, $`_\nu `$ and $`_{\overline{\nu }}`$ will evolve differently; they will be visibly different at each timestep, because of the weighting of the vector and axial-vector coupling constants which appear in the matrix element. In contrast, eq. (57) for bremsstrahlung must be solved only once. Since there is no difference in weighting between $`\nu _\mu `$ and $`\overline{\nu }_\mu `$, we can set $`_{\overline{\nu }}=_\nu `$ at every energy, at every timestep, as long as $`_{\overline{\nu }}=_\nu `$ at $`t=0`$. Of course, if we wish to consider $`_{\overline{\nu }}_\nu `$ initially, the two distributions would need to be evolved separately and simultaneously, coupled through the blocking and source terms on the right-hand side of the Boltzmann equation.
Equation (36) can also be used to find the total volumetric $`\nu _\mu \overline{\nu }_\mu `$ pair emissivity. To facilitate this we replace $`\epsilon _\nu `$ with $`ϵ`$ and insert $`\delta (ϵ(\epsilon _\nu +\epsilon _{\overline{\nu }}))𝑑ϵ`$. Assuming the neutrinos are radiated isotropically, we can use this delta function to do the integral over $`d^3\stackrel{}{q}_{\overline{\nu }}`$ and leave the total rate in terms of an integral over $`\epsilon _\nu `$ from zero to $`ϵ`$ and another over $`ϵ`$ from zero to infinity. Momentarily ignoring neutrino blocking in the final state, the former can be integrated easily. Making the same momentum, angle, and nucleon distribution function substitutions we used in deriving the single $`\nu _\mu `$ spectrum we can reduce the pair spectrum to an integral over $`u_+`$, $`u_{3c}`$, and $`q=ϵ/2T`$. We find that
$$Q_{\nu _\mu \overline{\nu }_\mu }=Ds\zeta T^{8.5}𝑑q𝑑u_{3c}𝑑u_+q^4e^qu_+^{1/2}\mathrm{\Phi }(ϵ,u_+,u_{3c}),$$
(58)
where
$$D=\frac{8}{15}\frac{G^2g_A^2}{\sqrt{2}\pi ^9}\left(\frac{f}{m_\pi }\right)^4m^{9/2},$$
(59)
and $`\mathrm{\Phi }(ϵ,u_+,u_{3c})`$ is defined in eq. (49). Note that eq. (58) allows us to easily calculate the pair differential volumetric emissivity ($`dQ_{\nu _\mu \overline{\nu }_\mu }/dϵ`$). For $`Q_{\nu _\mu \overline{\nu }_\mu }^{nn}`$ and $`Q_{\nu _\mu \overline{\nu }_\mu }^{pp}`$, $`s=1/4`$. As with the single $`\nu _\mu `$ spectrum, for $`Q_{\nu _\mu \overline{\nu }_\mu }^{np}`$ multiply eq. (58) by $`7/3`$ and set $`s=1`$. Finally, $`Q_{\nu _\mu \overline{\nu }_\mu }^{tot}`$=$`Q_{\nu _\mu \overline{\nu }_\mu }^{nn}`$\+ $`Q_{\nu _\mu \overline{\nu }_\mu }^{pp}`$+$`Q_{\nu _\mu \overline{\nu }_\mu }^{np}`$.
#### 1 The Non-Degenerate Nucleon Limit
In the non-degenerate nucleon limit, the term $`_1_2(1_3)(1_4)`$ reduces to $`e^{\eta _1}e^{\eta _2}e^{2(u_++u_{})}`$ which is independent of angle. This tremendous simplification allows for easy integration over $`u_+`$ and $`u_{3c}`$ in eqs. (46), (57), and (58). The total volumetric emissivity of a single $`\nu _\mu \overline{\nu }_\mu `$ pair in this limit, ignoring $`\nu _\mu `$ and $`\overline{\nu }_\mu `$ blocking in the final state, is
$$Q_{\nu _\mu \overline{\nu }_\mu }1.04\times 10^{30}\zeta (X\rho _{14})^2\left(\frac{T}{\mathrm{MeV}}\right)^{5.5}\mathrm{ergs}\mathrm{cm}^3\mathrm{s}^1.$$
(60)
For $`nn`$ and $`pp`$ bremsstrahlung, $`X`$ is the number fraction of neutrons ($`X_n`$) or protons ($`X_p`$), respectively. For the mixed-nucleon process ($`np`$), $`X^2`$ becomes $`(28/3)X_nX_p`$. Figure 1 compares the non-degenerate nucleon limit (eq. 60) with the arbitrary nucleon degeneracy generalization (eq. 58) in the case of neutron-neutron ($`nn`$) bremsstrahlung, as a function of the neutron degeneracy $`\eta _n=\mu _n/T`$. The filled square shows the degenerate limit obtained by ref. . Note that at $`\eta _n0`$, the fractional difference between the two is just $`12`$%. At realistic neutron degeneracies within the core ($`\eta _n2`$), this difference approaches 30%.
The single differential $`\nu _\mu `$ emissivity can be written in terms of the pair emissivity :
$`{\displaystyle \frac{dQ}{d\epsilon _\nu }}`$ $`=`$ $`C\left({\displaystyle \frac{Q_{\nu _\mu \overline{\nu }_\mu }}{T^4}}\right)\epsilon _\nu ^3{\displaystyle _1^{\mathrm{}}}{\displaystyle \frac{e^{2q_\nu x}}{x^3}}(x^2x)^{1/2}𝑑x`$ (61)
$`=`$ $`C\left({\displaystyle \frac{Q_{\nu _\mu \overline{\nu }_\mu }}{T^4}}\right)\epsilon _\nu ^3{\displaystyle _{q_\nu }^{\mathrm{}}}{\displaystyle \frac{e^q}{q}}K_1(q)(qq_\nu )^2𝑑q,`$ (62)
where $`C=2310/20481.128`$, $`q_\nu =\epsilon _\nu /2T`$, $`q=ϵ/2T`$, and $`K_1`$ is the standard modified Bessel function of imaginary argument. A useful fit to eq. (62), good to better than 3% over the full range of relevant neutrino energies is
$$\frac{dQ}{d\epsilon _\nu }0.234\frac{Q_{\nu _\mu \overline{\nu }_\mu }}{T}\left(\frac{\epsilon _\nu }{T}\right)^{2.4}e^{1.1\epsilon _\nu /T}.$$
(63)
Using eq. (31), we obtain the contribution to the Boltzmann equation including Pauli blocking of $`\nu _\mu `$ and $`\overline{\nu }_\mu `$ neutrinos in the final state:
$$\frac{_\nu }{t}=𝒞s\zeta _0^{\mathrm{}}𝑑\epsilon _{\overline{\nu }}(\epsilon _{\overline{\nu }}^2/ϵ)K_1\left(\frac{\beta ϵ}{2}\right)e^{\beta ϵ/2}\left\{(1_\nu )(1_{\overline{\nu }})_\nu _{\overline{\nu }}e^{\beta ϵ}\right\}.$$
(64)
where
$$𝒞=\frac{G^2m^{4.5}}{\pi ^{6.5}}\left(\frac{f}{m_\pi }\right)^4g_A^2T^{2.5}e^{\eta _1}e^{\eta _2}\frac{2G^2g_A^2}{\pi ^{3.5}}\left(\frac{f}{m_\pi }\right)^4\frac{m^{1.5}}{T^{.5}}n_1n_2.$$
(65)
In obtaining eq. (65), we have used the thermodynamic identity in the non-degenerate limit,
$$e^{\eta _i}=\left(\frac{2\pi }{mT}\right)^{3/2}\frac{n_i}{2},$$
(66)
where $`n`$ is the number density of nucleons considered and $`i`$ is 1 or 2 for neutrons or protons, depending on which nucleon bremsstrahlung process is considered.
## IV Results
The numerical algorithm we have developed accepts arbitrary initial $`\nu _\mu `$ and $`\overline{\nu }_\mu `$ phase-space distributions. Using the scattering formalism developed in the previous section, we evolve two initial distribution functions: (1) a broad Gaussian in energy centered at 40 MeV with a maximum of $`_\nu =0.80`$ and a full-width at half-maximum of $``$28.6 MeV, and (2) a Fermi-Dirac distribution at a temperature 2$`\times `$ the temperature of the surrounding matter and with zero chemical potential. While the former is unphysical in the context of supernova calculations, it illustrates the effects of blocking on both the average energy transfer and the rates for each scattering process. Furthermore, its evolution is more dynamic than the Fermi-Dirac distribution. As a result, the way in which the distribution is spread and shifted in time is more apparent. The essential differences between the two processes are then more easily gleaned. The latter initial distribution is motivated by consideration of the environment within the core of a supernova. The $`\nu _\mu `$ and $`\overline{\nu }_\mu `$ distribution functions, having been generated as pairs via $`e^+e^{}\nu _\mu \overline{\nu }_\mu `$ and nucleon-nucleon bremsstrahlung should have approximately zero chemical potential. Furthermore, even in the dense core, the $`\nu _\mu `$s will diffuse outward in radius and, hence, from higher to lower temperatures. By starting with a Fermi-Dirac distribution at twice the temperature of the matter at that radius, we learn more about how equilibration might effect the emergent $`\nu _\mu `$ spectrum in an actual collapse or protoneutron star cooling calculation.
For the production and emission processes, we start with zero neutrino occupancy and let each build to an equilibrium distribution of $`\nu _\mu `$s and $`\overline{\nu }_\mu `$s. As a check to the calculation, the asymptotic distribution should be Fermi-Dirac at the temperature of the ambient matter with zero neutrino chemical potential. Throughout these simulations, we take the factor $`\zeta `$ in eq. (57) for nucleon-nucleon bremsstrahlung to be 0.5. (This factor represents our ignorance of the importance of the nucleon momentum transfer terms.)
We repeat these calculations for four temperature, density, and composition points (StarA, StarB, StarC, and StarD) taken from the one-dimensional collapse calculation profile, Star , corresponding to four radii below the shock ($`80`$ km). Roughly, these points have densities $`10^{14}`$, $`10^{13}`$, $`10^{12}`$, and $`10^{11}`$ g cm<sup>-3</sup>. The actual numbers are shown in Table I.
### A Scattering
Figures 2 and 3 show the evolution of a Gaussian distribution at $`t=0`$ to an equilibrium Fermi-Dirac distribution at the temperature of the surrounding matter due to $`\nu _\mu `$-neutron ($`\nu _\mu n`$) and $`\nu _\mu `$-electron ($`\nu _\mu e^{}`$) scattering, respectively. The equilibrium distribution has a non-zero neutrino chemical potential set by the initial total number of $`\nu _\mu `$s, which is conserved to better than .001% throughout the calculation. Multiple curves on each plot show snapshots of $`_\nu `$ in time from $`t=0`$ to 1000 microseconds ($`\mu `$s). Both calculations were carried out at the thermodynamic point StarB whose characteristics are shown in Table I. StarB is indicative of the core of a supernova, a region of moderate to high temperatures ($`T15`$ MeV) and densities of $`10^{13}`$ g cm<sup>-3</sup>. These two figures illustrate the fundamental differences between $`\nu _\mu e^{}`$ and $`\nu _\mu n`$ scattering as thermalization processes. Curve A in Fig. 2 and curve C in Fig. 3 indicate that at high $`\nu _\mu `$ energies ($`\epsilon _\nu `$ 30 MeV) $`\nu _\mu n`$ scattering is a much more effective thermalization mechanism. At $`\epsilon _\nu 40`$ MeV both curves show the distribution is within $``$30% of equilibrium. Importantly, however, curve A is at 0.33 $`\mu `$s for $`\nu _\mu n`$ scattering whereas curve C is at 3.30 $`\mu `$s for $`\nu _\mu e^{}`$ scattering. Curve C, in Fig. 2 for $`\nu _\mu n`$ scattering, also at $`t=3.30`$ $`\mu `$s, shows that above $``$25 MeV the distribution has almost equilibrated. For $`\nu _\mu e^{}`$ scattering, similar evolution at high neutrino energies takes approximately 25 $`\mu `$s. These simple estimates reveal that $`\nu _\mu n`$ scattering is about 10 times faster than $`\nu _\mu e^{}`$ scattering at equilibrating $`\nu _\mu `$s with energies greater than approximately 25 MeV.
This situation is reversed at low $`\epsilon _\nu `$s. Comparing curve E at $`t=33.0`$ $`\mu `$s in both Fig. 2 and Fig. 3, we can see that at $`10`$ MeV both distributions have filled to approximately the same percentage of the asymptotic, equilibrium $`_\nu `$. However, below $`\epsilon _\nu `$8 MeV, $`\nu _\mu n`$ scattering has not filled $`_\nu `$ to the extent $`\nu _\mu e^{}`$ scattering has. In fact, the rate at which these low energy states are filled by $`\nu _\mu n`$ scattering is very low; the energy transfer ($`\omega `$) is much smaller than the incident $`\nu _\mu `$ energy. In this regime, the Fokker-Planck approximation for the time evolution of $`_\nu `$ in energy space may be applicable. In marked contrast, Fig. 3 indicates how effective $`\nu _\mu e^{}`$ scattering is at filling the lowest $`\epsilon _\nu `$ states. Curves F from Figs. 2 and 3, taken at 1000 $`\mu `$s, show that though the distribution has reached equilibrium via $`\nu _\mu e^{}`$ scattering, for $`\nu _\mu n`$ scattering the very lowest energy states remain unfilled. For each of the four points in the Star profile we consider, $`\nu _\mu n`$ scattering dominates at high energies ($`20`$ MeV), whereas $`\nu _\mu e^{}`$ scattering dominates at low $`\nu _\mu `$ energies ($`10`$ MeV) and particularly for $`\epsilon _\nu 3`$ MeV.
Figures 4 and 5 depict the evolution of $`_\nu `$ via $`\nu _\mu n`$ and $`\nu _\mu e^{}`$ scattering, respectively, for an initial Fermi-Dirac distribution at 2$`\times `$ the temperature of the surrounding neutrons and electrons and with zero neutrino chemical potential. This calculation was carried out at StarC (see Table I), which is representative of the outer core, in the semi-transparent regime, where the neutrinos begin to decouple from the matter (near the neutrinosphere). The same systematics highlighted in the discussion of the evolution of the initial Gaussian distribution for StarB are borne out in these figures. Curves A and B on both plots, denoting 0.10 and 0.33 milliseconds (ms) of elapsed time, respectively, confirm that above $`\epsilon _\nu 15`$ MeV $`\nu _\mu n`$ scattering dominates thermalization.
Figures 6 and 7 show $`\omega _{in}`$ and $`\omega _{out}`$, as defined in eqs. (6) and (7), for $`\nu _\mu n`$ scattering and $`\nu _\mu e^{}`$ scattering, respectively. The separate curves portray the evolution in time of the thermal average energy transfers as the distributions evolve to equilibrium (cf. Figs. 4 and 5). As one would expect from kinematic arguments, the magnitudes of both $`\omega _{in}`$ and $`\omega _{out}`$ for $`\nu _\mu n`$ scattering are much less than those for $`\nu _\mu e^{}`$ scattering. Though the energy transfers are much smaller, even at the highest energies, $`\nu _\mu n`$ scattering still dominates $`\nu _\mu e^{}`$ scattering in thermalizing the $`\nu _\mu `$ distribution because the rate for scattering is so much larger. At low neutrino energies, however, both average energy transfers for neutron scattering go to zero, whereas they approach large negative values ($`20`$ MeV) for electron scattering. At these low energies, the fact that the rate for $`\nu _\mu n`$ scattering is larger than for $`\nu _\mu e^{}`$ scattering fails to compensate for the vanishing energy transfer. For example, at $`\epsilon _\nu =3`$ MeV and $`t=33`$ ms, the energy transfer for $`\nu _\mu e^{}`$ scattering is more than 100 times that for $`\nu _\mu n`$ scattering.
In order to fold in information about both the rate of scattering and the average thermal energy transfer, we plot $`\mathrm{\Gamma }_D`$ and $`\mathrm{\Gamma }_E`$ (eqs. 8 and 9) in Fig. 8 for all four points considered in the Star profile. We show here a snapshot of the rates for both scattering processes for a Fermi-Dirac distribution initially at twice the local matter temperature, with zero neutrino chemical potential. Note that the spikes in $`\mathrm{\Gamma }_D`$ indicate the neutrino energy at which $`\omega _{out}=0`$ (cf. Figs. 6 and 7). In general, we find that as $`\epsilon _\nu 0`$, $`\mathrm{\Gamma }_D`$ and $`\mathrm{\Gamma }_E`$ go to zero for $`\nu _\mu `$-neutron scattering, whereas $`\mathrm{\Gamma }_D`$ approaches a constant and $`\mathrm{\Gamma }_E`$ gets very large for $`\nu _\mu e^{}`$ scattering . This is a consequence of the fact that, regardless of $`_\nu `$, $`\omega _{out}0`$ for $`\nu _\mu n`$ scattering as $`\epsilon _\nu 0`$, as shown in Fig. 6. For $`\nu _\mu e^{}`$ scattering the situation is different. As Fig. 7 reveals, $`\omega _{out}`$ approaches $`20`$ MeV at $`\epsilon _\nu =0`$. As expected from our analysis of the evolution of $`_\nu `$ in Figs. 4 and 5, at approximately 40 MeV the thermalization rate for $`\nu _\mu n`$ scattering for StarB is about an order of magnitude greater than that for $`\nu _\mu e^{}`$ scattering. Specifically, the $`\mathrm{\Gamma }_D`$’s cross at $``$15 MeV, whereas the $`\mathrm{\Gamma }_E`$’s cross at $``$20 MeV. Below these energies, both $`\nu _\mu n`$ rates drop off precipitously as a consequence of the fact that $`\omega _{out}0`$. Below $`\epsilon _\nu 5`$ MeV, the thermalization rate for $`\nu _\mu e^{}`$ scattering dominates by 2-5 orders of magnitude. As evidenced by the other panels in Fig. 8, this same trend holds in the other regions of the stellar profile. In general, the rates drop over the whole energy range for both processes as the density and temperature decrease, but the same systematics hold. In fact, for StarA, StarC, and StarD the $`\mathrm{\Gamma }_E`$ and $`\mathrm{\Gamma }_D`$ crossing points for both processes are lower than those for StarB. As a result of the higher temperature at this radius ($`T14.5`$ MeV) $`\nu _\mu e^{}`$ scattering is important in thermalizing slightly higher energy neutrinos than at the other radii. For StarC and StarD, specifically, both rates cross at neutrino energies less than 12 MeV.
These results demonstrate that $`\nu _\mu `$-nucleon scattering is an important thermalization process from the dense core through the semi-transparent regime for $`\nu _\mu `$s with energies greater than approximately 15 MeV. The addition of this energy transfer mechanism implies that the $`\nu _\mu `$s stay energetically coupled to the surrounding matter longer than has been previously estimated . We can approximate the radius at which the $`\nu _\mu `$s energetically decouple from the matter (the $`E_\mu `$-sphere) by observing when the diffusion timescale is approximately equal to the equilibration timescale given by $`\mathrm{\Gamma }_D^1=\tau _D`$, as defined in eq. (8). Using this crude approximation we find that by including $`\nu _\mu `$-nucleon energy transfer the $`E_\mu `$-sphere is pushed outward in radius by approximately 3 kilometers. This difference in radius corresponds to a 1-2 MeV drop in the matter temperature in the model Star. The average energy of the emergent spectrum is roughly correlated with the local matter temperature of the $`E_\mu `$-sphere. Therefore, we conclude that $`\nu _\mu `$-nucleon energy transfer in full transport calculations will likely soften the emergent $`\nu _\mu `$ spectrum.
### B Emission and Absorption
Figure 9 shows the total integrated volumetric emissivity as a function of radius in the model Star for nucleon-nucleon bremsstrahlung in the non-degenerate nucleon limit (eq. 60), its generalization for arbitrary nucleon degeneracy (eq. 58), and the emissivity for $`e^+e^{}`$ annihilation (eq. 33). Note that not one of these expressions contains neutrino blocking terms and that the general bremsstrahlung rate crosses that for $`e^+e^{}`$ annihilation at $`23`$ kilometers where $`\rho 6\times 10^{12}`$ g cm<sup>-3</sup>, $`T11`$ MeV, and $`Y_e0.13`$. While this plot gives a general idea of where $`e^+e^{}`$ annihilation should begin to compete with nucleon-nucleon bremsstrahlung, it fails to include the differential nature of the production in energy. In addition, it does not include absorption or blocking effects, which quantitatively alter the relative strength of the emission.
To begin to understand the import of these terms and the character of each pair production process, we include Figs. 10 and 11, which show the time evolution of $`_\nu `$ via nucleon-nucleon bremsstrahlung and electron-positron annihilation, respectively, for the point StarC, initialized with zero $`\nu _\mu `$ and $`\overline{\nu }_\mu `$ phase-space occupancies. The final equilibrium distribution is Fermi-Dirac at the temperature of the surrounding matter, with zero neutrino chemical potential. Comparing curve C on both graphs, which marks 10.0 milliseconds (ms) of elapsed time, one can see that bremsstrahlung dominates production below $`15`$ MeV. Indeed, bremsstrahlung overshoots its equilibrium distribution at energies below 10 MeV before finally filling the higher $`\epsilon _\nu `$ states. In contrast, electron-positron annihilation fills the higher states first and moves slowly toward the low-lying neutrino energies, taking a factor of 10 more time at this thermodynamic point to reach equilibrium.
In Figs. 12 and 13, we plot $`\mathrm{\Gamma }_{in}`$ and $`\mathrm{\Gamma }_{out}`$, as defined in eqs. (4) and (5), for both production processes at the point StarB. As one would predict from our simple observations of the time evolution of $`_\nu `$, the bremsstrahlung rates are much faster ($``$2 orders of magnitude) than the $`e^+e^{}`$ annihilation rates at low neutrino energies. At StarC, $`e^+e^{}`$ annihilation competes with bremsstrahlung above $`\epsilon _\nu 15`$ MeV. For StarB, however, at a matter density an order of magnitude greater than that for StarC, the energy at which nucleon-nucleon bremsstrahlung becomes more important than $`e^+e^{}`$ annihilation is $`60`$ MeV. In this regime, where $`T1214`$ MeV and $`\rho 10^{13}`$ g cm<sup>-3</sup>, we find that bremsstrahlung dominates neutrino pair-production via electron-positron annihilation. A close look at the evolution of the total thermal average neutrino energy ($`\epsilon _\nu `$) reveals that $`_\nu `$ reaches its asymptotic equilibrium distribution via nucleon-nucleon bremsstrahlung in $`1`$ ms. Electron-positron annihilation takes $`50`$ ms to fill all but the very lowest energy states. This trend continues as the matter becomes more dense. For StarA, well beneath the neutrinospheres at $`\rho 10^{14}`$ g cm<sup>-3</sup>, the rates for bremsstrahlung and electron-positron annihilation never cross. In fact, the former produces an equilibrium Fermi sea of $`\nu _\mu `$’s in $`50`$ $`\mu `$s, whereas the latter takes $`10^3`$ seconds. This difference of 8 orders of magnitude in timescale, however, is a bit misleading. Similar to $`\nu _\mu n`$ scattering, $`e^+e^{}`$ annihilation has trouble filling only the very lowest neutrino energy states. In actuality, at the highest energies, both $`\mathrm{\Gamma }_{in}`$ and $`\mathrm{\Gamma }_{out}`$ for $`e^+e^{}`$ annihilation come within 3-4 orders of magnitude of the rates for bremsstrahlung at the same energy. Still, the difference is striking. As the temperature drops from StarB (14 MeV) to StarA (10 MeV) and the density increases by an order of magnitude, $`\eta _e`$ goes from 3.79 to 15.75. Consequently, Pauli blocking of electrons in the final state suppresses the process $`\nu _\mu \overline{\nu }_\mu e^+e^{}`$, and the phase-space density of positrons is depleted to such an extent that $`e^+e^{}\nu _\mu \overline{\nu }_\mu `$ is suppressed as well. We conclude that beneath the neutrinospheres and specifically for $`\rho 10^{13}`$ g cm<sup>-3</sup>, nucleon-nucleon bremsstrahlung is the primary and dominant $`\nu _\mu \overline{\nu }_\mu `$ source. Near the neutrinosphere, within the gain region and behind the shock, between 30 km and 60 km at $`\rho 10^{12}`$ g cm<sup>-3</sup> and $`T68`$ MeV, bremsstrahlung competes with $`e^+e^{}`$ annihilation at all neutrino energies and is the primary production process for the low-lying $`\epsilon _\nu `$ and $`\epsilon _{\overline{\nu }}`$ states.
The addition of nucleon-nucleon bremsstrahlung will have quantitative implications for the $`\nu _\mu `$ and $`\nu _\tau `$ emergent spectra. Specifically, they should be softer and brighter. Burrows et al. confirm this with their study of static supernova and protoneutron star atmospheres, having included nucleon-nucleon bremsstrahlung in the non-degenerate limit. In addition to observing a systematic softening, they also find that the $`\nu _\mu `$ spectrum is a factor of 2 more luminous at $`\epsilon _\nu =10`$ MeV.
## V Summary and Conclusions
Our results for equilibration via $`\nu _\mu `$-electron scattering and $`\nu _\mu `$-nucleon scattering indicate that the latter competes with or dominates the former as a thermalizer for neutrino energies $`10`$ MeV for $`\rho 1\times 10^{11}`$ g cm<sup>-3</sup> at all temperatures. At neutrino energies $`30`$ MeV the difference at all densities and temperatures is approximately an order of magnitude. For the production and absorption processes, we find that nucleon-nucleon bremsstrahlung, at the average energy of an equilibrium Fermi-Dirac distribution at the local temperature, is 5 and 2 orders of magnitude faster than $`e^+e^{}`$ annihilation at StarA ($`T10`$ MeV, $`\rho 10^{14}`$ g cm<sup>-3</sup>) and StarB ($`T15`$ MeV, $`\rho 10^{13}`$ g cm<sup>-3</sup>), respectively. Only for $`\rho 10^{12}`$ g cm<sup>-3</sup> and $`T6`$ MeV does $`e^+e^{}\nu _\mu \overline{\nu }_\mu `$ begin to compete with bremsstrahlung at all energies. We conclude from this study that the emergent $`\nu _\mu `$ and $`\nu _\tau `$ spectrum is (1) brighter and (2) softer than previously estimated. The former results from the inclusion of the new pair emission process, nucleon-nucleon bremsstrahlung. The latter is a consequence of both the increased energy coupling between the nuclear and neutrino fluids through $`\nu _\mu `$-nucleon scattering and the fact that bremsstrahlung dominates $`e^+e^{}`$ annihilation near the neutrinospheres at the lowest neutrino energies. While the full transport problem, including $`\nu _\mu `$-nucleon scattering energy redistribution and nucleon-nucleon bremsstrahlung, must be solved in order to delineate precisely what consequences these processes have for the emergent $`\nu _\mu `$ spectrum, these calculations demonstrate that they should not be omitted.
## VI Acknowledgments
The authors thank Sanjay Reddy for helpful correspondence. A.B. and T.A.T. acknowledge support under NSF Grant No. AST96-14794 and J.E.H. acknowledges funding from the Fundacão de Amparo a Pesquisa do Estado de São Paulo.
## A Neutrino-electron scattering
Each of the retarded polarization functions in eqs. (25-27) can be written in terms of one-dimensional integrals over electron energy ($`\epsilon _e`$), which we label $`I_n`$ ;
$$\mathrm{Im}\mathrm{\Pi }_L^R(q,\omega )=\frac{q_\mu ^2}{2\pi |q|^3}\left[I_2+\omega I_1+\frac{q_\mu ^2}{4}I_0\right],$$
(A1)
$$\mathrm{Im}\mathrm{\Pi }_T^R(q,\omega )=\frac{q_\mu ^2}{4\pi |q|^3}\left[I_2+\omega I_1+\left(\frac{q_\mu ^2}{4}+\frac{q^2}{2}+m^2\frac{q^2}{q_\mu ^2}\right)I_0\right],$$
(A2)
$$\mathrm{Im}\mathrm{\Pi }_A^R(q,\omega )=\frac{m^2}{2\pi |q|}I_0,$$
(A3)
and
$$\mathrm{Im}\mathrm{\Pi }_{VA}^R(q,\omega )=\frac{q_\mu ^2}{8\pi |q|^3}\left[\omega I_0+2I_1\right].$$
(A4)
The authors of were able to express the $`I_n`$’s in terms of polylogarithmic integrals such that
$$I_0=Tz\left(1\frac{\xi _1}{z}\right),$$
(A5)
$$I_1=T^2z\left(\eta _e\frac{z}{2}\frac{\xi _2}{z}\frac{e_{}\xi _1}{zT}\right),$$
(A6)
and
$$I_2=T^3z\left(\eta _e^2z\eta _e+\frac{\pi ^2}{3}+\frac{z^2}{3}+2\frac{\xi _3}{z}2\frac{e_{}\xi _2}{Tz}+\frac{e_{}^2\xi _1}{T^2z}\right),$$
(A7)
where $`\eta _e=\mu _e/T`$ is the electron degeneracy, $`z=\beta \omega `$, $`\omega `$ is the energy transfer, and
$$e_{}=\frac{\omega }{2}+\frac{q}{2}\sqrt{14\frac{m^2}{q_\mu ^2}}.$$
(A8)
In eqs. (A5-A7), the $`\xi _n`$’s are differences between polylogarithmic integrals; $`\xi _n=\mathrm{Li}_n(\alpha _1)\mathrm{Li}(\alpha _2)`$, where
$$\mathrm{Li}_n(y)=_0^y\frac{\mathrm{Li}_{n1}(x)}{x}𝑑x,$$
(A9)
and Li$`{}_{1}{}^{}(x)=\mathrm{ln}(1x)`$. The arguments necessary for computing the integrals are $`\alpha _1=\mathrm{exp}[\beta (e_{}+\omega )\eta _e]`$ and $`\alpha _2=\mathrm{exp}(\beta e_{}\eta _e)`$.
## B electron-positron annihilation
The production kernel is defined by
$$R^p(\epsilon _\nu ,\epsilon _{\overline{\nu }},\mathrm{cos}\theta )=\frac{1}{2\epsilon _\nu \epsilon _{\overline{\nu }}}\frac{d^3\stackrel{}{p}}{(2\pi )^32\epsilon }\frac{d^3\stackrel{}{p}^{}}{(2\pi )^32\epsilon ^{}}_e^{}_{e^+}\left(\frac{1}{4}\underset{s}{}||^2\right)(2\pi )^4\delta ^4(\mathrm{P}).$$
(B1)
The differential production spectrum for final state $`\nu _\mu `$s can then be written as
$$\frac{dQ}{d\epsilon _\nu }=(1_\nu )\frac{\epsilon _\nu ^3}{(2\pi )^6}𝑑\mathrm{\Omega }_0^{\mathrm{}}\epsilon _{\overline{\nu }}^2𝑑\epsilon _{\overline{\nu }}_1^1𝑑\mu ^{}_0^{2\pi }𝑑\varphi R^p(\epsilon _\nu ,\epsilon _{\overline{\nu }},\mathrm{cos}\theta )(1_{\overline{\nu }}),$$
(B2)
where $`d\mathrm{\Omega }`$ is the differential solid angle for the $`\nu _\mu `$ neutrino, $`\mu ^{}=\mathrm{cos}\theta ^{}`$ is the cosine of the $`\overline{\nu }_\mu `$ angular coordinate, and $`\varphi `$ is the azimuthal angle between $`\nu _\mu `$ and $`\overline{\nu }_\mu `$. Expanding the production kernel in a Legendre series in the scattering angle, $`\mathrm{cos}\theta =\mu \mu ^{}+[(1\mu ^2)(1\mu ^{\mathrm{\hspace{0.17em}2}})]^{1/2}\mathrm{cos}\varphi `$,
$$R^p(\epsilon _\nu ,\epsilon _{\overline{\nu }},\mathrm{cos}\theta )=\frac{1}{2}\underset{l}{}(2l+1)\mathrm{\Phi }_l^p(\epsilon _\nu ,\epsilon _{\overline{\nu }})P_l(\mathrm{cos}\theta )\frac{1}{2}\mathrm{\Phi }_0^p(\epsilon _\nu ,\epsilon _{\overline{\nu }})+\frac{3}{2}\mathrm{\Phi }_1^p(\epsilon _\nu ,\epsilon _{\overline{\nu }})\mathrm{cos}\theta .$$
(B3)
$`\mathrm{\Phi }_0^p`$, in eqs. (30) and (B3), is given by
$$\mathrm{\Phi }_0^p(\epsilon _\nu ,\epsilon _{\overline{\nu }})=\frac{G^2}{\pi }_0^{\epsilon _\nu +\epsilon _{\overline{\nu }}}𝑑\epsilon _e^{}_{e^+}H_0(\epsilon _\nu ,\epsilon _{\overline{\nu },}\epsilon ),$$
(B4)
where $`_{e^+}`$ is a function of $`\epsilon ^{}(=\epsilon _\nu +\epsilon _{\overline{\nu }}\epsilon )`$ and
$$H_0(\epsilon _\nu ,\epsilon _{\overline{\nu }},\epsilon )=(V+A)^2J_0^I(\epsilon _\nu ,\epsilon _{\overline{\nu }},\epsilon )+(VA)^2J_0^{II}(\epsilon _\nu ,\epsilon _{\overline{\nu }},\epsilon ).$$
(B5)
Each $`J_0`$ in eq. (B5) is a polynomial in $`\epsilon _\nu `$, $`\epsilon _{\overline{\nu }}`$, and $`\epsilon `$ of dimension \[energy\]. They are related to each other by
$$J_0^I(\epsilon _\nu ,\epsilon _{\overline{\nu }},\epsilon )=J_0^{II}(\epsilon _{\overline{\nu }},\epsilon _\nu ,\epsilon ).$$
(B6)
Both $`J_0^I`$ and $`J_0^{II}`$ can be found in ref. (correcting for the typo in their eq. C67). From eqs. (30) and (B5) we see that the differences between the spectra for $`\nu _\mu `$s and $`\overline{\nu }_\mu `$s for a given temperature and electron degeneracy ($`\eta _e`$) arise solely from the relative weighting constants $`(V+A)^2`$ and $`(VA)^2`$ in eq. (B5) for $`J_0^I`$ and $`J_0^{II}`$, respectively. Indeed, in this approximation the same can be said for the difference in the spectrum between $`\nu _e`$ and $`\nu _\mu `$ neutrinos.
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# 𝒩=4 Super Yang-Mills Low-Energy Effective Action at Three and Four Loops
## Abstract
We investigate the low-energy effective action in N=4 super Yang-Mills theory with gauge group $`SU(n)`$ spontaneously broken down to its maximal torus. Using harmonic superspace technique we prove an absence of any three- and four-loop corrections to non-holomorphic effective potential depending on $`N=2`$ superfield strengths. A mechanism responsible for vanishing arbitrary loop corrections to low-energy effective action is discussed.
Supersymmetry imposes the significant restrictions on a structure of effective action in field models. It is naturally to expect that the most strong restrictions have to arise in maximally extended rigid supersymmetric model, that is in $`N=4`$ super-Yang-Mills theory.
Recently Dine and Seiberg found that a dependence of $`N=4`$ supersymmetric Yang-Mills low-energy effective action on $`N=2`$ superfield strengths $`𝒲`$ and $`\overline{𝒲}`$ is exactly fixed only by general properties of the quantum theory under consideration like finiteness and scale independence . According to ref the leading low-energy contributions to effective action in $`N=4`$ SYM with gauge group $`SU(2)`$ spontaneously broken down to $`U(1)`$ are described by non-holomorphic effective potential $`(𝒲,\overline{𝒲})`$ of the form
$`(𝒲,\overline{𝒲})=c\mathrm{log}\left({\displaystyle \frac{𝒲^2}{\mathrm{\Lambda }^2}}\right)\mathrm{log}\left({\displaystyle \frac{\overline{𝒲}^2}{\mathrm{\Lambda }^2}}\right)`$ (1)
Here $`\mathrm{\Lambda }`$ is some scale. The effective potential (1) possesses by two remarkable properties. First, the corresponding effective action
$`{\displaystyle d^4xd^8\theta (𝒲,\overline{𝒲})}`$ (2)
is scale independent. Second, any quantum corrections, if they exist at all, are included into a single constant $`c`$.
The explicit calculations of the non-holomorphic effective potential and finding the coefficient $`c`$ in one-loop approximation have been carried out in refs \[2-4\]. Extension of the above results for $`N=4`$ Yang-Mills theory with gauge group $`SU(n)`$, $`n>2`$, spontaneously broken down to its maximal torus have been developed in refs \[5-8\]. General structure of low-energy effective action in $`N=2,4`$ superconformal invariant field models was investigated in ref .
In the paper Dine and Seiberg presented the qualitative arguments based on principle of naturalness (see application of this principle to SUSY theories in ref ) that effective potential (1) gets neither perturbative nor non-perturbative quantum corrections beyond one loop. Therefore the expression (1) together with the results of one-loop calculations of the coefficient $`c`$ \[2-4\] determines exact low-energy effective action in $`N=4`$ $`SU(2)`$ Yang-Mills theory. The above arguments have also been discussed in refs . Another approach leading to the same conclusion about structure of low-energy effective action was developed in recent papers .
We would like to pay an attention that a mechanism providing an absence of higher loop corrections to non-holomorphic effective potential in $`N=4`$ Yang-Mills theory is unknown up to now. The firm results concern only two-loop approximation where the corresponding corrections are prohibited by $`N=2`$ supersymmetry (see also direct two-loop calculations in refs ).
An interesting aspect of $`N=4`$ Yang-Mills theory with gauge group $`SU(n)`$, $`n>2`$, has been recently pointed out in refs . The symmetry arguments do not prohibit an appearance of some new invariant structures, besides logarithmic, in non-holomorphic effective potential which are absent at $`n=2`$. The direct calculations \[5-8\] do not confirm such structures in one-loop approximation. However a question concerning their appearance at higher loops is open.
In this paper we are going to develop a technique for investigating a structure of the non-holomorphic effective potential at higher loops, to find a mechanism providing a cancellation of higher-loop contributions, and to clarify situation concerning the non-logarithimic corrections to low-energy effective action in $`N=4`$ Yang-Mills theory with gauge group $`SU(n)`$, $`n>2`$, spontaneously broken to its maximal torus. To be more precise, we investigate a structure of three- and four-loop supergraphs and show how $`N=4`$ supersymmetry provides an efficient mechanism of supergraph cancellations.
We consider $`N=4`$ Yang-Mills theory formulated in terms of $`N=2`$ superfields and get $`N=2`$ Yang-Mills theory coupled to hypermultiplet in adjoint representation. The most convenient and simple way to carry out quantum calculations in $`N=2`$ SUSY models is given by harmonic superspace approach \[16-18\] which is used in the paper. The various implementations of this approach to effective action in $`N=2`$ SUSY theories are discussed in refs .
The starting point of our consideration is the classical action of the $`N=4`$ super-Yang-Mills theory in $`q`$-hypermultiplet realization written in harmonic superspace
$`S`$ $`=`$ $`{\displaystyle \frac{1}{g^2}}\mathrm{tr}{\displaystyle d^{12}z\underset{n=2}{\overset{\mathrm{}}{}}\frac{(\mathrm{i})^n}{n}𝑑u_1\mathrm{}𝑑u_n\frac{V^{++}(z,u_1)\mathrm{}V^{++}(z,u_n)}{(u_1^+u_2^+)(u_2^+u_3^+)\mathrm{}(u_n^+u_1^+)}}+`$ (3)
$`+`$ $`{\displaystyle 𝑑\zeta ^{(4)}𝑑u\stackrel{˘}{q}^+(D^{++}+iV^{++})q^+}.`$
The denotions introduced in the paper are employed here and further.
The calculations are carried out in framework of $`N=2`$ background field method . We make background-quantum splitting by the rule
$$V^{++}V^{++}+gv^{++}$$
(4)
and construct the corresponding Faddeev-Popov ghost action in the form
$$S_{gh}=\mathrm{tr}𝑑\zeta ^{(4)}𝑑u𝐛(^{++})^2𝐜\mathrm{i}g\mathrm{tr}𝑑u𝑑\zeta ^{(4)}^{++}𝐛[v^{++},𝐜].$$
(5)
The background-dependent superpropagators in the theory with action of $`N=2`$ gauge multiplet and $`N=2`$ matter hypermultiplet (3) and action of ghosts (5) have been obtained in and look like
$`<v_\tau ^{++}(1)v_\tau ^{++}(2)>`$ $`=`$ $`{\displaystyle \frac{\mathrm{i}}{\stackrel{}{\mathrm{}}}}\stackrel{}{(𝒟_1^+)^4}\left\{\delta ^{12}(z_1z_2)\delta ^{(2,2)}(u_1,u_2)\right\}`$
$`<q_\tau ^+(1)\stackrel{˘}{q}_\tau ^+(2)>`$ $`=`$ $`{\displaystyle \frac{\mathrm{i}}{\stackrel{}{\mathrm{}}}}\stackrel{}{(𝒟_1^+)^4}\left\{\delta ^{12}(z_1z_2){\displaystyle \frac{1}{(u_1^+u_2^+)^3}}\right\}\stackrel{}{(𝒟_2^+)^4}`$
$`<\omega _\tau (1)\omega _\tau ^\mathrm{T}(2)>`$ $`=`$ $`{\displaystyle \frac{\mathrm{i}}{\stackrel{}{\mathrm{}}}}\stackrel{}{(𝒟_1^+)^4}\left\{\delta ^{12}(z_1z_2){\displaystyle \frac{(u_1^{}u_2^{})}{(u_1^+u_2^+)^3}}\right\}\stackrel{}{(𝒟_2^+)^4}`$
$`<𝐜_\tau (1)𝐛_\tau (2)>`$ $`=`$ $`{\displaystyle \frac{\mathrm{i}}{\stackrel{}{\mathrm{}}}}\stackrel{}{(𝒟_1^+)^4}\left\{\delta ^{12}(z_1z_2){\displaystyle \frac{(u_1^{}u_2^{})}{(u_1^+u_2^+)^3}}\right\}\stackrel{}{(𝒟_2^+)^4}.`$ (6)
with the operator $`\stackrel{}{\mathrm{}}`$ of the form
$`\stackrel{}{\mathrm{}}`$ $`=`$ $`𝒟^m𝒟_m+{\displaystyle \frac{\mathrm{i}}{2}}(𝒟^{+\alpha }𝒲)𝒟_\alpha ^{}+{\displaystyle \frac{\mathrm{i}}{2}}(\overline{𝒟}_{\dot{\alpha }}^+\overline{𝒲})\overline{𝒟}^{\dot{\alpha }}{\displaystyle \frac{\mathrm{i}}{4}}(𝒟^{+\alpha }𝒟_\alpha ^+𝒲)D^{}`$ (7)
$`+{\displaystyle \frac{\mathrm{i}}{8}}[𝒟^{+\alpha },𝒟_\alpha ^{}]𝒲+{\displaystyle \frac{1}{2}}\{\overline{𝒲},𝒲\}`$
Here index $`\tau `$ denotes that corresponding superfields taken in $`\tau `$-frame where covariant derivatives $`𝒟_\alpha ^i`$, $`\overline{𝒟}_{\dot{\alpha }}^i`$ are introduced to be independent of harmonic coordinates. The propagators look like (6) just in this frame (see details in ).
Our aim consists in calculations of three- and four-loop contributions to non-holomorphic effective potential. We consider the case of the gauge group $`SU(n)`$ spontaneously broken down to its maximal Abelian subgroup. The corresponding background strength $`𝒲`$ is a diagonal matrix of the form
$`𝒲=diag(𝒲_1,𝒲_2\mathrm{}𝒲_n);{\displaystyle \underset{i=1}{\overset{n}{}}}𝒲_i=0.`$ (8)
Since non-holomorphic effective potential depends only on background superfield strengths but not on their derivatives we omit everywhere all terms including derivatives of $`𝒲`$, $`\overline{𝒲}`$ in (6,7). Hence the operator $`\stackrel{}{\mathrm{}}`$ in propagators (6) looks like
$`\stackrel{}{\mathrm{}}`$ $`=`$ $`𝒟^m𝒟_m+𝒲\overline{𝒲},`$ (9)
or, in the manifest form,
$`\stackrel{}{\mathrm{}}`$ $`=`$ $`\left(\begin{array}{ccc}𝒟^m𝒟_m+𝒲_1\overline{𝒲}_1& 0& \mathrm{}\\ 0& \mathrm{}& \mathrm{}\\ 0& \mathrm{}& 𝒟^m𝒟_m+𝒲_n\overline{𝒲}_n\end{array}\right).`$ (13)
As a result we face a problem of calculating three- and four-loop supergraphs in the theory with constant background superfield strengths and operator $`\stackrel{}{\mathrm{}}`$ given by (8) and (9) respectively.
First of all, let us note that arbitrary $`L`$-loop supergraph provides non-zero contribution to non-holomorphic effective potential if and only if number of $`D`$-factors contained in it is equal to $`8L`$ or greater since contracting of any loop to a point by the rule
$`\delta ^8(\theta _1\theta _2)(D^+(u_1))^4(D^+(u_2))^4\delta ^8(\theta _1\theta _2)=(u_1^+u_2^+)^4\delta ^8(\theta _1\theta _2)`$ requires 8 $`D`$-factors. Then, an arbitrary supergraph with $`P_v`$ propagators of $`N=2`$ gauge superfield, $`P_m`$ propagators of matter hypermultiplet, $`P_c`$ propagators of ghosts, $`V_m`$ vertices containing interaction with matter and $`V_c`$ ones including interaction with ghosts contains the following number of $`D`$-factors:
$`N_D=4P_v+8P_m+8P_c4V_m4V_c`$ (14)
because of structure of propagators given in (6) with the operator $`\stackrel{}{\mathrm{}}`$ has the form (9) and vertices corresponding to actions (3,5) (recall that transformation of vertex of the form $`𝑑\zeta ^{(4)}`$ to an integral over whole superspace by the rule $`𝑑\zeta ^{(4)}(D^+)^4=d^{12}z`$ requires four $`D`$-factors). Then, it is easy to see that number of propagators of ghosts and matter is equal to number of vertices including interaction with ghosts and matter respectively since each pure ghost (matter) loop contains equal number of vertices and propagators, i.e. $`P_m=V_m`$, $`P_c=V_c`$. As a result, number of $`D`$-factors in arbitrary supergraph is equal to $`N_D=4P`$ where $`P`$ is a full number of propagators in corresponding supergraph. Hence any $`L`$-loop supergraph can contribute to non-holomorphic effective potential if and only if $`P2L`$. For example, two-loop supergraph should contain 4 and more propagators, and since number of propagators in two-loop supergraphs is no more than three we see that there is no two-loop contribution to $`(𝒲,\overline{𝒲})`$ in accordance with conclusion of where namely such an analysis was used to prove absence of two-loop non-holomorphic contribution. However a situation beyond two loops is much more complicated. In three- and four-loop supergraphs number of propagators should be no less than 6 or 8 respectively. Such supergraphs actually exist and are studied bellow.
Let us consider three-loop supergraphs corresponding to theory defined by actions (3,5). They are given by Figs. 1a, 1b. Here as usual wavy line is used for gauge propagator, solid line is for matter propagator, dashed line is for ghost propagator. The numbers near any supergraph at Fig. 1a will be explained later.
$`\frac{1}{2}`$ $`\frac{1}{2}`$ $`\frac{1}{2}`$
$`1`$ $`1`$ $`Fig.1a`$ $`1`$
$`Fig.1b`$
We are going to show that total contribution of the supergraphs given by Fig. 1a to non-holomorphic effective potential vanishes due to $`N=4`$ supersymmetry. To clarify a mechanism providing manifestation of $`N=4`$ supersymmetry in the supergraphs formulated in terms of $`N=2`$ superfields we introduce a notion of $`N=4`$ superpartner supergraphs. Three supergraphs are called $`N=4`$ superpartners if they have the following structure. One of the supergraphs contains the gauge loop given by Fig. 2a and some system of the propagators associated with this loop. Another supergraph contains the matter loop given by Fig. 2c instead of gauge one and the same system of the propagators as in first case. Third supergraph contains the ghost loop given by Fig. 2b instead of gauge one and the same system of the propagators as in first case. The examples of systems of superpartner supergraphs are given also by Figs. 3a – 3c, Figs. 4a – 4c, Fig. 5. Appearance of such a set of supergraphs turned out to be typical for higher loop contributions. The simplest set of $`N=4`$ superpartner supergraphs arising at one-loop order is given by Fig. 2a – Fig. 2c. We show that sum of three these supergraphs for background dependent superpropagators is equal to zero in the case of constant background superfield strengths.
$`Fig.2a`$ $`Fig.2b`$ $`Fig.2c`$
Structure of the supergraphs containing $`N=4`$ superpartners in the case when propagators do not depend on background superfields was studied in details by GIKOS . However we prove that the same result takes place when we consider non-holomorphic effective potential using the background field dependent propagators (7,10).
To evaluate contributions from supergraphs given by Figs. 2a – 2c one reminds that in $`\tau `$-frame (where the propagators look like (6)) covariant harmonic derivatives $`^{++},^{}`$ coincide with standard harmonic derivatives . Then, the commutation relations $`[^{++},𝒟_A^+]=0`$ with $`𝒟_A^+`$ be either vector or spinor covariant derivative take place in any frame . At the same time we pay attention to the fact that $`[𝒟_\gamma ^+,𝒟_{\alpha \dot{\alpha }}]=iϵ_{\gamma \alpha }\overline{𝒟}_{\dot{\alpha }}^+\overline{𝒲}`$. Therefore one can put $`[𝒟_\gamma ^+,\stackrel{}{\mathrm{}}]=0`$ in the sector of non-holomorphic effective potential.
Let us consider the supergraphs given by Figs.2a – 2c in more details. The contribution from the supergraph given at Fig.2a is equal to (see (6))
$`I_a`$ $`=`$ $`{\displaystyle }d^8\theta _1d^8\theta _2{\displaystyle }du_1du_2du_3dw_1dw_2dw_3{\displaystyle \frac{1}{\stackrel{}{\mathrm{}}}}\delta _{12}^8(𝒟_2^+)^4(𝒟_3^+)^4\delta _{12}^8V^{++}(1)V^{++}(2)\times `$ (15)
$`\times `$ $`{\displaystyle \frac{1}{\stackrel{}{\mathrm{}}}}{\displaystyle \frac{1}{(u_1^+u_2^+)(u_2^+u_3^+)(u_1^+u_3^+)(w_1^+w_2^+)(w_2^+w_3^+)(w_1^+w_3^+)}}\times `$
$`\times `$ $`\delta ^{(2,2)}(u_2,w_2)\delta ^{(2,2)}(u_3,w_3)`$
We take into account that $`𝒟^+`$ commutes with $`\stackrel{}{\mathrm{}}`$ and use the relation
$`\delta _{12}^8(𝒟_2^+)^4(𝒟_3^+)^4\delta _{12}^8=(u_2^+u_3^+)^4\delta _{12}^8`$. Integration over $`w_2`$ and $`w_3`$ leads to
$`I_a`$ $`=`$ $`{\displaystyle d^8\theta 𝑑u_1𝑑u_2𝑑u_3𝑑w_1\frac{1}{\stackrel{2}{\stackrel{}{\mathrm{}}}}V^{++}(1)V^{++}(2)\frac{(u_2^+u_3^+)^2}{(u_1^+u_2^+)(u_1^+u_3^+)(w_1^+u_2^+)(w_1^+u_3^+)}}`$ (16)
This expression coincides with the result obtained in , the only difference consists in presence of $`\stackrel{}{\mathrm{}}`$ instead of $`\mathrm{}`$. However since superfield strength $`𝒲`$ does not depend on harmonic coordinates we can carry out the trick which was used in . We express $`(u_2^+u_3^+)^2`$ as $`D_2^{++}D_3^{++}[(u_2^{}u_3^{})(u_2^+u_3^+)]`$ and integrate by parts to transfer $`D_2^{++}D_3^{++}`$ to other terms of $`I_a`$. After that the $`I_a`$ takes the form
$`I_a`$ $`=`$ $`{\displaystyle d^8\theta 𝑑u_1𝑑w_1\frac{1}{\stackrel{2}{\stackrel{}{\mathrm{}}}}V^{++}(1)V^{++}(2)\frac{u_1^+w_1^{}}{u_1^+w_1^+}}`$ (17)
Analogous consideration allows to show that contributions from supergraphs given by Fig.2b, Fig.2c are respectively equal to
$`I_b`$ $`=`$ $`2{\displaystyle d^8\theta 𝑑u_1𝑑w_1\frac{1}{\stackrel{2}{\stackrel{}{\mathrm{}}}}V^{++}(1)V^{++}(2)\frac{(u_1^+w_1^{})(u_1^{}w_1^+)}{(u_1^+w_1^+)^2}}`$ (18)
and
$`I_c`$ $`=`$ $`2{\displaystyle d^8\theta 𝑑u_1𝑑w_1\frac{1}{\stackrel{2}{\stackrel{}{\mathrm{}}}}V^{++}(1)V^{++}(2)\frac{1}{(u_1^+w_1^+)^2}}`$ (19)
Here $`V^{++}(1)`$, $`V^{++}(2)`$ are external gauge lines (which are contracted to some systems of propagators or vertices in multiloop supergraphs containing graphs given at Fig.2a – Fig.2c as subdiagrams). These contributions are analogous to corresponding expressions given in with the only difference in presence of $`\stackrel{}{\mathrm{}}`$ instead of $`\mathrm{}`$. It is evident that $`I_1+I_2+I_3=0`$ in accordance with , i.e. these supergraphs cancel each other in sector of constant background fields. This effect provides vanishing of sum of supergraphs given by Fig. 3a – Fig. 3c, Fig. 4a – Fig. 4c and Fig. 5.
We note that vanishing of this sum is caused by $`N=4`$ supersymmetry and does not depend on structure of (Abelian) gauge group. Therefore such a situation is common for any $`SU(n)`$ gauge group broken down to its maximal Abelian subgroup independently of value of $`n`$.
Let us return back to the supergraphs given by Fig. 1a. Each supergraph has some combinatoric factor. The straightforward calculations show that they are proportional to each other with the coefficients 1/2 or 1 which are written near the corresponding supergraphs. We call these coefficients the relative factors.
The supergraphs given by Fig. 1a can be equivalently regroupped into the sets of $`N=4`$ superpartners as shown on Figs. 3a - 3c where the relative factors are also written near the supergraphs. All supergraphs on Fig. 3a are $`N=4`$ superpartners as well as the supergraphs on Figs. 3b, 3c respectivelly. As a result one gets the three sets of the $`N=4`$ superpartner supergraphs. Each such a set can be studied by the same methos as the supergraphs on Figs. 2a - 2c. It leads to conclusion that sum of the contributions of the supergraphs in every set is equal to zero.
$`\frac{1}{2}`$ $`\frac{1}{2}`$ $`Fig.3a`$ $`\frac{1}{2}`$
$`\frac{1}{2}`$ $`\frac{1}{2}`$ $`Fig.3b`$ $`\frac{1}{2}`$
$`\frac{1}{2}`$ $`\frac{1}{2}`$ $`Fig.3c`$ $`\frac{1}{2}`$
Now we turn to remaining three-loop supergraphs including six propagators which are given by Fig. 1b. There are also two pairs of supergraphs analogous to this pair, they include ghost and gauge propagators instead of matter propagators. After $`D`$-algebra transformations contributions of both these supergraphs and their analogs in which matter loop is replaced by gauge and ghost loop respectively turn to be proportional to the following integral over internal momenta:
$`J`$ $`=`$ $`tr{\displaystyle }d^8\theta {\displaystyle }{\displaystyle \frac{d^4k_1d^4k_2d^4k_3}{(2\pi )^{12}}}{\displaystyle \frac{1}{(k_1^2+𝒲\overline{𝒲})(k_2^2+𝒲\overline{𝒲})(k_3^2+𝒲\overline{𝒲})}}\times `$ (20)
$`\times `$ $`{\displaystyle \frac{1}{((k_1+k_2)^2+𝒲\overline{𝒲})((k_1+k_3)^2+𝒲\overline{𝒲})((k_2+k_3)^2+𝒲\overline{𝒲})}}`$
Here $`tr`$ is a matrix trace. We use the expression for $`N=2`$ background superfield strength in the form (8), the operator $`\stackrel{}{\mathrm{}}`$ in the form (9). The integral in the expression (20) is formally logarithmically divergent (although supergraphs without matter legs and legs including derivatives of gauge strengths have superficial degree of divergence equal to zero , the corresponding divergences vanish due to supersymmetry). Therefore we carry out dimensional regularization via changing integration over $`d^4k_i`$ by integration over $`d^{4+ϵ}k_i`$ (with $`i=1,2,3`$). Straightforward calculation of $`J`$ leads to the result
$`J=tr{\displaystyle \frac{1}{(16\pi ^2)^3}}{\displaystyle d^8\theta (\frac{2}{ϵ}+\mathrm{log}(\frac{𝒲\overline{𝒲}}{\mu ^2}))}`$ (21)
However pole part vanishes due to known properties of integration over anticommuting variables (see f.e. ). Since $`𝒲,\overline{𝒲}`$ are the diagonal matrices they commute with each other, and we can use identity $`d^8\theta \mathrm{log}(\frac{𝒲\overline{𝒲}}{\mu ^2})=d^8\theta (\mathrm{log}(\frac{𝒲}{\mu })+\mathrm{log}(\frac{\overline{𝒲}}{\mu }))`$. This expression vanishes due to chirality of $`𝒲`$. Therefore the supergraphs given by Fig. 1b give zero contribution. The same situation takes place for analogous supergraphs where the matter propagators are replaced by gauge and ghost ones. Hence we conclude that three-loop contribution to non-holomorphic effective potential is equal to zero.
Now let us consider four-loop supergraphs. It turns to be that the situation we observed at three-loop order takes place also at four-loop order, i.e. each supergraph either has $`N=4`$ superpartners sum together with which it is equal to zero or gives the contribution proportional to $`d^{12}z\mathrm{log}(\frac{𝒲\overline{𝒲}}{\mu ^2})=0`$.
First of all, at four-loop order we get two systems of supergraphs given by Figs. 4.1 – 4.6 and Fig.5 which can be separated into sets of $`N=4`$ superpartners. The straightforward calculations show that the combinatoric factors of different supergraphs at Figs. 4.1-4.6 are proportional to each other with the relative factors 1 or 1/2. These factors are written near of each supergraph.
$`Fig\mathrm{.4.1}`$ $`\frac{1}{2}`$ $`Fig\mathrm{.4.2}`$ $`\frac{1}{2}`$ $`Fig\mathrm{.4.3}`$ $`\frac{1}{2}`$
$`Fig\mathrm{.4.4}`$ $`1`$ $`Fig\mathrm{.4.5}`$ $`1`$ $`Fig\mathrm{.4.6}`$ $`1`$
The scheme of separation of supergraphs given by Fig. 4.1 – Fig. 4.6 into sets of $`N=4`$ superpartners is given by Fig. 4a – Fig. 4c where the relative factor is also manifestly shown near corresponding supergraph. We again get three sets of the $`N=4`$ superpartner supergraphs as in three-loop case (see Figs. 3a - 3c). Each line on the Figs. 4a – 4c contains the superpartner supergraphs.
$`\frac{1}{2}`$ $`\frac{1}{2}`$ $`Fig.4a`$ $`\frac{1}{2}`$
$`\frac{1}{2}`$ $`\frac{1}{2}`$ $`Fig.4b`$ $`\frac{1}{2}`$
$`\frac{1}{2}`$ $`\frac{1}{2}`$ $`Fig.4c`$ $`\frac{1}{2}`$
It is evident that the total contribution of the supergraphs given by Figs. 4a – 4c is equal to that one of the supergraphs given by Figs. 4.1 – 4.6. Since the supergraphs within every set are the superpartners their total contribution vanishes. It is proved by the same method as three loop order.
Another system of $`N=4`$ superpartner supergraphs at four-loop order is given by Fig.5.
$`Fig.5`$
Since these supergraphs are $`N=4`$ superpartners, it is easy to see that sum of contributions of these supergraphs is equal to zero.
The remaining four-loop supergraphs with eight propagators are given by Figs. 6a – 6c.
$`Fig.6a`$ $`Fig.6b`$ $`Fig.6c`$
The contributions of these supergraphs after $`D`$-algebra transformations are proportional to the following integrals over internal momenta respectively:
$`J_{6a}`$ $`=`$ $`tr{\displaystyle }d^8\theta {\displaystyle }{\displaystyle \frac{d^4k_1d^4k_2d^4k_3d^4k_4}{(2\pi )^{16}}}{\displaystyle \frac{1}{(k_1^2+𝒲\overline{𝒲})(k_2^2+W\overline{W})(k_3^2+𝒲\overline{𝒲})}}\times `$
$`\times `$ $`{\displaystyle \frac{1}{((k_1+k_2)^2+𝒲\overline{𝒲})((k_1+k_2+k_3)^2+𝒲\overline{𝒲})((k_1+k_4)^2+𝒲\overline{𝒲})}}\times `$
$`\times `$ $`{\displaystyle \frac{1}{(k_4^2+𝒲\overline{𝒲})((k_2+k_3k_4)^2+𝒲\overline{𝒲})}}`$
$`J_{6b}`$ $`=`$ $`tr{\displaystyle }d^8\theta {\displaystyle }{\displaystyle \frac{d^4k_1d^4k_2d^4k_3d^4k_4}{(2\pi )^{16}}}{\displaystyle \frac{1}{(k_1^2+𝒲\overline{𝒲})(k_2^2+𝒲\overline{𝒲})(k_3^2+𝒲\overline{𝒲})}}\times `$
$`\times `$ $`{\displaystyle \frac{1}{((k_1+k_2)^2+𝒲\overline{𝒲})((k_1+k_3)^2+𝒲\overline{𝒲})((k_3+k_4)^2+𝒲\overline{𝒲})}}\times `$
$`\times `$ $`{\displaystyle \frac{1}{(k_4^2+W\overline{W})((k_1+k_3+k_4)^2+𝒲\overline{𝒲})}}`$
$`J_{6c}`$ $`=`$ $`tr{\displaystyle }d^8\theta {\displaystyle }{\displaystyle \frac{d^4k_1d^4k_2d^4k_3d^4k_4}{(2\pi )^{16}}}{\displaystyle \frac{1}{(k_1^2+𝒲\overline{𝒲})(k_2^2+𝒲\overline{𝒲})(k_3^2+𝒲\overline{𝒲})}}\times `$ (22)
$`\times `$ $`{\displaystyle \frac{1}{((k_1+k_2)^2+𝒲\overline{𝒲})((k_2+k_3)^2+𝒲\overline{𝒲})((k_2+k_3k_4)^2+𝒲\overline{𝒲})}}\times `$
$`\times `$ $`{\displaystyle \frac{1}{(k_4^2+𝒲\overline{𝒲})((k_1+k_2+k_4)^2+𝒲\overline{𝒲})}}`$
After dimensional regularization and integration the $`J_{6a}`$, $`J_{6b}`$, $`J_{6c}`$ turn to be equal to
$`J_{6a}=J_{6b}=J_{6c}=tr{\displaystyle \frac{1}{(16\pi ^2)^4}}{\displaystyle d^8\theta (\frac{2}{ϵ}+\mathrm{log}(\frac{𝒲\overline{𝒲}}{\mu ^2}))}`$ (23)
This expression is completely analogous to (21). As a result we get the same conclusion. Each of supergraphs given by Figs. 6a – 6c is proportional to $`d^{12}z\mathrm{log}(\frac{𝒲\overline{𝒲}}{\mu ^2})`$ where $`𝒲`$ is diagonal matrix of the form (8) and $`\overline{𝒲}`$ is its conjugate. The same situation takes place for the analogous supergraphs where the matter superpropagators are replaced by gauge and ghost ones. Hence their contributions also vanish. By the way, we convinced that terms of the form $`\frac{𝒲_a𝒲_b}{𝒲_c𝒲_d}`$ supposed in do not arise in non-holomorphic effective potential at least at three and four loops.
We investigated all four-loop supergraphs including eight internal lines. Namely for this number of superpropagators all $`D`$-factors are used for contracting loops to points in $`\theta `$-space. The four-loop supergraphs with nine superpropagators are also present but after $`D`$-algebra transformations in such supergraphs the extra factor $`(𝒟^+)^4`$ remains. It can act only on background superfield strengths. Therefore four-loop supergraphs with nine propagators cannot contribute to non-holomorphic effective potential.
To conclude, we have considered three- and four-loop supergarphs contributing to non-holomorphic effective potential and proved that only two situations are possible for these supergraphs: (i) either contribution of this supergraph is proportional to $`d^{12}z\mathrm{log}(\frac{𝒲\overline{𝒲}}{\mu ^2})`$, and such a structure vanishes due to properties of integral in superspace (ii) or such a supergraph has $`N=4`$ superpartner supergraphs, and sum of contribution from three $`N=4`$ superpartner supergraphs is equal to zero because of $`N=4`$ supersymmetry. This result is common for any unitary gauge group broken down to its maximally symmetric torus since vanishing of sum of superpartner supergraphs is caused only by $`N=4`$ supersymmetry, not by structure of gauge group.
We found that the mechanism of vanishing of corrections to non-holomorphic effective potential at three-loop order essentially differs from that one at two loops. Absence of two-loop contributions is stipulated only by $`N=2`$ supersymmetry. $`N=4`$ supersymmetry begins to work efficiently at three loops and higher and manifests itself by means of $`N=4`$ superpartner supergraphs. However we proved that the situation at four loops is completely analogous to one in the previous order. The mechanism of vanishing the three- and four-loop contributions to non-holomorphic effective potential looks like very generic and one can expect that it works at any loop.
Detailed study of the structure of above three- and four- loop supergraphs and explicit results of their calculations will be published elsewhere.
Acknowledgements. The authors are glad to acknowledge S.M. Kuzenko for valuable discissions and remarks and A.A. Tseytlin for interest to the work and remarks. The work was supported in part by the RFBR grant 99-02-16617, by the DFG-RFBR grant 99-02-04022, by the INTAS grant No 96-0308 and by GRACENAS grant 97-6.2-34. I.L.B. is grateful to the NATO research grant PST.CLG 974965, to the FAPESP grant and to Institute of Physics, University of São Paulo for hospitality.
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# Collective Singlet Excitations and Evolution of Raman Spectral Weights in the 2D Spin Dimer Compound SrCu2(BO3)2
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## Abstract
Raman light scattering of the two-dimensional quantum spin system SrCu<sub>2</sub>(BO<sub>3</sub>)<sub>2</sub> shows a rich structure in the magnetic excitation spectrum, including several well-defined bound state modes at low temperature, and a scattering continuum and quasielastic light scattering contributions at high temperature. The key to the understanding of the unique features of SrCu<sub>2</sub>(BO<sub>3</sub>)<sub>2</sub> is the presence of strong interactions between well-localized triplet excitations in the network of orthogonal spin dimers realized in this compound.
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Low-dimensional quantum spin systems with a quantum disordered ground state and a finite spin gap form a subject which is of both fundamental and applied interest to the physics community. This is due to the fascinating and diverse physics of the spin-liquid state, allowing to address issues related to the nature of quasiparticle excitations and the role of strong interaction effects. It is known that a gapped singlet ground state is realized in one dimension (1D) in dimerized or frustrated spin chains, and in even-leg spin ladders. Of even greater interest is the existence of a gap in two-dimensional (2D) spin systems, because of its potential relevance for the description of high-temperature superconductivity. Unfortunately, only very few 2D systems, such as CaV<sub>4</sub>O<sub>9</sub>, have been investigated, in which the spin gap is relatively large and not caused by anisotropies. Recently, the compound SrCu<sub>2</sub>(BO<sub>3</sub>)<sub>2</sub> with a layered structure was identified as a 2D S=1/2 Heisenberg system with a unique exchange topology leading to an exact dimer ground state, thus providing the opportunity to study the excitation spectrum of such a model quantum many-body system.
SrCu<sub>2</sub>(BO<sub>3</sub>)<sub>2</sub> has a tetragonal unit cell (D<sub>2d</sub>) with Cu<sup>2+</sup> ions that carry a localized spin S=1/2. The spins form dimers which consist of neighboring pairs of planar rectangular CuO<sub>4</sub> plaquettes with a Cu-Cu distance of 2.9 Å. The strength of the antiferromagnetic intradimer exchange coupling is estimated to be J<sub>1</sub>=100 K . The spin dimers are connected orthogonally by a triangular BO<sub>3</sub> unit. The distance between next-nearest Cu neighbors is 5.1 Å, and the interdimer exchange J<sub>2</sub>$``$68 K . The orthogonal arrangement of dimers thus represents a 2D frustrated quantum spin system . A sketch of the Cu-dimers is shown in the inset of Fig.1. The ground state in this exchange topology depends critically on the ratio J<sub>2</sub>/J<sub>1</sub> . For small J<sub>2</sub> the system consists of isolated dimers and the ground state is a product of singlets, while for small J<sub>1</sub> the model can be mapped on a 2D square lattice of spins and a Néel-ordered state is expected. It has been shown that the critical ratio of coupling constants that separates a gapfull and a gapless state is $`J_2/J_1=0.7`$ . The experimental ratio of $`J_2/J_10.68`$ is just below this value, placing this material in the dimerized phase, close to the Néel boundary. Thermodynamic measurements support the existence of a dimer ground state and a spin gap of $`\mathrm{\Delta }`$=34 K . The magnetic specific heat and susceptibility both show maxima (T$`{}_{}{}^{\chi }{}_{max}{}^{}`$=15 K, T$`{}_{}{}^{c_p}{}_{max}{}^{}`$=8 K) and a rapid decrease toward lower temperatures with exponential tails at $`T0`$.
We have investigated the magnetic excitation spectrum of single crystals of SrCu<sub>2</sub>(BO<sub>3</sub>)<sub>2</sub> by Raman scattering experiments. At low temperatures (T$``$$`\mathrm{\Delta }`$) we observe well-defined low-energy singlet modes which can be interpreted as collective bound state excitations of two and three elementary triplets. We demonstrate that the appearance of the collective modes reflects the strong triplet-triplet interactions present in the system. The spectral weight evolution for intermediate (T$``$$`\mathrm{\Delta }`$) and high (T$`>`$$`\mathrm{\Delta }`$) temperatures can be described by the temperature dependence of the magnetic susceptibility and the magnetic specific heat.
The samples were investigated in quasi-backscattering geometry with light polarizations in the $`ab`$-plane of the freshly cleaved crystal. The ($`a^{}b^{}`$) axes are rotated by 45 within this plane. The experiments used the $`\lambda `$ = 488-nm excitation line of an Ar<sup>+</sup> ion laser and a laser fluence below 20 W/cm<sup>2</sup>. The scattered light was analyzed using a XY-Dilor Raman spectrometer and a back-illuminated CCD detector. Measurements in a magnetic field were performed in 90 scattering geometry. In the analysis of our results the main emphasis has been on the role of interactions and their influence on the magnetic excitation spectrum. A detailed analysis of the phonon spectrum will be given elsewhere.
In the low energy region with Raman shifts comparable to the triplet gap $`\mathrm{\Delta }`$=24cm<sup>-1</sup>, drastic changes and a large shift of spectral weight occur with decreasing temperature. This evolution leads finally to the appearance of several new modes. In Fig.1 Raman spectra of SrCu<sub>2</sub>(BO<sub>3</sub>)<sub>2</sub> are shown to illustrate these effects at different temperatures and in two scattering configurations with light polarizations within the ab-plane of the crystal. The spectra in the upper panel (a) present data in $`(ab)`$:B<sub>2</sub> and the lower panel (b) shows data in $`(a^{}b^{})`$:B<sub>1</sub> scattering configuration. Symmetry components are denoted with respect to the D<sub>2d</sub> point group using $`b^{}`$=$`a`$+$`b`$.
At low temperatures (T$``$$`\mathrm{\Delta }`$) four well-defined modes with energies $`E^S`$= 30, 46, 56 and 70 cm<sup>-1</sup> appear. The dominant intensity of these modes is observed in the $`(a^{}b^{})`$ scattering configuration. They neither split nor shift in an applied magnetic field up to 6 T as shown in the inset of Fig.1(b) and therefore are assigned to spin singlets. The only effect of the magnetic field is observed as a shift of a weak intensity shoulder at 24 cm<sup>-1</sup> toward lower frequencies. This signal corresponds to the elementary spin gap. At higher temperatures (T$``$$`\mathrm{\Delta }`$) all modes are strongly damped. They are replaced by a continuum of scattering with a center of gravity near 50 cm<sup>-1</sup> (Fig.1(a),T=7 K), corresponding roughly to the energy 2$`\mathrm{\Delta }`$. For even higher temperatures (T$`>`$$`\mathrm{\Delta }`$) quasielastic scattering with a Lorentzian spectral function is detected. The latter two scattering intensities are observed in the $`(ab)`$ scattering configuration.
To understand these dramatic changes of the Raman spectra with temperature we compare them with observations related to the triplet excitation spectrum and then map the temperature dependence of the scattering intensity onto the corresponding thermodynamic data, i.e. the magnetic susceptibility and the magnetic part of the specific heat. To gain a deeper insight into the structure of the spectrum we will also identify the interactions leading to the formation of singlet bound states in the relevant Heisenberg model at T=0, and present estimates for their binding energies.
Recent ESR and neutron scattering investigations on SrCu<sub>2</sub>(BO<sub>3</sub>)<sub>2</sub> at low temperatures observed triplet excitations with a spin gap of $`\mathrm{\Delta }`$=34 K. This triplet branch has a very small dispersion of only 2 K pointing to extremely localized excitations . In addition, a second triplet branch $`\mathrm{\Delta }^{}`$=55 K with a larger dispersion of 17 K was observed. This branch can be interpreted as a triplet bound state of two elementary triplets (see discussion below). Frustration due to the interdimer coupling J<sub>2</sub> can lead to the reduction of the ratio $`\mathrm{\Delta }^{}`$/$`\mathrm{\Delta }`$=1.62 below 2 (corresponding to non-interacting magnons) .
The four modes that we observe in Raman scattering in SrCu<sub>2</sub>(BO<sub>3</sub>)<sub>2</sub> for T$``$$`\mathrm{\Delta }`$ are clearly related in energy and temperature scale to the spin gap of the compound. A phonon origin of these modes can be ruled out due to their temperature dependence and the nonexistence of a structural phase transition below room temperature. On the other hand, in low-dimensional spin systems with strong triplet-triplet interactions, well defined modes can appear below the scattering continuum . In this case the light scattering process is better described as spin conserving ($`\mathrm{\Delta }`$S=0) scattering on singlet bound states. Such states composed of two elementary triplets have been observed, e.g., in the low-temperature dimerized phases of $`\mathrm{CuGeO}_3`$ and $`\mathrm{NaV}_2\mathrm{O}_5`$ .
To illustrate the mechanism of bound state formation at T=0 in the 2D Heisenberg model $`H=\mathrm{\Sigma }_{i,j}J_{ij}𝐒_i𝐒_j`$ with $`J_{ij}>0`$ defined in the inset of Fig.1(a) (the Shastry-Sutherland model ), we have derived the effective bosonic representation in terms of triplets ($`𝐭_i^{}`$), excited above the singlets, formed on the stronger ($`J_1`$) bonds:
$`H=J_1{\displaystyle \underset{i}{}}𝐭_i^{}𝐭_i+V{\displaystyle \underset{i,j}{}}(𝐭_i^{}\times 𝐭_i)(𝐭_j^{}\times 𝐭_j)+`$ (1)
$`{\displaystyle \underset{iA}{}}\left\{i\mathrm{\Gamma }(𝐭_i^{}\times 𝐭_i)(𝐭_{i+\widehat{x}}^{}𝐭_{i\widehat{x}}^{}+\text{h.c.})\right\}+{\displaystyle \underset{iB}{}}\{\widehat{x}\widehat{y}\},`$ (2)
where $`V=\mathrm{\Gamma }=J_2/2`$. The site indices run over the square lattice formed by the dimers (singlet pairs of spins connected by $`J_1`$ bonds), and $`i,j`$ stands for nearest neighbors, while sub-lattice A(B) is formed by the vertical (horizontal) dimers. An important feature of Eq.(1) is the absence of quantum fluctuations (i.e. $`𝐭_i^{}𝐭_j^{}`$ terms), reflecting the fact that the ground state is an exact product of singlets. The triplet excitations would have been completely localized in the absence of the $`\mathrm{\Gamma }`$ term in Eq.(1), and hopping appears only in order $`(\mathrm{\Gamma }/J_1)^6(J_2/J_1)^6`$, leading to a small bandwidth compared to the gap .
Next we present estimates for the energies of the collective, two-particle excitations. The zero-momentum two-magnon bound state with S=0, constructed by exciting two triplets, has the form: $`|\mathrm{\Psi }^S=_{i,j}\psi _{i,j}^S𝐭_i^{}𝐭_j^{}|0`$. Starting from the limit $`J_2/J_11`$, the dominant contribution to binding comes from the two-particle interaction ($`V`$ term in Eq.(1)), with corrections of order $`(J_2/J_1)^2`$ and higher. This interaction provides an effective attraction between two triplets. Assuming the triplets are localized the binding problem can be solved exactly, and the singlet binding energy defined as $`ϵ^S=2\mathrm{\Delta }E_2^S`$, is $`ϵ^S=Z^2J_2`$. Here $`E_2^S`$ is the energy of the singlet bound state, and we have estimated the renormalization factor $`Z`$ by evaluating the appropriate lowest order diagrams to be $`Z0.75`$ at $`J_2/J_1=0.65`$. We have found, by comparing neutron scattering data for the one-particle dispersion from with the dispersion obtained by high-order perturbative expansions that the ratio $`J_2/J_1=0.65`$ is consistent with the experimental results. Indeed, for this ratio the theoretical bandwidth (BW=difference between the energy at $`𝐤=(\pi ,0)`$ and $`(0,0)`$) and gap are: $`BW0.04J_1,\mathrm{\Delta }0.36J_1`$ , consistent with the measured values $`BW0.3meV`$ and $`\mathrm{\Delta }3meV`$ . From the same analysis we estimate $`J_261K42cm^1`$. For the higher value of $`J_2/J_1=0.678`$ the bandwidth is approximately twice as large which would make this ratio inconsistent with the neutron scattering data. Putting everything together we have the estimate for the singlet binding energy $`ϵ^S=Z^2J_224cm^1`$. Our experimental value for the lowest singlet mode in Fig.1(b) with $`E_2^S`$=30 $`cm^1`$ is $`ϵ_{Raman}^S18cm^1`$.
The analysis based on the localized picture somewhat overestimates the binding since corrections of order $`(J_2/J_1)^2`$ and higher have been neglected. Such high-order corrections lead to the development of a strong dispersion as well as create several two-particle singlet modes with different energies, classified according to the group $`D_{2d}`$. Detailed calculations have recently been performed , providing quantitatively accurate results for the energies in the different symmetry sectors. Only the modes with $`\mathrm{\Gamma }_3(xy)`$ symmetry contribute to the Raman intensity, and, as $`T0`$, their contribution grows for the ($`a^{}b^{}`$) and vanishes for the ($`ab`$) geometry , consistent with our observations. In addition, the calculated energies are in excellent agreement with the two lowest ($`30cm^1`$ and $`46cm^1`$) singlet modes observed in our experiment, which we therefore interpret as two-particle bound states.
In SrCu<sub>2</sub>(BO<sub>3</sub>)<sub>2</sub> the four modes in Fig.1(b) show ratios $`E^S`$/$`\mathrm{\Delta }`$=1.25 – 2.9. Two of the well-pronounced peaks (at $`56cm^1`$ and $`70cm^1`$) are above the two-particle threshold ($`2\mathrm{\Delta }`$). The smallness of the one-particle bandwidth makes it possible to resolve such higher energy peaks, since the weight of the scattering background is expected to be small. The pronounced sharpness of the peaks indicates that they can be interpreted as three-particle excitations, whose existence is indeed possible due to the localized nature of the states. A three-particle singlet bound state can be constructed as: $`|\mathrm{\Phi }^S=_{i,j,k}\varphi _{i,j,k}^S(𝐭_i^{}\times 𝐭_j^{})𝐭_k^{}|0`$, and its energy on the same level of approximation as discussed above (localized triplets) is: $`E_3^S=3\mathrm{\Delta }Z^3J_254cm^1`$, in very good agreement with the energy of the mode at $`56cm^1`$. The appearance of additional levels (as well as finite dispersion) at higher order is also expected. The total Raman scattering intensity, reflecting the presence of two- and three- particle singlets with spectral weights $`I_2`$ and $`I_3`$, is: $`I_B(\omega )=I_2\delta (\omega E_2^S)+I_3\delta (\omega E_3^S)`$. In order to estimate the ratio of the two intensities we notice that while light couples directly to the singlet two-magnon state, the three-magnon state has to be created via the action of the $`\mathrm{\Gamma }`$ (magnon number non-conserving) term in Eq.(1) (i.e. this term provides a vertex correction to the Raman operator $`R𝐭_i^{}𝐭_j^{}`$). Consequently, we estimate at T=0, $`I_3/I_2(J_2/J_1)^20.4`$, in good agreement with the experimental ratio of 0.32 (at T=1.5 K). Notice that in the hypothetical case $`\mathrm{\Gamma }=0`$, when the excitations are completely localized, the three-particle states do not contribute to the Raman intensity. Thus the $`\mathrm{\Gamma }`$ term plays an important role providing both a finite bandwidth and coupling to higher bound state modes. Let us also mention that three-particle bound states have been predicted theoretically in quasi-1D systems . Our work presents the first experimental evidence for their existence in a (2D) dimerized spin system.
The bound state modes experience strong damping with increasing temperature (see Fig. 1(b)), due to scattering on thermally excited triplet states. The decrease of the scattering intensity $`I_B`$ with temperature is governed by the density of excitations: $`I_B(T)(1Ae^{\frac{\mathrm{\Delta }}{k_BT}}),`$ where $`A`$ is a constant. In Fig.2(a) the integrated intensity of the bound states is shown together with a fit based on the equation above ($`A`$=215.7) and fixing $`\mathrm{\Delta }`$=34 K from experimental data. The macroscopic temperature-induced population of triplet states has a direct effect on the excitation spectrum. For intermediate temperatures (5-100 K, typical spectrum shown at 7 K in Fig.1(a)) a broad continuum of scattering is observed that is of magnetic origin . However, in comparison with a 2D antiferromagnet with a maximum of the two-magnon continuum at E<sub>max</sub>=2.8 J its energy is very small. The reason for the maximum appearing at only 50 cm$`{}_{}{}^{1}`$2$`\mathrm{\Delta }`$ is the smallness of the one-particle bandwidth in SrCu<sub>2</sub>(BO<sub>3</sub>)<sub>2</sub>. The temperature dependence of the two-magnon scattering intensity shows a pronounced increase in the temperature range where the intensity of the bound states is dropping (see the open triangles in Fig.2(b)). At higher temperatures (not shown here) it decreases again and is then superimposed on the intensity of quasielastic scattering. In Fig.2(b) we compare the integrated scattering intensity with the behavior of the magnetic susceptibility $`\chi `$(T). Calculations of the latter quantity exist, but their agreement with experiment is not fully satisfactory and consequently we take $`\chi `$(T) directly from experiment . Notice that both the two-magnon continuum intensity and $`\chi `$(T) exhibit strong variations on the same temperature scale, even though we are not aware of a rigorous sum rule relating the two quantities.
Quasielastic scattering connected with fluctuations of the magnetic energy density is found at high temperatures (T$`>`$$`\mathrm{\Delta }`$) in our Raman scattering experiments. It has the expected Lorentzian spectral function . To determine the evolution as a function of temperature we use the hydrodynamic form of the respective correlation function , which includes the magnetic specific heat $`c_m(T)`$ and the thermal diffusion constant $`D_T`$. In the high temperature approximation ($`\mathrm{}\omega /k_BT1`$) the result is : $`I(\omega )\frac{k_B}{\mathrm{}}\frac{c_mT^2D_Tk^2}{\omega ^2+(D_Tk^2)^2}`$ , where $`k`$ is the scattering wave-vector. A fit to this equation can then be used to estimate $`c_m(T)`$ from the integrated intensity, scaled by T<sup>2</sup>. In Fig.2(c) the result of this procedure is found in good agreement with the measured specific heat.
In conclusion, we have shown that the low-energy excitation spectrum of the 2D compound SrCu<sub>2</sub>(BO<sub>3</sub>)<sub>2</sub> has a rich and complex structure. Our main result is that at low temperature (T$``$$`\mathrm{\Delta }`$) the spectrum contains several well pronounced singlet modes which are interpreted as collective two-particle and novel three-particle singlet bound states of strongly localized triplets. We have demonstrated that since quantum fluctuations are absent in the ground state (due to the unique dimer arrangement), the triplet-triplet interactions can lead to large binding energies and scattering intensities, and consequently make the collective modes observable. At intermediate (T$``$$`\mathrm{\Delta }`$) and high (T$``$$`\mathrm{\Delta }`$) temperatures, where the quantum (interaction) effects are not important, the light scattering is dominated by a two-magnon continuum and quasielastic contributions. We relate the spectral weight evolution in these regimes to thermodynamic quantities.
We acknowledge fruitful discussions with C. Pinettes, P.H.M. van Loosdrecht, G.S. Uhrig, C. Gros, R. Valentí, W. Brenig, W. Atkinson, P. Hirschfeld, and D. Tanner. Financial support by DFG/SFB 341 and NSF Grant DMR-9357474 (V.N.K.) is gratefully acknowledged.
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# 1 Introduction
## 1 Introduction
An important advantage of the dynamical symmetry (DS) approach in nuclear theory is the possibility to describe consistently various collective bands of heavy deformed nuclei . Generally, the DS concept is based on the assumption that the physical system possesses a “primary” symmetry with respect to a given group, called DS group. The Hamiltonian of the system reduces this symmetry to the group of invariance of the system (which for the nuclear system coincides with the angular momentum group SO(3)) and thus the energy spectrum is generated . The Lie algebra of the DS group is then reduced to the algebra of the group of invariance and is referred to as spectrum generating algebra. The basic idea of DS approach in heavy deformed nuclei is that their collective bands can be united into one or several multiplets, appearing in this reduction . It provides a natural way to study the interaction between a particular couple of bands as well as the attendant spectroscopic characteristics of nuclei.
Various classification schemes with band coupling have been developed on the basis of DS approach. Well known models, such as the Interacting Boson Model (IBM) , the symplectic models and the Fermion Dynamical Symmetry Model , provide a good overall description of nuclear collective phenomena, covering the different regions of vibrational, rotational and transitional nuclei.
On the other hand, some models, based on the SU(3) dynamical symmetry, reproduce successfully the particular characteristics of rotational bands in deformed nuclei. Such models are the Pseudo-SU(3) Model , which has microscopic motivations, as well as the Vector-Boson Model (VBM) with SU(3) dynamical symmetry , which allows a relevant phenomenological treatment of the SU(3) multiplets in nuclei.
While in the SU(3) limit of the IBM the possible irreducible representations (irreps) ($`\lambda ,\mu `$) are restricted by the total number of bosons describing the specific nucleus, in the VBM the possible SU(3) irreps ($`\lambda ,\mu `$) are not restricted by the underlying theory. However, it has been shown recently that some favored regions of ($`\lambda ,\mu `$) multiplets in the VBM could be outlined through the numerical analysis of the experimental data available for the ground ($`g`$) and the $`\gamma `$\- collective bands of even–even deformed nuclei. (The favored multiplets provide the best model descriptions.) As a result, a systematic behavior of the SU(3) symmetry properties of rotational nuclei has been established in terms of the VBM. It suggests the presence of a transition between a scheme, in which the $`g`$ and the $`\gamma `$ bands are coupled into one and the same ($`\lambda ,\mu `$) irrep and a scheme, where these two bands belong to different irreps. In addition it has been supposed that the fine systematic properties of rotational spectra could be interpreted as a manifestation of a more general dynamical symmetry.
As a first step in the recovering of the dynamical mechanism causing such a transition, one should study the way in which the SU(3) symmetry is reduced in the ($`\lambda ,\mu `$)- plane. In particular, it is of interest to reproduce the limits, in which the quantum numbers $`\lambda `$ and $`\mu `$ go to infinity, i.e. the cases, in which the SU(3) irreps are not finite anymore. These limits correspond to the so called SU(3) contraction process, in which the algebra of SU(3) goes to the algebra of the semi-direct product T$`{}_{5}{}^{}`$SO(3), i.e. $`SU(3)T_5SO(3)`$ (T<sub>5</sub> is the group of 5-dimensional translations generated by the components of the SU(3)- quadrupole operators) . Generally, the contraction limit corresponds to a singular linear transformation of the basis of a given Lie algebra. The transformed structure constants approach a well-defined limits and a new Lie algebra, called contracted algebra, results . The original and the contracted algebra are not isomorphic.
On the above basis it is expected that in the SU(3) contraction limit the space of the SU(3) irreps should undergo a respective limiting transition. As a result the SU(3) multiplets should be disintegrated to sets of various independent bands. It is, therefore, reasonable to consider this limit as a natural way in which the band-mixing interactions vanish. It is important to remark that the SU(3) contraction process is a situation in which a compact group goes to a non-compact one. Hence, one could try to interpret the vanishing $`g`$$`\gamma `$ band-mixing interaction as a transition from a compact to a non-compact DS group.
In the present work we realize the above considerations through the formalism of the VBM. Our purpose is to examine the various directions in the ($`\lambda ,\mu `$)-plane by investigating the respective changes in the structure of the SU(3) multiplets in terms of model defined spectroscopic characteristics of rotational nuclei. As such, we consider here the SU(3) energy splitting and the $`g`$$`\gamma `$ interband transitions, which carry important information about the link between the two bands. It is known that the energy splitting of the multiplet determines to a great extent the systematic behavior of the SU(3) dynamical symmetry in deformed nuclei .
In the VBM relatively simple analytic expressions for the energies and the transition probabilities can be derived both for the lowest $`L=2`$ states of any ($`\lambda ,\mu `$)- multiplet and for all the states of any ($`\lambda ,2`$) multiplet. The analytic expressions for the $`L=2`$ states allow one to examine the SU(3) characteristics of nuclei in terms of two-dimensional surfaces in the ($`\lambda ,\mu `$)-plane, while in the ($`\lambda ,2`$) direction one is able to investigate the behavior of the full set of states in the multiplet, i.e. the states with $`L2`$. On the other hand, the $`L2`$ states of the irreps with $`\mu >2`$ can be treated numerically.
In such a way, a relevant combination of analytic and numerical analyses could be applied in order to reveal the systematic behavior of all the states of SU(3) irreps in the ($`\lambda ,\mu `$)-plane including the limiting cases of SU(3) group contraction. The collective scheme of the VBM is constructed by using the irreps with $`\lambda \mu `$ and comprises the following two SU(3) contraction limits:
(i) $`\lambda \mathrm{}`$, with $`\mu `$ finite;
(ii) $`\lambda \mathrm{}`$, $`\mu \mathrm{}`$, with $`\mu \lambda `$.
Below we provide a detailed study of the most important spectroscopic characteristics of the $`g`$ and the $`\gamma `$ band in the above limiting cases. It will be shown that our approach gives a reasonable interpretation of the corresponding experimental data and leads to rather clear conclusions about the rearrangement of collective rotational bands in heavy deformed nuclei.
In Sec. II the $`g`$$`\gamma `$ band coupling scheme of the VBM is briefly presented. In Sec. III we derive analytic expressions for the energies and the $`B(E2)`$-transition probabilities for the $`2_g`$ and $`2_\gamma `$ states of an arbitrary ($`\lambda ,\mu `$) multiplet. Using them, we obtain the analytic behavior of the energy splitting and the physically meaningful transition ratios in the SU(3) contraction limits (i) and (ii). In Sec. IV we derive expressions for the splitting and the transition ratios for the full set of states ($`L2`$) in the ($`\lambda ,2`$) multiplets and obtain their analytic form in the first limiting case ($`\lambda \mathrm{}`$; $`\mu =2`$). In Sec. V all analytic results are examined numerically. Also, there we provide a numerical study of the second limiting case ($`\lambda \mathrm{}`$, $`\mu \mathrm{}`$; $`\mu \lambda `$) for the states with $`L2`$. The results are discussed together with an analysis of experimental data. In Sec. VI the conclusions are given.
## 2 $`g`$$`\gamma `$ band coupling in the VBM
The Vector-Boson Model (VBM) with SU(3) dynamical symmetry is founded on the assumption that the low-lying collective states of deformed even–even nuclei can be described by means of two distinct kinds of vector bosons, whose creation operators $`𝝃^\mathbf{+}`$ and $`𝜼^\mathbf{+}`$ are O(3) vectors and in addition transform according to two independent SU(3) irreps of the type $`(\lambda ,\mu )=(1,0)`$ . The vector bosons provide a relevant construction of the SU(3) angular momentum and quadrupole operators like the bosons in the Schwinger realization of SU(2) . Therefore, they can be considered as natural building blocks of a model scheme with SU(3) dynamical symmetry. Also, the vector bosons can be interpreted as quanta of elementary collective excitations of the nucleus .
In this model an SU(3)-symmetry reducing Hamiltonian is constructed by using three basic O(3) scalars, which belong to the enveloping algebra of SU(3) :
$$V=g_1L^2+g_2LQL+g_3A^+A.$$
(1)
Here $`g_1`$, $`g_2`$ and $`g_3`$ are free parameters; $`L`$ and $`Q`$ are the angular momentum and quadrupole operators respectively; and $`A^+=𝝃_{}^{\mathbf{+}}{}_{}{}^{2}𝜼_{}^{\mathbf{+}}{}_{}{}^{2}(𝝃^\mathbf{+}𝜼^\mathbf{+})^2`$.
The basis states
$$|\begin{array}{c}(\lambda ,\mu )\\ \alpha ,L,M\end{array},$$
(2)
corresponding to the $`SU(3)O(3)`$ group reduction, are constructed by means of the above vector–boson operators and are known as the basis of Bargmann–Moshinsky .
The quantum number $`\alpha `$ in Eq. (2) distinguishes the various O(3) irreps, $`(L,M)`$, appearing in a given SU(3) irrep $`(\lambda ,\mu )`$ and labels the different bands of an SU(3) multiplet. It is an integer number determined through the following inequality
$$\mathrm{max}\{0,\frac{1}{2}(\mu L)\}\alpha \mathrm{min}\{\frac{1}{2}(\mu \beta ),\frac{1}{2}(\lambda +\mu L\beta )\},$$
(3)
where
$$\beta =\{\begin{array}{cc}0,\hfill & \lambda +\mu L\text{ even}\hfill \\ 1,\hfill & \lambda +\mu L\text{ odd}\hfill \end{array}$$
In the VBM the $`g`$\- and the lowest $`\gamma `$\- band belong to one and the same SU(3) multiplet, where $`\lambda `$ and $`\mu `$ are even and $`\lambda \mu `$. These bands are labeled by two neighboring integer values of the quantum number $`\alpha `$. (More precisely, the states of the $`g`$\- band are labeled by the largest value of $`\alpha `$ appearing in (3), while the $`\gamma `$\- band corresponds to the next smaller $`\alpha `$\- value.) The so defined multiplet is split with respect to $`\alpha `$.
The above scheme provides a good description of the energy levels and of the B(E2) transition ratios within and between the $`g`$\- and $`\gamma `$\- bands . The other collective bands, in particular the lowest $`\beta `$-band, do not belong to the same irrep. Therefore, they are not considered in the framework of this model.
## 3 The $`L=2`$ states in ($`\lambda ,\mu `$)-plane
### 3.1 $`L=2`$ energy splitting
Here we consider the $`L=2`$ energy levels of the $`g`$\- and the $`\gamma `$\- band in terms of the VBM. For any ($`\lambda ,\mu `$) multiplet ($`\mu 2`$), the $`2_g`$ and $`2_\gamma `$ states are the only possible ones appearing at angular momentum $`L=2`$. They are labeled by the quantum number $`\alpha `$ as follows \[See inequality (3)\]: $`\alpha _1=\mu /21`$ for $`2_\gamma `$ and $`\alpha _2=\mu /2`$ for $`2_g`$. Hence, for the $`L=2`$ states the Hamiltonian matrix is always two-dimensional and the corresponding eigenvalue equation has the form:
$$\mathrm{det}\left(\begin{array}{cc}V_{1,1}\omega ^{(2)}& V_{1,2}\\ V_{2,1}& V_{2,2}\omega ^{(2)}\end{array}\right)=0,$$
(4)
where $`\omega ^{(2)}\omega ^{L=2}`$ are the eigenvalues and
$$V_{j,j^{}}\alpha _j,2|V|\alpha _j^{},2=\begin{array}{c}(\lambda ,\mu )\\ \alpha _j,2,2\end{array}\left|V\right|\begin{array}{c}(\lambda ,\mu )\\ \alpha _j^{},2,2\end{array},$$
(5)
with $`j,j^{}=1,2`$, are the corresponding Hamiltonian matrix elements. We have derived these matrix elements in the form:
$`V_{1,1}`$ $`=`$ $`({\displaystyle \frac{\mu }{2}}1),2|V|({\displaystyle \frac{\mu }{2}}1),2=6g_1+6g_2(2\lambda +2\mu +3)+g_3P(\lambda ,\mu ),`$ (6)
$`V_{2,2}`$ $`=`$ $`{\displaystyle \frac{\mu }{2}},2|V|{\displaystyle \frac{\mu }{2}},2=6g_16g_2(2\lambda +2\mu +3)+g_3Q(\lambda ,\mu ),`$ (7)
$`V_{1,2}`$ $`=`$ $`({\displaystyle \frac{\mu }{2}}1),2|V|{\displaystyle \frac{\mu }{2}},2=12g_2\mu 2g_3\mu (\mu 2),`$ (8)
$`V_{2,1}`$ $`=`$ $`{\displaystyle \frac{\mu }{2}},2|V|({\displaystyle \frac{\mu }{2}}1),2=12g_2\lambda +2g_3\lambda (\lambda +2\mu +2),`$ (9)
where
$`P(\lambda ,\mu )`$ $`=`$ $`\lambda (\mu 2)(\mu +2)(\lambda +2\mu +2)+\mu (\mu 2)(\mu +1)(\mu +3),`$ (10)
$`Q(\lambda ,\mu )`$ $`=`$ $`\lambda \mu ^2(\lambda +2\mu +2)+\mu (\mu 1)(\mu +1)(\mu +2).`$ (11)
The energy levels $`E_2^g`$ and $`E_2^\gamma `$, corresponding to the $`2_g`$ and $`2_\gamma `$ states respectively, are determined as
$`E_2^g`$ $`=`$ $`\omega _1^{(2)}\omega ^{(0)},`$ (12)
$`E_2^\gamma `$ $`=`$ $`\omega _2^{(2)}\omega ^{(0)},`$ (13)
where
$$\omega _i^{(2)}=\frac{1}{2}\left\{V_{1,1}+V_{2,2}+(1)^i\sqrt{(V_{1,1}+V_{2,2})^24(V_{1,1}V_{2,2}V_{1,2}V_{2,1})}\right\},$$
(14)
$`i=1,2`$, are the solutions of the eigenvalue equation (4), and $`\omega ^{(0)}=g_3\mu ^2(\lambda +\mu +1)^2`$ is the zero-level eigenvalue, corresponding to the ground state $`0_g`$. After using Eqs. (6)–(9) we obtain the following analytic expressions for $`E_2^g`$ and $`E_2^\gamma `$:
$`E_2^g`$ $`=`$ $`6g_12Fg_32\sqrt{Ag_2^2+Bg_3^2Cg_2g_3},`$ (15)
$`E_2^\gamma `$ $`=`$ $`6g_12Fg_3+2\sqrt{Ag_2^2+Bg_3^2Cg_2g_3},`$ (16)
where
$`A`$ $`=`$ $`A(\lambda ,\mu )=9[(2\lambda +2\mu +3)^24\lambda \mu ];`$ (17)
$`B`$ $`=`$ $`B(\lambda ,\mu )=[\lambda (\lambda +2\mu +2)+\mu (\mu +1)]^2`$ (18)
$``$ $`\lambda \mu (\lambda +2\mu +2)(\mu 2);`$
$`C`$ $`=`$ $`C(\lambda ,\mu )=6(2\lambda +2\mu +3)[\lambda (\lambda +2\mu +2)+\mu (\mu +1)]`$ (19)
$``$ $`6\lambda \mu (\lambda +3\mu );`$
$`F`$ $`=`$ $`F(\lambda ,\mu )=\lambda (\lambda +2\mu +2)+2\mu (\mu +1).`$ (20)
Hence, we derive a model expression for the energy splitting of the SU(3) multiplet. It is known that the splitting can be characterized by the ratio :
$$\mathrm{\Delta }E_2=\frac{E_2^\gamma E_2^g}{E_2^g}.$$
(21)
In terms of Eqs. (15) and (16) the quantity $`\mathrm{\Delta }E_2`$ obtains the following analytic form:
$$\mathrm{\Delta }E_2=\frac{2}{(3g_1Fg_3)/\sqrt{Ag_2^2+Bg_3^2Cg_2g_3}1}.$$
(22)
The expressions, obtained so far, allow us to study analytically the $`g`$$`\gamma `$ band-mixing interaction and the energy splitting at $`L=2`$ in the ($`\lambda ,\mu `$)-plane. In particular we are able to reproduce analytically the SU(3) contraction limits:
(i) $`\lambda \mathrm{}`$, with $`\mu `$ finite;
(ii) $`\lambda \mathrm{}`$, $`\mu \mathrm{}`$, with $`\mu \lambda `$. Since the difference $`\lambda \mu `$ is always finite, we take for definiteness $`\mu =\lambda `$.
In each of these limits we estimate the $`\lambda `$\- and/or $`\mu `$\- dependence of the matrix elements (6)–(9), as well as the analytic behavior of the splitting ratio $`\mathrm{\Delta }E_2`$.
In case (i) the matrix elements are determined by the corresponding highest degrees of $`\lambda `$. Thus for $`\mu >2`$ the Hamiltonian matrix $`(V_{i,j})`$ obtains the following asymptotic form:
$$(V)_\lambda \mathrm{}\left(\begin{array}{cc}\lambda ^2& \\ \lambda ^2& \lambda ^2\end{array}\right),$$
(23)
where the upper off-diagonal element (denoted by $``$) does not depend on $`\lambda `$. Then the relative contribution of the off-diagonal (band-mixing) terms in the eigenvalue equation (4) decreases with the increase of $`\lambda `$ as $`\lambda ^2/\lambda ^4=1/\lambda ^2`$. For $`\mu =2`$ the term $`V_{1,1}`$ is proportional to $`\lambda `$ instead of $`\lambda ^2`$ \[See Eqs. (6) and (10)\], so that in this particular case the off-diagonal contribution decreases as $`1/\lambda `$.
In the same limiting case the functions (17)–(20) have the following asymptotic behavior:
$$A_\lambda \mathrm{}=36\lambda ^2;B_\lambda \mathrm{}=\lambda ^4;C_\lambda \mathrm{}=12\lambda ^3;F_\lambda \mathrm{}=\lambda ^2.$$
After applying them in Eq. (22), we find the analytic limit of the splitting ratio (22):
$$\underset{\lambda \mathrm{}}{lim}\mathrm{\Delta }E_2=\frac{2}{g_3/|g_3|1}.$$
(24)
We remark that the application of the VBM in rare earth nuclei and actinides requires $`g_3<0`$ , which gives in (24)
$$\underset{\lambda \mathrm{}}{lim}\mathrm{\Delta }E_2=\mathrm{}.$$
(25)
Therefore, in this case the SU(3)-multiplet is completely split.
Consider now the limiting case (ii), $`\lambda =\mu \mathrm{}`$. Then the asymptotic form of the matrix $`(V_{i,j})`$ is:
$$(V)_{\lambda =\mu \mathrm{}}\left(\begin{array}{cc}\lambda ^4& \lambda ^2\\ \lambda ^2& \lambda ^4\end{array}\right).$$
(26)
Here we find that the relative magnitude of the band-mixing interaction decreases as $`\lambda ^4/\lambda ^8=1/\lambda ^4`$, i.e., more rapidly in comparison to the previous case.
Furthermore, in the limiting case (ii) one has:
$$A_{\lambda =\mu \mathrm{}}=108\lambda ^2;B_{\lambda =\mu \mathrm{}}=13\lambda ^4;C_{\lambda =\mu \mathrm{}}=72\lambda ^3;F_{\lambda =\mu \mathrm{}}=5\lambda ^2.$$
Then the SU(3) splitting ratio goes to:
$$\underset{\lambda =\mu \mathrm{}}{lim}\mathrm{\Delta }E_2=\frac{2}{(5/\sqrt{13})g_3/|g_3|1}.$$
(27)
For $`g_3<0`$ we obtain
$$\underset{\lambda =\mu \mathrm{}}{lim}\mathrm{\Delta }E_2=2/(5/\sqrt{13}1)=5.17.$$
(28)
Therefore, in this case the band-mixing interaction vanishes, while the energy splitting between the two bands remains finite.
### 3.2 Transition ratios in the $`L=2`$ states
Here we turn to the electromagnetic transition probabilities for the states $`2_g`$, $`2_\gamma `$ and $`0_g`$. In particular it is of interest to consider the following B(E2) transition ratios:
$`R_1(2)`$ $`=`$ $`{\displaystyle \frac{B(E2;2_\gamma 2_g)}{B(E2;2_g0_g)}};`$ (29)
$`R_2(2)`$ $`=`$ $`{\displaystyle \frac{B(E2;2_\gamma 2_g)}{B(E2;2_\gamma 0_g)}}.`$ (30)
The first of them, $`R_1(2)`$, gives the relative magnitude of the $`g`$$`\gamma `$ interband transition probability with respect to the ground intraband one. Thus it naturally characterizes the link between the two bands within the multiplet. The second ratio represents one of the widely used collective characteristics of nuclei related to Alaga rules. Both quantities (29) and (30) can be obtained from the experimental data on deformed nuclei and therefore have a direct physical meaning.
In order to derive analytic expressions for the above ratios we calculate the matrix elements of the quadrupole operator $`Q_0`$ between the eigenstates
$`|\omega _i^{(2)}`$ $`=`$ $`C_{i1}^{(2)}|\begin{array}{c}(\lambda ,\mu )\\ \mu /21,2,2\end{array}+C_{i2}^{(2)}|\begin{array}{c}(\lambda ,\mu )\\ \mu /2,2,2\end{array},i=1,2;`$ (35)
$`|\omega ^{(0)}`$ $`=`$ $`C^{(0)}|\begin{array}{c}(\lambda ,\mu )\\ \mu /2,0,0\end{array},`$ (38)
of the VBM Hamiltonian (1). (It should be remembered that the eigenvalues $`\omega _1^{(2)}`$, $`\omega _2^{(2)}`$ and $`\omega ^{(0)}`$ correspond to the $`2_g`$, $`2_\gamma `$ and $`0_g`$ states respectively.) After applying analytically the formalism developed in we obtain the following matrix elements:
$`\omega _1^{(2)}|Q_0|\omega _2^{(2)}`$ $`=`$ $`{\displaystyle \frac{4}{7}}{\displaystyle \frac{\lambda (C_{21}^{(2)})^2+\mu (C_{22}^{(2)})^2+(2\lambda +2\mu +3)C_{21}^{(2)}C_{22}^{(2)}}{C_{11}^{(2)}C_{22}^{(2)}C_{21}^{(2)}C_{12}^{(2)}}};`$ (39)
$`\omega _1^{(0)}|Q_0|\omega _1^{(2)}`$ $`=`$ $`\sqrt{6}C^{(0)}{\displaystyle \frac{\mu C_{22}^{(2)}\lambda C_{21}^{(2)}}{C_{11}^{(2)}C_{22}^{(2)}C_{12}^{(2)}C_{21}^{(2)}}};`$ (40)
$`\omega ^{(0)}|Q_0|\omega _2^{(2)}`$ $`=`$ $`\sqrt{6}C^{(0)}{\displaystyle \frac{\lambda C_{11}^{(2)}\mu C_{12}^{(2)}}{C_{11}^{(2)}C_{22}^{(2)}C_{12}^{(2)}C_{21}^{(2)}}}.`$ (41)
The wave-function coefficients are determined as
$`C_{i1}^{(2)}`$ $`=`$ $`(f_{11}^{(2)}+2h_{i2}f_{21}^{(2)}+h_{i2}^2f_{22}^{(2)})^{\frac{1}{2}};`$ (42)
$`C_{i2}^{(2)}`$ $`=`$ $`h_{i2}C_{i1}^{(2)},i=1,2`$ (43)
$`C^{(0)}`$ $`=`$ $`(f^{(0)})^{\frac{1}{2}}.`$ (44)
Here
$`f_{11}^{(2)}=\begin{array}{c}(\lambda ,\mu )\\ \mu /21,2\end{array}|\begin{array}{c}(\lambda ,\mu )\\ \mu /21,2\end{array}={\displaystyle \frac{1}{30}}R(\lambda ,\mu ){\displaystyle \underset{l=0}{\overset{\mu 2}{}}}\left(\begin{array}{c}\mu /2\\ l/2\end{array}\right)S^l(\lambda ,\mu ){\displaystyle \frac{(l+1)(\mu l)}{\mu ^2(\lambda +l+6)}}`$ (51)
$`\times [(\mu l)(\lambda +2)(\lambda +3)(\lambda +5)\mu \lambda (\lambda +4)(\lambda +l+6)];`$ (52)
$`f_{21}^{(2)}`$ $`=`$ $`\begin{array}{c}(\lambda ,\mu )\\ \mu /2,2\end{array}|\begin{array}{c}(\lambda ,\mu )\\ \mu /21,2\end{array}={\displaystyle \frac{1}{15}}R(\lambda ,\mu ){\displaystyle \underset{l=0}{\overset{\mu 2}{}}}\left(\begin{array}{c}\mu /2\\ l/2\end{array}\right)S^l(\lambda ,\mu ){\displaystyle \frac{(l+1)(\mu l)}{\mu }};`$ (59)
$`f_{22}^{(2)}`$ $`=`$ $`\begin{array}{c}(\lambda ,\mu )\\ \mu /2,2\end{array}|\begin{array}{c}(\lambda ,\mu )\\ \mu /2,2\end{array}={\displaystyle \frac{1}{15}}R(\lambda ,\mu ){\displaystyle \underset{l=0}{\overset{\mu }{}}}\left(\begin{array}{c}\mu /2\\ l/2\end{array}\right)S^l(\lambda ,\mu )(l+1)(l+2);`$ (66)
$`f^{(0)}`$ $`=`$ $`\begin{array}{c}(\lambda ,\mu )\\ \mu /2,0\end{array}|\begin{array}{c}(\lambda ,\mu )\\ \mu /2,0\end{array}=R(\lambda ,\mu ){\displaystyle \underset{l=0}{\overset{\mu }{}}}\left(\begin{array}{c}\mu /2\\ l/2\end{array}\right)S^l(\lambda ,\mu ){\displaystyle \frac{(\lambda +\mu l)(\lambda +l+4)}{\mu (\lambda +3)(\lambda +\mu +4)}},`$ (73)
are the corresponding overlap integrals obtained by the general expression in , with
$`R(\lambda ,\mu )`$ $`=`$ $`(\lambda +3)!!(\mu !!)^2,`$ (74)
$`S^l(\lambda ,\mu )`$ $`=`$ $`((l1)!!)^2{\displaystyle \frac{(\lambda +\mu l2)!!(\lambda +\mu +4)!!}{(\lambda +l+4)!!}}.`$ (75)
In addition \[see Eqs. (13)–(15) in ref. \]
$`h_{i2}={\displaystyle \frac{(V_{11}\omega _i^{(2)})}{V_{12}}}=[3g_2(2\lambda +2\mu +3)+g_3[(\lambda +\mu )^2+2\lambda +\mu ]+`$
$`(1)^i\sqrt{A(\lambda ,\mu )g_2^2+B(\lambda ,\mu )g_3^2C(\lambda ,\mu )g_2g_3}]/(6g_2\mu g_3\mu (\mu 2)),`$ (76)
with $`A(\lambda ,\mu )`$, $`B(\lambda ,\mu )`$ and $`C(\lambda ,\mu )`$ being defined in Eqs. (17)–(19).
By using the general expression for the B(E2) transition probability between two of the above eigenstates
$`B(E2;L_\nu L_{\nu ^{\prime \prime }}^{})`$ $`=`$ $`\left(\begin{array}{ccc}L^{}& 2& L\\ L& 0& L\end{array}\right)^2\left|\omega _{\nu ^{\prime \prime }}^{(L^{})}|Q_0|\omega _\nu ^{(L)}\right|^2`$ (79)
($`L,L^{}=0,2`$; $`\nu ,\nu ^{\prime \prime }=g,\gamma `$), we have studied analytically the transition ratios (29) and (30) in the two limiting cases considered in the previous subsection.
We have analyzed the explicit expressions for the overlap integrals (52)–(73) and the $`h_{i2}`$-factors (76). In this way we have deduced that in both limits, (i) and (ii), the overlap integrals increase to infinity. On the other hand one can verify that this behavior is compensated consistently in the ratios (29) and (30), where the total contribution of the integrals and the $`h_{i2}`$-factors is finite.
Thus, for the case (i) ($`\lambda \mathrm{}`$, with $`\mu `$ finite) we have obtained the following analytic limits of the transition ratios $`R_1(2)`$ and $`R_2(2)`$:
$`\underset{\lambda \mathrm{}}{lim}{\displaystyle \frac{B(E2;2_\gamma 2_g)}{B(E2;2_g0_g)}}`$ $`=`$ $`0;`$ (80)
$`\underset{\lambda \mathrm{}}{lim}{\displaystyle \frac{B(E2;2_\gamma 2_g)}{B(E2;2_\gamma 0_g)}}`$ $`=`$ $`{\displaystyle \frac{10}{7}}\left({\displaystyle \frac{\mu +2}{2\mu }}\right)^2.`$ (81)
So, in this case we find that the relative magnitude of the $`g`$$`\gamma `$ interband transition is zero, while the ratio $`R_2(2)`$ obtains finite values depending on the quantum number $`\mu `$. We remark that for $`\mu =2`$ one has $`R_2(2)=10/7`$, which is the standard Alaga value.
In case (ii) ($`\mu =\lambda \mathrm{}`$) we obtain the following limits:
$`\underset{\lambda =\mu \mathrm{}}{lim}{\displaystyle \frac{B(E2;2_\gamma 2_g)}{B(E2;2_g0_g)}}`$ $`=`$ $`{\displaystyle \frac{10}{7}}{\displaystyle \frac{(c_2^2+4c_2+1)^2}{(c_2^2+c_2+1)(c_21)^2}}0.172;`$ (82)
$`\underset{\lambda =\mu \mathrm{}}{lim}{\displaystyle \frac{B(E2;2_\gamma 2_g)}{B(E2;2_\gamma 0_g)}}`$ $`=`$ $`{\displaystyle \frac{10}{7}}{\displaystyle \frac{(c_2^2+4c_2+1)^2(c_1^2+c_1+1)}{(c_2^2+c_2+1)^2(c_11)^2}}0.304,`$ (83)
with $`c_1=4\sqrt{13}7.606`$ and $`c_2=4+\sqrt{13}0.394`$.
In this case one finds that both ratios, $`R_1(2)`$ and $`R_2(2)`$, remain finite.
We remark that all obtained limits do not depend on the model parameters. (It is assumed that $`g_1`$, $`g_2`$, and $`g_3`$ are finite, with $`g_2<0`$ and $`g_3<0`$.)
## 4 The ($`\lambda ,2`$)-direction
### 4.1 SU(3) splitting in $`L2`$ states
For the $`(\lambda ,2)`$\- irreps the $`g`$\- and the $`\gamma `$-bands are the only possible ones appearing in the corresponding SU(3) multiplets. They are labeled by the quantum numbers $`\alpha _2=1`$ and $`\alpha _1=0`$ respectively \[See inequality (3)\]. In the even angular momentum states the Hamiltonian matrix is always two-dimensional, while for the odd states of the $`\gamma `$\- band one has a single matrix element. Hence for the $`(\lambda ,2)`$\- multiplets one is able to derive analytic expressions for the spectroscopic characteristics of the full set of states ($`L2`$) in a way similar to that of the previous section. That is why we do not explain in detail all steps of analytic calculations and report only the final results in this direction.
So, for a given $`(\lambda ,2)`$ multiplet the energy levels $`E^g(L)`$ and $`E^\gamma (L)`$ of the $`g`$ and the $`\gamma `$-band can be written in the following form:
$`E^g(L)`$ $`=`$ $`\stackrel{~}{B}+\stackrel{~}{A}L(L+1)|\stackrel{~}{B}|R^{(L)},`$ (84)
$`E^\gamma (L_{even})`$ $`=`$ $`\stackrel{~}{B}+\stackrel{~}{A}L(L+1)+|\stackrel{~}{B}|R^{(L)};`$ (85)
$`E^\gamma (L_{odd})`$ $`=`$ $`2\stackrel{~}{B}+(\stackrel{~}{A}+g_3)L(L+1),`$ (86)
where
$`\stackrel{~}{A}`$ $`=`$ $`\stackrel{~}{A}(g_1,g_2,g_3)=g_1(2\lambda +5)g_2g_3,`$ (87)
$`\stackrel{~}{B}`$ $`=`$ $`\stackrel{~}{B}(\lambda ,g_2,g_3)=6(2\lambda +5)g_22(\lambda +3)^2g_3,`$ (88)
and
$$R^{(L)}=\sqrt{1+aL(L+1)+bL^2(L+1)^2}.$$
(89)
with
$`a=a(\lambda ,g_2,g_3)`$ $`=`$ $`{\displaystyle \frac{4}{\stackrel{~}{B}^2}}\left\{(\lambda +3)[(\lambda +3)g_36g_2]g_33(g_36g_2)g_2\right\},`$ (90)
$`b=b(\lambda ,g_2,g_3)`$ $`=`$ $`{\displaystyle \frac{1}{\stackrel{~}{B}^2}}\left(g_36g_2\right)^2.`$ (91)
Now we introduce the following energy ratio
$$\mathrm{\Delta }E_L=\frac{E_L^\gamma E_L^g}{E_2^g},$$
(92)
which is more general compared to Eq.(21) and characterizes the magnitude of the energy splitting in any even angular momentum state of a given SU(3) multiplet.
By using Eqs.(84) and (85) we obtain $`\mathrm{\Delta }E_L`$ in the following analytic form
$$\mathrm{\Delta }E_L=\frac{2|\stackrel{~}{B}|R^{(L)}}{6\stackrel{~}{A}|\stackrel{~}{B}|R^{(2)}+\stackrel{~}{B}},$$
(93)
which in the SU(3) contraction limit goes to
$$\underset{\stackrel{\lambda \mathrm{}}{\mu =2}}{lim}\mathrm{\Delta }E_L=\frac{2}{g_3/|g_3|1}.$$
(94)
For $`g_3<0`$ one has
$$\underset{\stackrel{\lambda \mathrm{}}{\mu =2}}{lim}\mathrm{\Delta }E_L=\mathrm{}.$$
(95)
Thus we find that for all even states of a given $`(\lambda ,2)`$\- multiplet the SU(3) splitting goes to infinity in the same way \[see also Eqs.(24) and (25)\].
### 4.2 Transition ratios in the ($`\lambda ,2`$)-direction
For the ($`\lambda ,2`$)-direction the B(E2) transitions between the states of a given multiplet can be examined through the following \[more general compared to (29) and (30)\] transition ratios:
$`R_1(L)`$ $`=`$ $`{\displaystyle \frac{B(E2;L_\gamma L_g)}{B(E2;L_g(L2)_g)}},L=even,`$ (96)
$`R_2(L)`$ $`=`$ $`{\displaystyle \frac{B(E2;L_\gamma L_g)}{B(E2;L_\gamma (L2)_g)}},L=even,`$ (97)
$`R_3(L)`$ $`=`$ $`{\displaystyle \frac{B(E2;L_\gamma (L+1)_g)}{B(E2;L_\gamma (L1)_g)}},L=odd.`$ (98)
The first two ratios, $`R_1(L)`$ and $`R_2(L)`$, have the same physical meaning as the ratios (29) and (30) of the previous section. The third ratio, $`R_3(L)`$, involves the odd angular momentum states in the study. In such a way we investigate the transition characteristics of the full set of states in a given SU(3) multiplet.
In the case of a ($`\lambda ,2`$)-multiplet the Hamiltonian eigenstates are constructed as
$`|\omega _i^{(L)}`$ $`=`$ $`C_{i1}^{(L)}|\begin{array}{c}(\lambda ,2)\\ 0,L\end{array}+C_{i2}^{(L)}|\begin{array}{c}(\lambda ,2)\\ 1,L\end{array},i=1,2,L=even;`$ (103)
$`|\omega _{odd}^{(L)}`$ $`=`$ $`C_{odd}^{(L)}|\begin{array}{c}(\lambda ,2)\\ 0,L\end{array},L=odd.`$ (106)
The necessary transition matrix elements are derived in the form
$`\omega _1^{(L)}|Q_0|\omega _2^{(L)}`$ $`=`$ $`{\displaystyle \frac{12}{(L+1)(2L+3)}}[(\lambda +2L)(C_{21}^{(L)})^2`$ (107)
$`+`$ $`[2\lambda +5+L(L1)]C_{21}^{(L)}C_{22}^{(L)}`$
$`+`$ $`L(L1)(C_{22}^{(L)})^2]/[C_{11}^{(L)}C_{22}^{(L)}C_{21}^{(L)}C_{12}^{(L)}];`$
$`\omega _1^{(L2)}|Q_0|\omega _1^{(L)}`$ $`=`$ $`{\displaystyle \frac{6}{\sqrt{L(2L1)}}}[(\lambda +4L)C_{11}^{(L2)}C_{22}^{(L)}`$ (108)
$``$ $`(\lambda +2L)C_{12}^{(L2)}C_{21}^{(L)}`$
$`+`$ $`2C_{12}^{(L2)}C_{22}^{(L)}]/[C_{11}^{(L)}C_{22}^{(L)}C_{21}^{(L)}C_{12}^{(L)}];`$
$`\omega _1^{(L2)}|Q_0|\omega _2^{(L)}`$ $`=`$ $`{\displaystyle \frac{6}{\sqrt{L(2L1)}}}[(\lambda +4L)C_{11}^{(L2)}C_{12}^{(L)}`$ (109)
$``$ $`(\lambda +2L)C_{12}^{(L2)}C_{11}^{(L)}`$
$`+`$ $`2C_{12}^{(L2)}C_{12}^{(L)}]/[C_{11}^{(L)}C_{22}^{(L)}C_{21}^{(L)}C_{12}^{(L)}];`$
$`\omega _1^{(L+1)}|Q_0|\omega _2^{(L)}`$ $`=`$ $`{\displaystyle \frac{6}{(L+2)\sqrt{L+1}}}C_{\mathrm{odd}}^{(L)}[(2\lambda L+4)C_{22}^{(L+1)}`$ (110)
$`+`$ $`(\lambda L+1)C_{21}^{(L+1)}]/[C_{11}^{(L+1)}C_{22}^{(L+1)}C_{21}^{(L+1)}C_{12}^{(L+1)}];`$
$`\omega _1^{(L1)}|Q_0|\omega _2^{(L)}`$ $`=`$ $`{\displaystyle \frac{12}{(L+1)\sqrt{L}}}{\displaystyle \frac{1}{C_{\mathrm{odd}}^{(L)}}}\left[(\lambda L+3)C_{11}^{(L1)}(L1)C_{12}^{(L1)}\right],`$ (111)
with the wave-function coefficients
$`C_{i1}^{(L)}`$ $`=`$ $`\left(f_{11}^{(L)}+2h_{i2}^{(L)}f_{21}^{(L)}+(h^{(L)})_{i2}^2f_{22}^{(L)}\right)^{\frac{1}{2}},L=even;`$ (112)
$`C_{i2}^{(L)}`$ $`=`$ $`h_{i2}^{(L)}C_{i1}^{(L)},L=even,i=1,2`$ (113)
$`C_{odd}^{(L)}`$ $`=`$ $`(f_{odd}^{(L)})^{\frac{1}{2}},L=odd.`$ (114)
The corresponding overlap integrals are obtained in the form
$`f_{11}^{(L)}`$ $`=`$ $`\begin{array}{c}(\lambda ,2)\\ 0,L\end{array}|\begin{array}{c}(\lambda ,2)\\ 0,L\end{array}=S^L(\lambda ){\displaystyle \frac{[(L^2+L+1)\lambda +L^3+4L^2+2L+2]}{L(L1)}};`$ (119)
$`f_{21}^{(L)}`$ $`=`$ $`\begin{array}{c}(\lambda ,2)\\ 1,L\end{array}|\begin{array}{c}(\lambda ,2)\\ 0,L\end{array}=S^L(\lambda )(\lambda +L+4);`$ (124)
$`f_{22}^{(L)}`$ $`=`$ $`\begin{array}{c}(\lambda ,2)\\ 1,L\end{array}|\begin{array}{c}(\lambda ,2)\\ 1,L\end{array}=S^L(\lambda ){\displaystyle \frac{[2(\lambda L+2)(\lambda +L+4)+(L+1)(L+2)]}{(\lambda L+2)}};`$ (129)
$`f_{\mathrm{odd}}^{(L)}`$ $`=`$ $`\begin{array}{c}(\lambda ,2)\\ 0,L\end{array}|\begin{array}{c}(\lambda ,2)\\ 0,L\end{array}=S_{\mathrm{odd}}^L(\lambda ){\displaystyle \frac{(L+1)(L+2)(\lambda +2)}{2L(L1)}},`$ (134)
with
$`S^L(\lambda )`$ $`=`$ $`{\displaystyle \frac{2L!(\lambda L+2)!!(\lambda +L+1)!!}{(2L+1)!!}}`$
$`S_{\mathrm{odd}}^L(\lambda )`$ $`=`$ $`{\displaystyle \frac{2L!(\lambda L+1)!!(\lambda +L+2)!!}{(2L+1)!!}}.`$ (135)
Also we have
$`h_{i2}^{(L)}`$ $`=`$ $`[6g_2[(2\lambda +5)+L(L1)]+g_3[2(\lambda +3)^2L(L+1)]`$ (136)
$`+`$ $`(1)^i\stackrel{~}{B}R^{(L)}]/[12g_2L(L1)].`$
Eqs. (119)–(134) show that the overlap integrals increase to infinity with the increase of the quantum number $`\lambda `$. However, similarly to the previous section, one can verify that this behavior is compensated consistently in the ratios (96)–(98).
After using the above analytic form of the matrix elements (107)–(111) we have obtained the SU(3) contraction limits of the ratios $`R_1(L)`$, $`R_2(L)`$ and $`R_3(L)`$ \[Eqs. (96)–(98)\]:
$`\underset{\stackrel{\lambda \mathrm{}}{\mu =2}}{lim}{\displaystyle \frac{B(E2;L_\gamma L_g)}{B(E2;L_g(L2)_g)}}`$ $`=`$ $`0;`$ (137)
$`\underset{\stackrel{\lambda \mathrm{}}{\mu =2}}{lim}{\displaystyle \frac{B(E2;L_\gamma L_g)}{B(E2;L_\gamma (L2)_g)}}`$ $`=`$ $`6{\displaystyle \frac{(L1)(2L+1)}{(L+1)(2L+3)}};`$ (138)
$`\underset{\stackrel{\lambda \mathrm{}}{\mu =2}}{lim}{\displaystyle \frac{B(E2;L_\gamma (L+1)_g)}{B(E2;L_\gamma (L1)_g)}}`$ $`=`$ $`{\displaystyle \frac{(L1)}{(L+2)}}.`$ (139)
Thus, we find that for all the states ($`L2`$) of the Hamiltonian the relative magnitude of the $`g`$$`\gamma `$ interband transitions goes to zero. Also, we see that the ratios $`R_2(L)`$ and $`R_3(L)`$ go to the corresponding standard Alaga rules.
## 5 Results and Discussions
The theoretical results given above allow one to examine the mechanism of the SU(3) symmetry reduction in the space of the $`(\lambda ,\mu )`$ irreps as well as to identify its manifestation in reference to the experimental data on heavy deformed nuclei.
The analytic study of the Hamiltonian matrix elements shows \[Eqs. (23) and (26)\] how the increase in the quantum numbers $`\lambda `$ and/or $`\mu `$ is connected with the corresponding decrease in the $`g`$$`\gamma `$ band-mixing interaction within the framework of the SU(3) symmetry. Generally this result illustrates the behavior of the energy-mixing in the $`(\lambda ,\mu )`$-plane. In both limits, (i) and (ii), the $`g`$$`\gamma `$ mixing decreases asymptotically to zero. Similar limiting behavior of the $`L2`$ matrix elements in $`(\lambda ,2)`$-direction has been established in our previous work (See Sec. IV–C of Ref. ). Thus in all limiting cases the SU(3) symmetry disappears completely and the two bands do not belong to the same SU(3) multiplet anymore.
It is appropriate at this point to elucidate the meaning of the above consideration in terms of the SU(3) group contraction process . This process corresponds to a renormalization of the quadrupole operator, $`QQ/\sqrt{C_2}`$, with
$$C_2=(\lambda +2\mu )(\lambda +2\mu +6)+3\lambda (\lambda +2)$$
(140)
being the eigenvalue of the second order Casimir operator of SU(3). The following commutation relations between the angular momentum and the renormalized quadrupole operators are then valid:
$`[L_m,L_n]`$ $`=`$ $`\sqrt{2}C_{1m1n}^{1m+n}L_{m+n},`$ (141)
$`[L_m,Q_n]`$ $`=`$ $`\sqrt{6}C_{1m2n}^{2m+n}Q_{m+n},`$ (142)
$`[Q_m,Q_n]`$ $`=`$ $`3\sqrt{10}C_{2m2n}^{1m+n}{\displaystyle \frac{L_{m+n}}{C_2}}.`$ (143)
They differ from the standard SU(3) commutation relations by the factor $`C_2`$ in the right-hand side of (143). Taking into account Eq. (140), one finds that in both limits (i) and (ii), considered in the present work, the commutator (143) vanishes and the commutation relations of the algebra of the triaxial rotor group T$`{}_{5}{}^{}`$SO(3) hold. In such a way the vanishing $`g`$$`\gamma `$ band-mixing could be interpreted as a transition from a compact to a non-compact DS group.
Let us now analyze the behavior of the splitting and transition ratios of Secs. 3 and 4 in the $`(\lambda ,\mu )`$-plane. For this purpose we use the analytic expressions for numerical calculations. In the particular case of $`L2`$ states in ($`\mu =\lambda \mathrm{}`$) direction, which is not accessible analytically, we apply numerically the algorithm developed in Ref. . All calculations are carried out for the same set of fixed model parameters $`g_1=1`$, $`g_2=0.2`$, $`g_3=0.25`$. These values belong to the corresponding parameter regions obtained for a group of rare earth nuclei and actinides \[See Table 2 in Ref. \]. In this respect they can be considered as an overall set of model parameters. Also, it should be emphasized that in the SU(3) contraction limiting cases the various sets of (finite) parameter values give the same asymptotic behavior of the model quantities.
In Fig. 1 the splitting ratio $`\mathrm{\Delta }E_2`$ \[Eq. (22)\] is plotted as a function of the quantum numbers $`\lambda `$ and $`\mu `$. In the limiting case (i) ($`\lambda \mathrm{}`$, with $`\mu `$ finite) the two-dimensional surface shows a rapid increase of $`\mathrm{\Delta }E_2`$, while in case (ii) ($`\mu =\lambda \mathrm{}`$) the splitting ratio gradually saturates towards the constant value $`5.17`$ \[See Eq. (28)\]. In Fig. 2 the splitting ratio $`\mathrm{\Delta }E_L`$ \[Eq. (92)\] is plotted as a function of the quantum number $`\lambda `$ for $`L=2,4\mathrm{}12`$. In the case $`\lambda \mathrm{}`$, $`\mu =2`$ the energy splitting goes to infinity with almost equal values for all angular momenta (Fig. 2(a)). In case (ii) $`\mathrm{\Delta }E_L`$ trends to finite values which increase with $`L`$ (Fig. 2(b)). So, in the first limiting case the complete reduction of the SU(3) symmetry leads to a large energy separation between the bands in the multiplet, while in the second case (for finite angular momenta) the bands remain close to each other, but their mutual disposition does not depend on the Hamiltonian parameters anymore, so that it should not be associated with any band coupling.
The experimental $`\mathrm{\Delta }E_2`$ ratios of several rare earth nuclei and actinides are given in Table 1. They vary within the limits $`5\mathrm{\Delta }E_220`$, for the rare earths and $`13\mathrm{\Delta }E_225`$, for the actinides. The behavior of the splitting ratios is clear: The $`\mathrm{\Delta }E_2`$ ratio generally increases towards the middle of the rotational region. This is illustrated in Table 1 through the number of the nucleon pairs (or holes) in the valence shells, $`N`$. (The number $`N`$ is a well established characteristic of nuclear collectivity used in the IBM .) A clearly pronounced increase of $`\mathrm{\Delta }E_2`$ with increasing $`N`$ is observed for the isotopes of Sm, Gd, Er, Yb, and W. Similar behavior of the energy splitting is observed in the $`L>2`$ states of these nuclei . One concludes that the data show that the SU(3) splitting increases toward the midshell regions.
We turn now to the analysis of the interband transition ratios. In Fig. 3 the theoretical ratio $`R_1(2)`$, Eq. (29), is plotted as a two-dimensional function of the quantum numbers $`\lambda `$ and $`\mu `$. In the limiting case (i) ($`\lambda \mathrm{}`$, with $`\mu `$ finite) the $`R_1`$ surface shows a rapid decrease to zero. In case (ii) ($`\mu =\lambda \mathrm{}`$) $`R_1(2)`$ decreases gradually and saturates towards the constant value $`0.172`$ \[See Eq. (82)\]. In Fig. 4 the transition ratio $`R_1(L)`$ \[Eq. (96)\] is plotted as a function of the quantum number $`\lambda `$ for $`L=2,4\mathrm{}12`$. In the case $`\lambda \mathrm{}`$, $`\mu =2`$ it goes to zero with almost equal values for all angular momenta (Fig. 4(a)). In case (ii) $`R_1(L)`$ obtains finite values which decrease with $`L`$ (Fig. 4(b)). Thus in the first limiting case the $`g`$$`\gamma `$ interband transition link vanishes rapidly, while in the second case (for finite angular momenta) the relative magnitude of the interband transition probability remains non-zero. However, as in the energy splitting, the $`R_1`$-ratios do not depend on the Hamiltonian parameters anymore. Therefore, they should not be treated in terms of the SU(3) symmetry anymore.
The experimental $`R_1(2)`$-values for several rare earth nuclei and actinides are given in Table 1. Here one observes a rather spectacular decrease of $`R_1(2)`$ towards the midshell regions. The best examples (with the largest number of available data) occur in the cases of the Gd, Er and Yb isotopes. Note that the decrease in the experimental $`g`$$`\gamma `$ transition probabilities is well consistent with the corresponding increase in the SU(3) splitting. In this way, the experimental data strongly support the VBM predictions in the SU(3) contraction limit.
The two-dimensional surface obtained for the theoretical $`R_2(2)`$ ratio, Eq. (30), is shown in Fig. 5. We see that $`R_2(2)`$ gradually decreases with $`\lambda `$ and $`\mu `$ and trends towards the finite values, obtained in the limiting cases (i) and (ii) \[See Eqs. (81) and (83)\]. In Fig. 6 the transition ratio $`R_2(L)`$ \[Eq. (97)\] is plotted as a function of the quantum number $`\lambda `$ for $`L=2,4\mathrm{}12`$. In the case $`\lambda \mathrm{}`$, $`\mu =2`$ it gradually goes to the Alaga values for the corresponding angular momenta (Fig. 6(a), see also Eq. (138)). In case (ii) $`R_2(L)`$ trends to finite values which (in the numerically investigated $`\lambda `$\- range) exhibit a complicated behavior as a function of $`L`$ (Fig. 6(b)). In both limiting cases the lack of dependence on the Hamiltonian parameters indicates the complete reduction of the SU(3) symmetry. The experimental data on $`R_2(2)`$, given in Table 1, show a slightly expressed trend of decreasing towards the midshells, but one could not draw any definite conclusions about the systematic behavior of this quantity.
In Fig. 7 the transition ratio $`R_3(L)`$ \[Eq. (98)\] is plotted as a function of the quantum number $`\lambda `$ for the odd angular momenta $`L=3,5\mathrm{}11`$. For $`\lambda \mathrm{}`$, $`\mu =2`$ it goes to the corresponding Alaga values in a way similar to the $`R_2(L)`$\- ratio (Fig. 7(a), see also Eq. (139)). In the second direction, (ii), $`R_3(L)`$ saturates to finite values (Fig. 7(b)). It is clear that towards the SU(3) contraction limit the B(E2) transition characteristics of the odd $`\gamma `$ band states should be consistent with the even angular momentum ones.
In order to assess quantitatively the results presented so far, we provide a numerical analysis of the SU(3) multiplets in $`(\lambda ,\mu )`$\- plane on the basis of the experimental energy and transition ratios given in Table 1. More precisely we determined the quantum numbers $`\lambda `$ and $`\mu `$ together with the Hamiltonian parameters $`g_1`$ $`g_2`$ and $`g_3`$, by fitting them in the numerical procedure of Ref. so as to reproduce the experimental $`g`$\- and $`\gamma `$\- band levels up to $`L=8`$, for lantanides and $`L=18`$, for actinides together with the values of the experimental ratios $`\mathrm{\Delta }E_2`$, $`R_1(2)`$ and $`R_2(2)`$. (Only for two nuclei, <sup>160</sup>Gd and <sup>162</sup>Dy, the $`R_1(2)`$ ratios have not been used in the fits due to the uncertainty of the experimental data.)
The “favored” values of the quantum numbers $`\lambda `$ and $`\mu `$ obtained for the various isotopes are given in the fourth column of Table 1. Generally the quantum number $`\lambda `$ vary in the range $`14\lambda 68`$, while $`\mu `$ obtains the values $`2\mu 6`$.
So, one finds that while $`\mu `$ is closed in narrow limits, the quantum number $`\lambda `$ exhibits a well pronounced systematic behavior. The favored $`\lambda `$-values as well as the SU(3) splitting ratio $`\mathrm{\Delta }E_2`$ increase with the increase of the valence pair number $`N`$, i.e. towards the middle of the valence shells in rotational nuclei. For example, the small splitting observed in the nuclei <sup>152</sup>Sm ($`\mathrm{\Delta }E_2=7.9`$), <sup>154</sup>Gd ($`\mathrm{\Delta }E_2=7.1`$) and <sup>162</sup>Er ($`\mathrm{\Delta }E_2=7.8`$) which are situated near the beginning of the respective group of rotational isotopes, is associated with the small $`\lambda `$\- values, $`\lambda =1416`$ and the relatively large interband transition ratios $`R_1(2)=0.070.08`$. On the other hand, for the middshell nuclei <sup>172,174</sup>Yb with large splitting values, $`\mathrm{\Delta }E_2=1820`$ and small $`R_1(2)0.01`$, we obtain large $`\lambda `$-values, $`\lambda 6070`$. Also, large $`\lambda `$ values, $`\lambda 60`$ have been obtained for the <sup>234,238</sup>U isotopes with $`\mathrm{\Delta }E_2=2022`$. Well pronounced systematic behavior of the quantum number $`\lambda `$ is observed in the Er and Yb isotopes.
The above results are consistent with the numerical analyses carried out in . It should be mentioned that the involvement of the interband transition ratio $`R_1(2)`$ in the present fits leads to an increase in the quantum number $`\mu `$ above the $`\mu =2`$ values. Actually, some trend of small increase in $`\mu `$ (up to $`\mu =6`$) with the increase of $`\lambda `$ is indicated for various isotope groups (See, for example the Yb isotopes in Table 1). Nevertheless, in almost all nuclei under study the quantum number $`\lambda `$ is essentially larger than $`\mu `$, which is natural for the well deformed nuclei. Several exceptions are observed in the nuclei far from the midshell region, such as <sup>158</sup>Dy (with $`(\lambda ,\mu )=(16,6)`$) and <sup>186</sup>W (with $`(\lambda ,\mu )=(24,10)`$) where the interband transition ratios are very large $`R_1(2)0.10.2`$.
The above quantitative considerations show that the changes in the SU(3) characteristics of nuclei (especially the quantum number $`\lambda `$) towards the middle of given rotational region could be associated with the corresponding decrease in the $`g`$$`\gamma `$ band mixing interaction towards the SU(3) contraction limits. In terms of our study the strong $`g`$$`\gamma `$ splitting, observed near the middle of rotational regions, corresponds to the weak mutual perturbation of the bands. This is consistent with the respectively good rotational behavior of the $`g`$-band, which in this case could belong to a separate SU(3) multiplet. \[See the experimental energy ratios, $`R_4=E_4^g/E_2^g`$, given in Table 1.\]
We remark that present analyses are based mainly on data in the rare earth region. Actually Table 1 includes only nuclei for which the $`g`$$`\gamma `$ interband transition probabilities are measured. That is why only four actinides (<sup>230,232</sup>Th, <sup>234,238</sup>U) are considered. Nevertheless, they give an indication for similar behaviors of the splitting and the interband transition ratios as the ones in rare earth nuclei. On the other hand, the generally stronger energy splitting, observed in the actinide region (See also ), suggests a generally weaker $`g`$$`\gamma `$ coupling compared to the rare earth nuclei.
We are now able to discuss the physical significance of the considered SU(3) contraction limits as well as to depict the physically meaningful directions in the $`(\lambda ,\mu )`$\- plane, which could be appropriate for studying the transition between the different band coupling schemes. The theoretical analyses and experimental data show that the limiting case (i) ($`\lambda \mathrm{}`$, with $`\mu `$ finite) has a rather clear physical interpretation. It is consistent with the observed continuous increase of the $`g`$$`\gamma `$ band splitting and the corresponding continuous decrease of the interband transition probabilities towards the midshell regions in rare earth nuclei. The limiting case (ii) ($`\lambda =\mu \mathrm{}`$) does not have any similar direct interpretation. It suggests finite values for the splitting and the interband transition probabilities, while the bands do not interact in the framework of SU(3) symmetry. In addition, it is well known that the case $`\lambda =\mu `$ does not correspond to deformed nuclei, for which the inequality $`\mu <\lambda `$ is satisfied. Nevertheless the study of this limit is useful from the following viewpoint: It implies that the strong suppression of the band interaction as well as the transition between the different band coupling schemes could be realized at reasonable (finite) SU(3) splitting. Based on the above considerations, we deduce that the possibly interesting physically meaningful directions in the $`(\lambda ,\mu )`$–plane should be associated with a consistent increase in the quantum numbers $`\lambda `$ and $`\mu `$. Thus, any particular direction of interest could be easily estimated by using its intermediate behavior between the two considered limiting cases.
Some discussion concerning the Interacting Boson Model (IBM) classification scheme is appropriate at this point. In Ref. it has been suggested that for deformed nuclei both the VBM scheme (with $`g`$$`\gamma `$ band coupling) and the IBM one (with $`\beta `$$`\gamma `$ band coupling) could be considered as complementary schemes. It has also been pointed out that the SU(3) scheme of the VBM is naturally applicable to nuclei with weak energy splitting, while strong splitting invokes the SU(3) scheme of the IBM, in which the $`g`$\- band is situated in a separate irrep. Furthermore, the theoretical results and the experimental data given in the present work suggest that the VBM band coupling scheme is more appropriate near the ends of the rotational regions, while in the midshell regions the coupling scheme of the IBM is realized. In this respect the detailed comparison of both band-coupling mechanisms would be of interest. For example, the analytic expressions for the $`g`$$`\gamma `$ interband transition probabilities, obtained in the framework of the IBM in , would be useful. \[See Eqs. (5) and (6) of .\] They give a behavior of the transition ratios $`R_1(L)`$, $`R_2(L)`$ and $`R_3(L)`$ \[Eqs. (96), (97) and (98)\] in the infinite valence pair number limit ($`N\mathrm{}`$) similar to the behavior obtained in the limiting case (i) ($`\lambda \mathrm{}`$, with $`\mu =2`$) \[Eqs. (137), (138) (139)\] of the present VBM scheme.
As an extension of the present studies it would be worthwhile to examine, in a similar way, the link between the $`\gamma `$\- and the $`\beta `$\- band. Furthermore, besides the VBM scheme, one could refer in this case to the modifications of the IBM in which higher-order terms conserving the SU(3) symmetry are added . The consistent study (within both models) of the ways in which the SU(3) symmetry is reduced could give important information about the rearrangement of rotational bands into different SU(3) irreps.
In the above context, we emphasize that the analyses implemented in the presented paper give a general prescription to handle the fine behavior of the band coupling interactions in any collective algebraic scheme in heavy deformed nuclei. Actually, the group contraction process should play a major role in a transition between two different band coupling schemes. The transition from the compact SU(3) group to the non-compact T$`{}_{5}{}^{}`$SO(3) rotor group could be considered as a starting point in a process of reconstruction of various multiplets in a more general symplectic group of dynamical symmetry. (It is interesting to mention that the meaning of SU(3) contraction has been also discussed (though in a rather different aspect) in reference to a possible phase transition between a superconductor and rigid rotor collective motion of nuclei .)
## 6 Conclusions
We have derived analytic expressions for the energies and $`B(E2)`$ transition probabilities in the ground- and $`\gamma `$-band states of even deformed nuclei within the Vector-Boson Model with SU(3) dynamical symmetry. On this basis we applied both analytic and numeric analyses to examine the behavior of the corresponding energy splitting and B(E2) transition ratios in the two SU(3) contraction limits of the model, (i) ($`\lambda \mathrm{}`$, with $`\mu `$ finite), and (ii) ($`\lambda =\mu \mathrm{}`$). It has been shown that in both limits the $`g`$$`\gamma `$ band mixing decreases asymptotically to zero. In case (i) this is associated with the corresponding continuous increase in the splitting of the multiplet and the rapidly vanishing $`g`$$`\gamma `$ interband transition link. Case (ii) gives finite values for the energy splitting and the interband transition ratios which, however, should not be associated with any band coupling. The latter result implies that a strong reduction of the band interaction could be possible at finite SU(3) splitting. Thus, the present analyses outline the possible directions in the $`(\lambda ,\mu )`$-plane in which the $`g`$$`\gamma `$ band coupling is reduced.
The experimental data on the ground- and $`\gamma `$\- band states in deformed even–even nuclei show clearly a pronounced increase in the $`g`$$`\gamma `$ band splitting and a corresponding decrease in the interband transition probabilities towards the midshell regions. They suggest that the SU(3) contraction effects in the $`g`$$`\gamma `$ band coupling scheme should be sought in the best rotational nuclei, in which the mutual perturbation of the bands is weak. So, the experimental data and their quantitative estimation in the VBM framework strongly support our theoretical analyses.
Based on the presented investigation, we conclude that the transition from the $`g`$$`\gamma `$ band coupling scheme to a scheme in which the $`g`$-band is situated in a separate irrep should be realized towards the midshell regions. In this respect the complementarity of the classification schemes of the Vector-Boson Model with SU(3) dynamical symmetry and the IBM becomes clear. The consistent study of the rearrangement of collective bands in deformed nuclei, including the $`\beta `$\- excited bands, will be the subject of forthcoming work.
Acknowledgments
The authors are thankful to P. Van Isacker for stimulating discussions and D. N. Kadrev for the help in collecting the experimental data. This work has been supported by the Bulgarian National Fund for Scientific Research under contract no MU–F–02/98.
Figure Captions
Figure 1. The theoretical energy splitting ratio $`\mathrm{\Delta }E_2`$ \[Eq. (22)\] is plotted as a two-dimensional function of the quantum numbers $`\lambda `$ and $`\mu `$ for $`g_1=1`$, $`g_2=0.2`$ and $`g_3=0.25`$.
Figure 2. The theoretical energy splitting ratio $`\mathrm{\Delta }E_L`$ \[Eq. (92)\] is plotted as a function of the quantum number $`\lambda `$ for $`L=2,4\mathrm{}12`$ with $`g_1=1`$, $`g_2=0.2`$ and $`g_3=0.25`$ in the cases: (a) $`\mu =2`$; (b) $`\mu =\lambda `$.
Figure 3. The theoretical $`R_1(2)`$ ratio \[Eq. (29)\] is plotted as a two-dimensional function of the quantum numbers $`\lambda `$ and $`\mu `$ for $`g_2=0.2`$ and $`g_3=0.25`$.
Figure 4. The theoretical $`R_1(L)`$ ratio \[Eq. (96)\] is plotted as a function of the quantum number $`\lambda `$ for $`L=2,4\mathrm{}12`$ with $`g_2=0.2`$ and $`g_3=0.25`$ in the cases: (a) $`\mu =2`$; (b) $`\mu =\lambda `$.
Figure 5. The same as Fig. 3 but for the theoretical $`R_2(2)`$ ratio \[Eq. (30)\].
Figure 6. The same as Fig. 4 but for the theoretical $`R_2(L)`$ ratio \[Eq. (97)\].
Figure 7. The theoretical $`R_3(L)`$ ratio \[Eq. (98)\] is plotted as a function of the quantum number $`\lambda `$ for $`L=3,5\mathrm{}11`$ with $`g_2=0.2`$ and $`g_3=0.25`$ in the cases: (a) $`\mu =2`$; (b) $`\mu =\lambda `$.
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# Synchronization in populations of globally coupled oscillators with inertial effects
## I Introduction
The dynamical behavior of large populations of nonlinearly coupled oscillators may describe many phenomena in Physics, Biology and Medicine, . In particular synchronization of mean-field coupled phase oscillators with different natural frequencies is nicely illustrated by Kuramoto’s well-known and extesively analyzed model . To describe certain biological phenomena, inertial effects should be added to this model. In , Ermentrout revisited the special problem of self-synchronization in populations of fireflies of a certain kind (the Pteroptyx malaccae). Compared to observed behavior, the approach to oscillator synchronization as described by the Kuramoto model seems to be too fast. Thus a more appropriate adaptive frequency model has been proposed in , where the natural frequency of an oscillator is a new independent variable, which is allowed to vary with time. Thus an oscillator is described by its phase and frequency. From the mathematical standpoint, the new model is governed by a system of coupled second-order differential equations containing inertial terms, in contrast to the system of first-order differential equations governing the Kuramoto model. Indeed inertia slows down synchronization and this may result in better agreement between theory and experimental measurements. Other possible biological applications of Ermentrout-type models include after-effects in alterations of circadian cycles in mammalians, cf.
A different set of applications for oscillators with inertia include power systems described by the swing equations , or by Hamilton equations . An important technologically relevant application is the study of superconducting Josephson junctions arrays . Here inertial terms describe the effect of nonzero electrical capacitance. Such effect is often far from being negligible, and it is absent in the Kuramoto model, .
In this paper, we consider the model equations of Ref. ,
$`\dot{\theta _j}`$ $`=`$ $`\omega _j`$ (1)
$`m\dot{\omega _j}`$ $`=`$ $`\omega _j+\mathrm{\Omega }_j+Kr_N\mathrm{sin}(\psi _N\theta _j)+\xi _j(t),`$ (3)
$`j=1,\mathrm{},N,`$
where $`\theta _j`$, $`\omega _j`$ and $`\mathrm{\Omega }_j`$ denote phase, frequency and natural frequency of the $`j`$th oscillator, respectively. The natural frequencies are distributed with probability density $`g(\mathrm{\Omega })`$, which may have a single maximum (unimodal distribution), or several peaks (multimodal distribution). The positive parameters $`m`$ and $`K`$ are the “inertia” and the coupling strength, respectively. The complex order parameter defined by
$$r_Ne^{i\psi _N}=\frac{1}{N}\underset{j=1}{\overset{N}{}}e^{i\theta _j},$$
(4)
measures phase synchronization: $`r_N>0`$ if the oscillators are synchronized and $`r_N=0`$ if not. Finally, $`\xi _j`$’s are independent identically distributed Gaussian white noises, with $`\xi _j=0,\xi _i(t)\xi _j(s)=2D\delta _{ij}\delta (ts)`$. White noise terms were not included in . When the inertial terms vanish, $`m=0`$, Eqs. (3) and (4) are exactly the Kuramoto model.
Typically, N is very large, and oscillator synchronization is conveniently analyzed in the limiting case of infinitely many oscillators. In this limit, models with mean-field coupling are described by an evolution equation for the one-oscillator probability density, $`\rho (\theta ,\omega ,\mathrm{\Omega },t)`$, . For the present model this equation is
$`{\displaystyle \frac{\rho }{t}}`$ $`=`$ $`{\displaystyle \frac{D}{m^2}}{\displaystyle \frac{^2\rho }{\omega ^2}}`$ (5)
$``$ $`{\displaystyle \frac{1}{m}}{\displaystyle \frac{}{\omega }}[\left(\omega +\mathrm{\Omega }+Kr\mathrm{sin}(\psi \theta )\right)\rho ]\omega {\displaystyle \frac{\rho }{\theta }},`$ (6)
where the order parameter is now given by
$`re^{i\psi }={\displaystyle _0^{2\pi }}{\displaystyle _{\mathrm{}}^+\mathrm{}}{\displaystyle _{\mathrm{}}^+\mathrm{}}e^{i\theta }\rho (\theta ,\omega ,\mathrm{\Omega },t)g(\mathrm{\Omega })𝑑\mathrm{\Omega }𝑑\omega 𝑑\theta .`$ (7)
Equations (6) and (7) should be supplemented with appropriate initial and boundary data ($`\rho `$ is $`2\pi `$-periodic in $`\theta `$ and has suitable decay behavior as $`\omega \pm \mathrm{}`$) plus the normalization condition,
$$_0^{2\pi }_{\mathrm{}}^+\mathrm{}\rho (\theta ,\omega ,\mathrm{\Omega },t)𝑑\omega 𝑑\theta =1.$$
(8)
Differentiating $`_0^{2\pi }_{\mathrm{}}^+\mathrm{}\rho (\theta ,\omega ,\mathrm{\Omega },t)𝑑\omega 𝑑\theta `$ with respect to time, and then using Eq. (6) itself, together with periodicity in $`\theta `$ and decay in $`\omega `$, we find that the left side of (8) is time independent. Normalization to unity of the initial probability density then implies (8) for the solution of (6).
In this paper, we study oscillator synchronization and transition from incoherence to synchronization in the model (6) - (8). The incoherent solution of (6) - (8) (or simply incoherence) is a stationary solution which is independent of $`\theta `$. This solution asigns equal probability to all angles and has $`r=0`$ (no order), so it corresponds to lack of oscillator synchronization. There are synchronized solutions which branch off from incoherence as the coupling among oscillators is increased. These bifurcations describe the synchronization transitions, which we have analyzed and compared to the corresponding ones in the Kuramoto model. Our main results are that inertia: (i) may stabilize incoherence, making it harder to synchronize oscillators, and (ii) it may harden the synchronization transition. In the Kuramoto model ($`m=0`$) or with oscillators with identical natural frequencies, the synchronization transition is soft (supercritical bifurcation), whereas it may become hard (subcritical bifurcation) if the distribution of natural frequencies has a nonzero spread (unimodal Lorentzian distribution) or several peaks (e.g., a discrete bimodal distribution). The methods we have used in our analysis are similar to those previously employed in the Kuramoto model : linear stability of incoherence, bifurcation analysis, high-frequency singular perturbations and numerical solutions. An important difference is that now we do not have an explicit functional form for stationary solutions (as it was the case for the Kuramoto model). This has led us to use mode-coupling expansions of the solution and solving the corresponding mode-coupling equations. Solutions of these equations in close form are not always accesible, so that we have introduced some closure assumptions. The results of these uncontrolled assumptions have been compared to direct simulations or to approximate amplitude equations and found reasonable in the limit of small inertia.
The rest of the paper is as follows. In Section II, we find the incoherent solution and study its linear stability for several natural frequency distributions. Results are compared with those obtained in the massless case . It is found that the critical coupling needed to destabilize incoherence increases with $`m`$ for “unimodal” frequency distributions of the Lorentzian type. The critical coupling is independent of $`m`$ when $`g(\mathrm{\Omega })=\delta (\mathrm{\Omega })`$. In this case the time needed to reach synchronization increases as $`m`$ increases. If $`g(\mathrm{\Omega })=[\delta (\mathrm{\Omega }\mathrm{\Omega }_0)+\delta (\mathrm{\Omega }+\mathrm{\Omega }_0)]/2`$ (discrete bimodal distribution), the critical coupling may grow or decrease with $`m`$ depending on the values of $`\mathrm{\Omega }_0`$. In Section III, we construct other stationary solutions by two procedures: an amplitude expansion for solutions branching off from incoherence and a general expansion in Hermite polynomials which is appropriately truncated. An exact analytical solution is obtained if $`g(\mathrm{\Omega })=\delta (\mathrm{\Omega })`$ (cf. ), while analytical approximations for small $`m`$ and $`\mathrm{\Omega }`$ are available in the general case. We have observed that inertia tends to harden the synchronization transition: in the Kuramoto model ($`m=0`$) or with oscillators with identical natural frequencies, the synchronization transition is soft (supercritical bifurcation), whereas it becomes hard (subcritical bifurcation) in the cases of unimodal Lorentzian or discrete bimodal frequency distributions. In Section IV, we obtain approximations to stable time-dependent solutions of Eq. (6) in the “high frequency limit”, $`\mathrm{\Omega }\mathrm{}`$ . There are partially synchronized nonlinearly stable solutions of standing wave type, as in the Kuramoto model ). Finally, numerical results are presented in Section V, and compared to the approximate or exact solutions of previous Sections. Two Appendices at the end are devoted to technical details.
## II Linear stability of the incoherent solution
The incoherent solution is a $`\theta `$-independent stationary solution of (6). Its order parameter is $`r=0`$ according to (7). Then Eqs. (6), decay as $`\omega \pm \mathrm{}`$ and the normalization condition (8) yield the incoherent solution:
$$\rho _0(\omega ,\mathrm{\Omega })=\frac{1}{2\pi }\sqrt{\frac{m}{2\pi D}}e^{\frac{m}{2D}(\omega \mathrm{\Omega })^2}.$$
(9)
To analyze its linear stability, let us consider a small disturbance about incoherence,
$$\rho (\theta ,\omega ,\mathrm{\Omega },t)=\rho _0(\omega ,\mathrm{\Omega })+\epsilon \eta (\theta ,\omega ,\mathrm{\Omega },t)+O(\epsilon ^2),$$
(10)
where $`\epsilon 1`$. Normalization of $`\rho (\theta ,\omega ,\mathrm{\Omega },t)`$ then implies
$$_0^{2\pi }_{\mathrm{}}^+\mathrm{}\eta (\theta ,\omega ,\mathrm{\Omega },t)𝑑\omega 𝑑\theta =0.$$
(11)
We now introduce (10) into (6) and (7) and equate like terms in $`\epsilon `$. To order $`\epsilon `$, the result is
$`{\displaystyle \frac{\eta }{t}}+\omega {\displaystyle \frac{\eta }{\theta }}{\displaystyle \frac{1}{m}}{\displaystyle \frac{}{\omega }}[(\omega \mathrm{\Omega })\eta ]{\displaystyle \frac{D}{m^2}}{\displaystyle \frac{^2\eta }{\omega ^2}}=`$ (12)
$`{\displaystyle \frac{K\frac{\rho _0}{\omega }}{m}}{\displaystyle _0^{2\pi }}{\displaystyle _{\mathrm{}}^+\mathrm{}}{\displaystyle _{\mathrm{}}^+\mathrm{}}\eta (\varphi ,\omega ,\mathrm{\Omega },t)`$ (13)
$`\times \mathrm{sin}(\varphi \theta )g(\mathrm{\Omega })d\mathrm{\Omega }d\omega d\varphi .`$ (14)
We now insert a trial solution
$$\eta (\theta ,\omega ,\mathrm{\Omega },t)=e^{\lambda t}\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}b_n(\omega ;\mathrm{\Omega },\lambda )e^{in\theta }$$
(15)
(which is $`2\pi `$-periodic in $`\theta `$) into (14), thereby obtaining
$`{\displaystyle \frac{d^2b_n}{d\omega ^2}}+{\displaystyle \frac{m(\omega \mathrm{\Omega })}{D}}{\displaystyle \frac{db_n}{d\omega }}+{\displaystyle \frac{m(1m\lambda inm\omega )b_n}{D}}=`$ (16)
$`{\displaystyle \frac{\pi mK(i\delta _{n,1}i\delta _{n,1})\frac{\rho _0}{\omega }}{D}}1,b_n,`$ (17)
where we have defined the scalar product
$`\phi ,\psi ={\displaystyle _{\mathrm{}}^+\mathrm{}}{\displaystyle _{\mathrm{}}^+\mathrm{}}\overline{\phi (\omega ,\mathrm{\Omega })}\psi (\omega ,\mathrm{\Omega })g(\mathrm{\Omega })𝑑\mathrm{\Omega }𝑑\omega .`$ (18)
Notice that $`b_n=\overline{b_n}`$ and that $`b_n=0`$ because of the normalization condition (11).
Equation (17) can be transformed into a nonhomogeneous parabolic cylinder equation by the following change of variable:
$`b_n(\omega ;\mathrm{\Omega },\lambda )=\mathrm{exp}\left[{\displaystyle \frac{m(\omega \mathrm{\Omega })^2}{4D}}\right]\beta _n(w;\mathrm{\Omega },\lambda ),`$ (19)
$`w=\sqrt{{\displaystyle \frac{m}{D}}}(\omega \mathrm{\Omega }+2nDi).`$ (20)
Inserting (19) and (20) into (17), we obtain
$`{\displaystyle \frac{d^2\beta _n}{dw^2}}+\left[{\displaystyle \frac{1}{2}}{\displaystyle \frac{w^2}{4}}m(\lambda +in\mathrm{\Omega }+n^2D)\right]\beta _n=`$ (21)
$`i\pi K{\displaystyle \frac{\rho _0}{\omega }}e^{\frac{1}{4}(w2i\sqrt{mD})^2}1,e^{\frac{1}{4}(w2i\sqrt{mD})^2}\beta _1\delta _{n,1}.`$ (22)
(Recall that $`d\omega =\sqrt{D/m}dw`$ when using the definition of scalar product). Let us assume now that $`n\pm 1`$ and that $`\mathrm{\Omega }`$ is a fixed real number. Then the right hand side of (22) is zero and the resulting equation has the following eigenvalues
$`\lambda _{p,n}(\mathrm{\Omega })={\displaystyle \frac{p}{m}}n^2Din\mathrm{\Omega },p=0,1,2,\mathrm{},`$ (23)
associated to the eigenfunctions
$`\beta _{p,n}(w;\mathrm{\Omega },\lambda _{p,n})=D_p(w)=2^{\frac{p}{2}}e^{\frac{w^2}{4}}H_p\left({\displaystyle \frac{w}{\sqrt{2}}}\right),`$ (24)
which are independent of $`n`$ and $`\mathrm{\Omega }`$. In this formula, $`D_p(w)`$ and $`H_p(x)`$ are the parabolic cylinder function and the Hermite polynomial of index $`p`$, respectively . The eigenvalues $`\lambda _{p,n}(\mathrm{\Omega })`$ of (23), with $`n=\pm 1,\pm 2,\mathrm{}`$, $`p=0,1,\mathrm{}`$ and $`\mathrm{\Omega }`$ belonging to the support of $`g(\mathrm{\Omega })`$, constitute the continuous spectrum of the linear stability problem. In fact, a nonhomogeneous linear problem with a homogeneous part given by (22) cannot be solved for an arbitrary source term if $`\lambda =\lambda _{p,n}`$. Notice that the continuous spectrum lies to the left side of the imaginary axis if $`D>0`$ and $`n0`$. Then the “eigenvalues” (23) have negative real parts (and therefore correspond to stable modes). As we have already said, the neutrally stable modes with $`n=0`$ have zero amplitude due to the normalization condition (11).
If $`n=1`$, we can solve (17) by means of an expansion in eigenfunctions $`D_p(w)`$, $`p=0,1,2,\mathrm{}`$. To obtain the generalized Fourier coefficients of $`\beta _1`$, we multiply both sides of (17) by $`D_p(w)`$ and integrate over $`w`$. As $`_{\mathrm{}}^{\mathrm{}}D_p(x)D_n(x)𝑑x=\sqrt{2\pi }p!\delta _{pn}`$ (orthogonality condition, cf. §7.711.1 of ), the result is
$`\beta _1(\omega ;\mathrm{\Omega },\lambda )`$ $`=`$ $`{\displaystyle \frac{i\pi K}{m}}1,e^{\left(\frac{w}{2}i\sqrt{mD}\right)^2}\beta _1`$ (25)
$`\times `$ $`{\displaystyle \underset{p=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{_{\mathrm{}}^{\mathrm{}}e^{\left(\frac{w}{2}i\sqrt{mD}\right)^2}D_p\rho _0^{}𝑑w}{\sqrt{2\pi }p!\left(\frac{p}{m}+\lambda +i\mathrm{\Omega }+D\right)}}D_p(w),`$ (26)
where
$`\rho _0^{}(w)`$ $`=`$ $`{\displaystyle \frac{\rho _0}{\omega }}|_{\omega =\mathrm{\Omega }i2D+(D/m)^{\frac{1}{2}}w}`$ (27)
$`=`$ $`{\displaystyle \frac{m(wi2\sqrt{mD})}{(2\pi )^{\frac{3}{2}}D}}e^{\frac{1}{2}(wi2\sqrt{mD})^2}.`$ (28)
Once we have found $`\beta _1`$, we can calculate the scalar product $`1,e^{\left(\frac{w}{2}i\sqrt{mD}\right)^2}\beta _1`$. Since this scalar product appears as a factor in both sides of the resulting expression, we can divide by it, thereby obtaining an eigenvalue equation for $`\lambda `$:
$`1={\displaystyle \frac{i\pi K\sqrt{D}}{\sqrt{2\pi m^3}}}{\displaystyle \underset{p=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{p!}}{\displaystyle _{\mathrm{}}^+\mathrm{}}e^{\left(\frac{w}{2}i\sqrt{mD}\right)^2}D_p\rho _0^{}𝑑w`$ (29)
$`\times {\displaystyle _{\mathrm{}}^+\mathrm{}}e^{\left(\frac{w}{2}i\sqrt{mD}\right)^2}D_pdw`$ (30)
$`\times {\displaystyle _{\mathrm{}}^+\mathrm{}}{\displaystyle \frac{g(\mathrm{\Omega })}{\frac{p}{m}+\lambda +i\mathrm{\Omega }+D}}d\mathrm{\Omega }.`$ (31)
In Appendix A, we show that this equation may be rewritten as
$`1={\displaystyle \frac{K}{4\pi \sqrt{mD}}}{\displaystyle \underset{p=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{𝒜_p(\sqrt{mD})𝒜_p^{}(\sqrt{mD})}{p!}}`$ (32)
$`\times {\displaystyle _{\mathrm{}}^+\mathrm{}}{\displaystyle \frac{g(\mathrm{\Omega })}{\frac{p}{m}+\lambda +i\mathrm{\Omega }+D}}d\mathrm{\Omega },`$ (33)
where $`𝒜_p(x)`$ is defined as
$`𝒜_p(x)={\displaystyle _{\mathrm{}}^+\mathrm{}}D_p(w)e^{\left(\frac{w}{2}ix\right)^2}𝑑w.`$ (34)
The result of evaluating this integral is (cf. Appendix A):
$`𝒜_p(x)=i^p\sqrt{2\pi }e^{\frac{x^2}{2}}x^p(x>0).`$ (35)
Inserting (35) in (33), we obtain
$`1={\displaystyle \frac{Ke^{mD}}{2}}{\displaystyle \underset{p=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(mD)^p\left(1+\frac{p}{mD}\right)}{p!}}`$ (36)
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{g(\mathrm{\Omega })d\mathrm{\Omega }}{\lambda +D+i\mathrm{\Omega }+\frac{p}{m}}}.`$ (37)
As $`m0`$, this equation coincides with that obtained for the Kuramoto model, . Equation (37) can be rewritten in terms of incomplete gamma functions as follows :
$`{\displaystyle \frac{2D}{K}}=e^{mD}{\displaystyle _{\mathrm{}}^{\mathrm{}}}[mD\gamma (m(\lambda +D+i\mathrm{\Omega }),mD)`$ (38)
$`\gamma (1+(\lambda +D+i\mathrm{\Omega })m,mD)]{\displaystyle \frac{g(\mathrm{\Omega })d\mathrm{\Omega }}{(mD)^{m(\lambda +D+i\mathrm{\Omega })}}}`$ (39)
$`={\displaystyle _{\mathrm{}}^{\mathrm{}}}\left[1{\displaystyle \frac{(\lambda +i\mathrm{\Omega })me^{mD}\gamma (m(\lambda +D+i\mathrm{\Omega }),mD)}{(mD)^{(\lambda +D+i\mathrm{\Omega })m}}}\right]`$ (40)
$`\times g(\mathrm{\Omega })d\mathrm{\Omega },`$ (41)
where
$`\gamma (a,x)={\displaystyle _0^x}e^tt^{a1}𝑑t.`$ (42)
¿From now on, we analyze Eq. (37) for special frequency distributions:
(a) Unimodal frequency distribution, $`g(\mathrm{\Omega })=\delta (\mathrm{\Omega })`$.
In this case we show that if $`\text{Re}\lambda =0`$, then $`\text{Im}\lambda =0`$. Thus, the eigenvalues that may acquire a positive real part are real. Then the critical coupling is obtained by setting $`\lambda =0`$. By subtracting from Eq. (37) its complex conjugate, we obtain
$`0=\text{Im}(\lambda )f(\text{Im}\lambda ,m,D),`$ (43)
where
$$f(\text{Im}\lambda ,m,D)=\underset{p=0}{\overset{\mathrm{}}{}}\frac{(mD)^p\left(1+\frac{p}{mD}\right)}{p!\left[(\text{Im}\lambda )^2+\left(D+\frac{p}{m}\right)^2\right]}.$$
(44)
Notice than the even function $`f(\text{Im}\lambda ,m,D)`$ decreases monotonically with Im$`(\lambda )>0`$. On the other hand, $`f`$ tends to zero, as $`\text{Im}\lambda +\mathrm{}`$, and
$$f\frac{m}{D}(mD)^{mD}\gamma (mD,mD)>0,\text{as }\text{Im}\lambda 0.$$
(45)
Thus $`f`$ does not vanish at finite values of Im$`\lambda `$, and therefore the only solution of (43) is Im$`\lambda =0`$. Setting now $`\lambda =0`$, Equation (37) yields the critical coupling $`K=K_c`$,
$`{\displaystyle \frac{K_c}{2D}}e^{mD}\left[{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(mD)^n}{n!}}\right]=1,\text{and therefore}`$ (46)
$`K_c=2D.`$ (47)
Figures 1 and 2 show the largest eigenvalue $`\lambda `$ as a function of $`m`$ and $`K`$. To compute $`\lambda `$ numerically, we used from Eq. (10), and (15), that the amplitude of the order parameter is
$$rCe^{\lambda t},$$
(48)
close to incoherence. Then, the goal is to simulate the evolution of the system, choosing the initial condition sufficiently close to the incoherent solution, and obtain numerically the amplitude order parameter $`r(t)`$. Fig. 2 shows that different eigenvalue curves (for different $`m`$) intersect the horizontal axis, $`\lambda =0`$, at the same value of $`K`$, as expected from Eq. (47).
(b) Unimodal Lorentzian frequency distribution, $`g(\mathrm{\Omega })=\frac{\epsilon /\pi }{\epsilon ^2+\mathrm{\Omega }^2}`$.
Eq. (37) becomes
$`1={\displaystyle \frac{K}{2D}}e^{mD}[{\displaystyle \frac{D}{\lambda +D+\epsilon }}+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(mD)^n}{n!}}`$ (49)
$`\times {\displaystyle \frac{n+mD}{m(\lambda +D+\epsilon )+n}}]={\displaystyle \frac{K}{2}}me^{mD}`$ (50)
$`\times (mD)^{m(\lambda +D+\epsilon )}[(mD)^{m(\lambda +D+\epsilon )1}e^{mD}`$ (51)
$`\gamma (m(\lambda +D+\epsilon ),mD)({\displaystyle \frac{\lambda +D+\epsilon }{D}}1)]`$ (52)
An explicit solution for $`\lambda `$ cannot in general be found. Thus, we consider several limiting cases corresponding to physically interesting parameter choices. In the small noise limit, $`D1`$, we consider the cases $`(i)`$ $`m=O(1)`$ fixed, and $`(ii)`$ $`mD=1`$. It is remarkable that the expansion
$$\gamma (a,x)=e^xx^a\underset{n=0}{\overset{\mathrm{}}{}}\frac{x^n}{(a)_{n+1}},$$
(53)
where $`(a)_k=a(a+1)\mathrm{}(a+k1)`$, $`k=1,2,\mathrm{}`$, holds in both cases. If $`x=mD0`$, $`a=mD+m(\lambda +\epsilon )>x`$, (53) holds as a convergent expansion, . If $`x=mD=1`$, and $`a=1+m(\lambda +\epsilon )\mathrm{}`$ (with fixed $`\lambda `$ and $`\epsilon `$ of order 1), (53) holds as an asymptotic expansion, . Inserting (53) in Eq. (52), we obtain
$`{\displaystyle \frac{2D}{K}}=1+{\displaystyle \frac{xa}{a}}\left[1+{\displaystyle \frac{x}{a+1}}+{\displaystyle \frac{x^2}{(a+1)(a+2)}}+O\left({\displaystyle \frac{x^3}{a^3}}\right)\right]`$ (54)
Similarly to the unimodal case, it is possible to prove that $`\lambda `$ is always real. To this purpose, notice that replacing $`\lambda +\epsilon `$ in eq. (52) with $`\lambda `$ in eq. (37) (setting $`\mathrm{\Omega }=0`$) we obtain the same equation. The critical coupling $`K=K_c`$ is then found by setting $`\lambda =0`$ in (54). In case $`(i)`$, we have
$$K_c=2\epsilon (m\epsilon +1)+\frac{2(2+3m\epsilon )}{2+m\epsilon }D+O\left(D^2\right),$$
(55)
In the limit of vanishing mass, we recover the result $`K_c=2(D+\epsilon )`$ valid for the Kuramoto model . Another important limit is $`\epsilon =0`$, which reproduces the unimodal distribution. We find $`K_c=2D`$, independently of mass. Thus the spread in frequency distribution plays an important role in synchronizing populations of oscillators affected by inertia.
In case $`(ii)`$, Eq. (54) yields
$$1=\frac{Kx}{2Da(a+1)}+O\left(a^3\right),$$
(56)
from which
$$\lambda =\left[\epsilon +\frac{1}{2m}\left(3\pm \sqrt{1+2Km}\right)\right].$$
(57)
This quantity is always real and vanishes for $`K=K_c=2\epsilon (m\epsilon +3)+\frac{4}{m}`$. Note that $`K_c`$ grows roughly linearly with $`m`$. Thus oscillator synchronization is made harder by increasing inertia in the limit of vanishing noise. This behavior is slightly different from that described in . There numerical simulations seemed to show that incoherence remains stable up to a critical coupling, which was independent of $`m`$. The singular nature of the limit $`D0`$ makes the cause of this discrepancy unclear, although we should mention that no stability analysis was conducted in . In the opposite limit $`m\mathrm{}`$, $`\lambda \epsilon `$, and incoherence is always stable.
The stability diagram in the parameter space ($`\epsilon `$,$`K`$) is shown in Figure 3 for $`m=0.2`$, and compared to that of the Kuramoto model ($`m=0`$). This diagram is obtained from Eq. (52) with $`\lambda =0`$, for fixed $`D`$ and $`m`$. In this figure, we have also plotted the evolution of the order parameter amplitude for the parameter values marked in the stability diagram by (1) to (4). In all cases, the initial condition is taken sufficiently close to the incoherent solution, $`r=0`$.
(c) Bimodal frequency distribution, $`g(\mathrm{\Omega })=\frac{1}{2}[\delta (\mathrm{\Omega }\mathrm{\Omega }_0)+\delta (\mathrm{\Omega }+\mathrm{\Omega }_0)]`$.
Eq. (37) becomes
$`{\displaystyle \frac{Ke^{mD}}{2D}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(mD)^n}{n!}}{\displaystyle \frac{(n+mD)[m(\lambda +D)+n]}{[m(\lambda +D)+n]^2+m^2\mathrm{\Omega }_0^2}}=1.`$ (58)
In the high frequency limit, $`\mathrm{\Omega }_0\mathrm{}`$, we can find an analytical formula for $`\lambda `$ by inserting the following asymptotic expansion for the incomplete gamma function in (41) ,
$$\gamma (a,x)\frac{e^xx^a}{a},a\mathrm{},$$
(59)
where $`a=m(\lambda +D+i\mathrm{\Omega }_0)`$, and $`x=mD`$. The result is
$$1=\frac{K}{4}\left(\frac{1}{\lambda +D+i\mathrm{\Omega }_0}+\frac{1}{\lambda +Di\mathrm{\Omega }_0}\right),$$
(60)
which yields
$$\lambda =D+\frac{K}{4}+\frac{K}{4}i\sqrt{16\mathrm{\Omega }_0^2K^2}.$$
(61)
Re$`\lambda =0`$ gives the same critical coupling as the Kuramoto model, $`K_c=4D`$, for the same bimodal frequency distribution .
Figure 4 shows the stability diagram, which is obtained from Eq. (58) with Re$`\lambda =0`$, for $`D=1`$ and three different mass values $`m=0.1,\mathrm{\hspace{0.17em}1},\mathrm{\hspace{0.17em}6}`$. The stability diagram of the Kuramoto model ($`m=0`$) is also depicted. Notice how the curves corresponding to differerent masses tend to $`K=4D`$ as $`\mathrm{\Omega }_0\mathrm{}`$, as expected. Fig. 5 displays the evolution of the order parameter amplitude in different regions of the stability diagram corresponding to $`m=0.8`$ and $`D=1`$.
In conclusion, increasing $`m`$, $`D`$, $`\epsilon `$ (inertia, noise, and frequency spread) makes it more difficult to synchronize the oscillator population via stationary bifurcations from incoherence.
## III Mode-coupling equations and stationary solutions
Inspired by the previous linear stability analysis, we shall expand the distribution function using a basis of parabolic cylinder functions (or, equivalently, Hermite polynomials) of unit mean square norm ,
$`\rho (\theta ,\omega ,\mathrm{\Omega },t)=\left({\displaystyle \frac{2\pi D}{m}}\right)^{\frac{1}{4}}e^{\frac{m\omega ^2}{4D}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}c_n(\theta ,\mathrm{\Omega },t)\psi _n(\omega ).`$ (62)
In (62), we have defined
$`\psi _n(\omega )=\left(n!\sqrt{{\displaystyle \frac{2\pi D}{m}}}\right)^{\frac{1}{2}}D_n\left(\sqrt{{\displaystyle \frac{m}{D}}}\omega \right)`$ (63)
$`=\left(n!\mathrm{\hspace{0.17em}2}^n\sqrt{{\displaystyle \frac{2\pi D}{m}}}\right)^{1/2}H_n\left(\sqrt{{\displaystyle \frac{m}{2D}}}\omega \right)e^{m\omega ^2/4D},`$ (64)
so that $`_{\mathrm{}}^{\mathrm{}}\psi _n\psi _p𝑑\omega =\delta _{np}`$. The functions $`c_n(\theta ,\mathrm{\Omega },t)`$ are $`2\pi `$-periodic in $`\theta `$, and we have
$$_0^{2\pi }c_0(\theta ,\mathrm{\Omega },t)𝑑\theta =1,$$
(65)
as it follows from the normalization of $`\rho (\theta ,\omega ,\mathrm{\Omega },t)`$.
We shall find a system of mode-coupling equations for the coefficient functions $`c_n`$. Then we shall try to find stationary solutions for different frequency distributions $`g(\mathrm{\Omega })`$. This is not so easy in the general case, so that we shall follow a standard approach: We shall recognize in the system of equations for the coefficient functions a particular stationary solution corresponding to incoherence. Then we shall try to find a bifurcation equation for other stationary solutions branching off from incoherence. We shall see that even this requires different approximation schemes in order to succeed.
### A Mode-coupling equations
Let us insert Eq. (62) into the Fokker-Planck equation (6). We then obtain the following hierarchy of coupled partial differential equations for $`c_n(\theta ,\mathrm{\Omega },t)`$,
$`{\displaystyle \frac{c_0}{t}}=\sqrt{{\displaystyle \frac{D}{m}}}{\displaystyle \frac{c_1}{\theta }},`$ (66)
$`{\displaystyle \frac{c_1}{t}}=c_0{\displaystyle \frac{1}{m}}c_1\sqrt{2}\sqrt{{\displaystyle \frac{D}{m}}}{\displaystyle \frac{c_2}{\theta }},`$ (67)
$`\mathrm{}`$ (68)
$`{\displaystyle \frac{c_n}{t}}=\sqrt{n}c_{n1}{\displaystyle \frac{n}{m}}c_n\sqrt{n+1}\sqrt{{\displaystyle \frac{D}{m}}}{\displaystyle \frac{c_{n+1}}{\theta }}.`$ (69)
Here we have defined the operator
$`f=\sqrt{{\displaystyle \frac{D}{m}}}\left[{\displaystyle \frac{}{\theta }}{\displaystyle \frac{\mathrm{\Omega }+Kr\mathrm{sin}(\psi \theta )}{D}}\right]f.`$ (70)
In terms of the functions $`c_n`$, the equation for the order parameter becomes
$`re^{i\psi }={\displaystyle _0^{2\pi }}{\displaystyle _{\mathrm{}}^+\mathrm{}}e^{i\theta }c_0(\theta ,\mathrm{\Omega },t)g(\mathrm{\Omega })𝑑\mathrm{\Omega }𝑑\theta .`$ (71)
$`r\mathrm{sin}(\psi \theta )={\displaystyle _0^{2\pi }}{\displaystyle _{\mathrm{}}^+\mathrm{}}\mathrm{sin}(\varphi \theta )c_0(\varphi ,\mathrm{\Omega },t)g(\mathrm{\Omega })𝑑\mathrm{\Omega }𝑑\varphi .`$ (72)
### B Incoherence and bifurcating stationary solutions
The incoherent $`\theta `$-independent solution, $`c_n=C_n(\mathrm{\Omega })`$, depends only on $`\mathrm{\Omega }`$. It can readily be obtained from the above hierarchy of equations, by ignoring all derivatives. The result is
$$C_n(\mathrm{\Omega })=\frac{1}{2\pi }\frac{1}{\sqrt{n!}}\left(\frac{m}{D}\right)^{n/2}\mathrm{\Omega }^n.$$
(73)
Inserting this into the expansion in Eq. (62), we indeed recover Eq. (9). In particular setting $`\mathrm{\Omega }=0`$, this yields the simplest incoherent solution, $`C_n=1/(2\pi )`$, corresponding to the unimodal frequency distribution. Other interesting stationary solutions are partially synchronized distributions, which depend on $`\theta `$. Notice that eq. (66) implies that $`c_1`$ does not depend on $`\theta `$.
In the following we analyze stationary solutions bifurcating from incoherence for different frequency distributions.
(a) Discrete unimodal frequency distribution, $`g(\mathrm{\Omega })=\delta (\mathrm{\Omega })`$.
Let us look for stationary solutions with finitely many nonzero coefficients $`c_n=0`$ for $`n>N`$ and $`g(\mathrm{\Omega })=\delta (\mathrm{\Omega })`$. Eq. (69) yields $`c_N=0`$, and therefore, $`c_N=K_Ne^{\frac{K}{D}r\mathrm{cos}(\psi \theta )}`$ ($`K_N=`$constant). Inserting this result in (69) for $`n=N1`$, we find $`c_{N1}=\sqrt{N}c_N/m`$. The corresponding solution is not periodic in $`\theta `$ unless $`K_N=0`$. Repeating this argument, we obtain the stationary solution
$$c_0(\theta )=\frac{e^{\frac{K}{D}r\mathrm{cos}(\psi \theta )}}{_0^{2\pi }e^{\frac{K}{D}r\mathrm{cos}(\psi \theta )}𝑑\theta },$$
(74)
where $`c_n=0`$ for $`n>0`$ and the normalization condition has been used. Observe that the resulting distribution, $`\rho (\theta ,\omega )=\left(\frac{2\pi D}{m}\right)^{1/4}c_0(\theta )\psi _0(\omega )e^{m\omega ^2/4D}`$, is factorized with respect to its two arguments, $`\theta `$ and $`\omega `$, cf. . It is remarkable that $`c_0(\theta )`$ is independent of $`m`$, and coincides with that obtained for the Kuramoto model. Therefore, from (71) the order parameter does not depend on inertia, and the bifurcation diagram for $`r`$ is exactly the same as that for the Kuramoto model.
(b) General frequency distributions.
For general frequency distributions, we could try to find stationary solutions of the mode-coupling equations (69) which bifurcate from incoherence. It is however more direct to work with the stationary Fokker-Planck equation (6) as follows. We consider stationary solutions as functions of a fixed value of the synchronization parameter $`r`$. Then we expand these solutions in power series of $`r`$ and write a hierarchy of equations for the coefficient functions. The first coefficient function should be the incoherent solution of synchronization parameter $`r=0`$. Inserting the power series probability density function into Eq. (7), we find the amplitude equation for stationary solutions bifurcating from incoherence. This procedure is explained in Appendix B. We quote here the result:
$$r=\frac{Kr}{2D}\alpha +\frac{(Kr)^3}{6}\beta +O((Kr)^4),$$
(75)
where
$`\alpha =e^{mD}[{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{1}{1+\frac{\mathrm{\Omega }^2}{D^2}}}g(\mathrm{\Omega })d\mathrm{\Omega }`$ (76)
$`+{\displaystyle \underset{p=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(mD)^p(1+\frac{p}{mD})^2}{p!}}`$ (77)
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{1}{(1+\frac{p}{mD})^2+\frac{\mathrm{\Omega }^2}{D^2}}}g(\mathrm{\Omega })d\mathrm{\Omega }],`$ (78)
and $`\beta `$ can be calculated numerically by solving a system of differential equations. As the inertia vanishes, $`mD0+`$, $`\alpha `$ becomes
$`\alpha ={\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{g(\mathrm{\Omega })d\mathrm{\Omega }}{1+\frac{\mathrm{\Omega }^2}{D^2}}}{\displaystyle \frac{m}{D}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{\mathrm{\Omega }^2}{1+\frac{\mathrm{\Omega }^2}{D^2}}}g(\mathrm{\Omega })𝑑\mathrm{\Omega }+O(m^2D^2).`$ (79)
Notice that $`\alpha `$ in Equation (78) coincides with $`2D/K`$ in Eq. (37) provided $`\lambda =0`$ and $`g(\mathrm{\Omega })=g(\mathrm{\Omega })`$. This means that the critical coupling for bifurcation towards stationary synchronized states is obtained at $`K\alpha =2D`$, no matter what the symmetric frequency distribution may be.
Before interpreting the amplitude equation (75), we shall outline a procedure to obtain its coefficients based upon an uncontrolled closure assumption which is accurate for small values of $`mD`$. Consider the expansion (62). We expect that the coefficients $`c_n`$ of stationary solutions close to incoherence do not differ much from the coefficients of the latter. In view of the functional form (73), we anticipate that the coefficients $`c_n`$ approach zero as $`n\mathrm{}`$ faster for smaller values of the mass. Thus we shall now consider stationary solutions for general frequency distributions, such that $`c_n=0`$ in (69), for all $`n3`$. By solving equation (69) for such a stationary solution with $`n=2`$, we obtain
$$c_2(\theta ,\mathrm{\Omega })=\sqrt{\frac{m}{2D}}[\mathrm{\Omega }+Kr\mathrm{sin}(\psi \theta )]c_1(\mathrm{\Omega }).$$
(80)
We now insert this expression into (67). The resulting equation is solved for a $`c_0`$, which is $`2\pi `$-periodic in $`\theta `$ and obeys the normalization condition (8). We find
$`c_0(\theta ,\mathrm{\Omega })={\displaystyle \frac{e^{\frac{Kr}{D}\mathrm{cos}(\psi \theta )}\phi (\theta ,\mathrm{\Omega })}{Z(\mathrm{\Omega })}}`$ (81)
$`\phi (\theta ,\mathrm{\Omega })={\displaystyle _0^{2\pi }}[1mKr\mathrm{cos}(\psi \theta \eta )]`$ (82)
$`\times e^{\frac{1}{D}[\mathrm{\Omega }\eta +Kr\mathrm{cos}(\psi \theta \eta )]}d\eta ,`$ (83)
$`Z(\mathrm{\Omega })={\displaystyle _0^{2\pi }}e^{\frac{Kr}{D}\mathrm{cos}(\psi \theta )}\phi (\theta ,\mathrm{\Omega })𝑑\theta ,`$ (84)
In (81), $`r`$ should be determined so that
$`r={\displaystyle _0^{2\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}c_0(\theta ,\mathrm{\Omega })\mathrm{cos}(\psi \theta )g(\mathrm{\Omega })𝑑\mathrm{\Omega }𝑑\theta ,`$ (86)
holds. The function $`c_1(\mathrm{\Omega })`$ can be obtained by integrating (67) with respect to $`\theta `$ and using the normalization condition for $`c_0`$ together with the $`2\pi `$-periodicity in $`\theta `$ of $`c_0`$ and $`c_2`$:
$$c_1=\frac{1}{2\pi }\sqrt{\frac{m}{D}}\left[\mathrm{\Omega }+Kr_0^{2\pi }\mathrm{sin}(\psi \theta )c_0(\theta ,\mathrm{\Omega })𝑑\theta \right].$$
(87)
Then, the stationary distribution can be approximated by
$`\rho (\theta ,\omega ,\mathrm{\Omega })\left({\displaystyle \frac{2\pi D}{m}}\right)^{1/4}[c_0(\theta ,\mathrm{\Omega })\psi _0(\omega )+c_1(\mathrm{\Omega })\psi _1(\omega )`$ (88)
$`+c_2(\theta ,\mathrm{\Omega })\psi _2(\omega )]e^{m\omega ^2/4D}.`$ (89)
Notice that a nonvanishing $`c_1(\mathrm{\Omega })`$ in (62) implies that the probability density is no longer even in $`\omega `$ ($`\psi _1`$ is an odd function of $`\omega `$). By (87), this occurs in the synchronized phase for the case of nonidentical oscillators. For the bimodal frequency distribution, Fig. 6 shows $`c_1`$ as function of $`\mathrm{\Omega }_0`$ for two different values of $`m`$. Both the approximate expression (87) and results of direct numerical simulations are depicted. Notice that the agreement between our approximation and the numerical result improves as $`m`$ decreases. Thus we observe that for each fixed nonzero $`\mathrm{\Omega }`$ and each fixed $`\theta `$, the distribution function is no longer peaked at $`\omega =0`$. However, the instantaneous frequency distribution, defined by
$$_0^{2\pi }_{\mathrm{}}^{\mathrm{}}\rho (\theta ,\omega ,\mathrm{\Omega })g(\mathrm{\Omega })𝑑\mathrm{\Omega }𝑑\theta ,$$
may turn out to be even in $`\omega `$. This is certainly true for the approximate stationary distribution (89), for
$$_{\mathrm{}}^{\mathrm{}}c_1(\mathrm{\Omega })g(\mathrm{\Omega })𝑑\mathrm{\Omega }=\sqrt{\frac{m}{D}}\frac{Kr}{2\pi }_0^{2\pi }_{\mathrm{}}^{\mathrm{}}c_0(\theta ,\mathrm{\Omega })\mathrm{sin}(\psi \theta )g(\mathrm{\Omega })𝑑\mathrm{\Omega }𝑑\theta =0,$$
as it follows from the definition of the order parameter and the expansion (62).
Another indication that the exact instantaneous frequency distribution may be even in $`\omega `$ is that the average frequency tends to zero as $`t+\mathrm{}`$. This would occur if the stable stationary distribution is even in $`\omega `$ (although, admittedly, distributions which are not even in $`\omega `$ may have zero mean). The result can be shown directly from the Fokker-Planck equation (6). We multiply that equation by $`\omega g(\mathrm{\Omega })`$ and integrate with respect to all variables. Then we obtain the following equation for the mean value $`\omega `$:
$$\frac{d}{dt}\omega +\frac{\omega }{m}=_{\mathrm{}}^{\mathrm{}}\mathrm{\Omega }g(\mathrm{\Omega })𝑑\mathrm{\Omega }.$$
The right hand side of this expression is zero if $`g(\mathrm{\Omega })=g(\mathrm{\Omega })`$. Then $`\omega `$ tends to zero exponentially fast as $`t+\mathrm{}`$.
Let us now go back to the problem of obtaining the coefficients in the amplitude equation (75). Let us expand (81) in powers of $`r`$ and insert the result in (86). We obtain an approximation to the coefficients of the amplitude equation (86). $`\alpha `$ is again given by (75), and the expression for $`\beta `$ is
$`\beta ={\displaystyle _{\mathrm{}}^{\mathrm{}}}[{\displaystyle \frac{3\frac{\mathrm{\Omega }^2}{D^2}\frac{3}{2}+\frac{39}{4}m\frac{\mathrm{\Omega }^2}{D}}{(1+\frac{\mathrm{\Omega }^2}{D^2})^2(4+\frac{\mathrm{\Omega }^2}{D^2})D^3}}`$ (90)
$`+{\displaystyle \frac{3m^2\mathrm{\Omega }^2\left(1+\frac{\mathrm{\Omega }^2}{2D^2}\right)}{2D^3\left(1+\frac{\mathrm{\Omega }^2}{D^2}\right)^2}}]g(\mathrm{\Omega })d\mathrm{\Omega }.`$ (91)
We now interpret the amplitude equation (75) in the usual way. Notice that the nonzero solution is approximately given by
$$Kr\sqrt{\frac{3(2DK\alpha )}{K\beta D}}.$$
Then the critical value of $`K`$ is $`K^{}=2D/\alpha `$, and the sign of $`\beta `$ determines the direction of the bifurcating branch of stationary solutions. Assume $`\alpha >0`$. Then the bifurcating stationary solution exists for $`K>K^{}`$ if $`\beta <0`$ (supercritical bifurcation). If $`\beta >0`$, we have a subcritical bifurcation and the partially synchronized stationary solution exists for $`K<K^{}`$. The critical coupling tends to infinity as $`\alpha 0+`$, and the bifurcation does not occur at positive couplings if $`\alpha <0`$. Given the formulas (78) or (79), $`\alpha `$ could become negative, depending on the value of $`m`$. Let us now analyze the bifurcation diagram corresponding to two different frequency distributions.
* Unimodal Lorentzian frequency distribution.
For an unimodal Lorentzian frequency distribution, the coefficients $`\alpha `$ and $`\beta `$ of (79) and (91) are approximately given by
$`\alpha ={\displaystyle \frac{D}{D+\gamma }}(1m\epsilon ),`$ (92)
$`\beta ={\displaystyle \frac{3}{4}}{\displaystyle \frac{1}{(D+\epsilon )^2(2D+\epsilon )}}[1+{\displaystyle \frac{m}{D}}\epsilon (3D+\epsilon )].`$ (93)
Apparently, the sign of $`\alpha `$ could again be negative provided $`m\epsilon `$ is sufficiently large. However, evaluation of the exact expression (78) shows that $`\alpha >0`$. Figure 7 shows $`\alpha `$ as a function of $`\epsilon `$ for two different masses. Notice that the approximate expression for $`\alpha `$ fits better the numerical result as $`m`$ and $`\epsilon `$ decrease.
In contrast with $`\alpha `$, $`\beta `$ can really change sign for a Lorentzian frequency distribution. Fig. 8 shows the good qualitative agreement between the approximate expression for $`\beta `$ and its numerical evaluation from the exact equations (B22).
There is a critical mass $`m_c`$ ($`\epsilon `$ and $`D`$ are kept fixed) for which $`\beta =0`$. This mass separates the supercritical and subcritical bifurcation regions. Setting $`\beta =0`$ in eq. (93) yields
$$m_c=\frac{D}{\epsilon (3D+\epsilon )}.$$
(94)
Figure 9 shows the qualitative agreement between the approximate expression (94), and the numerical result obtained from the solution of (B22).
Note that in the massless case, the stationary solution always branches off supercritically, independently of $`\epsilon `$. In the presence of inertia and for appropriate values of $`\epsilon `$, we have found a subcritical bifurcation. This prediction is illustrated by Figures 10 (parameters corresponding to a supercritical bifurcation) and 11 (parameters corresponding to a subcritical bifurcation). Figure 10 is a typical supercritical bifurcation diagram (order parameter amplitude versus $`K`$). In the subcritical case, Figure 11 shows the different evolution of the order parameter amplitude for two different initial conditions. Notice the bistability between incoherence and the stable synchronized solution typical of a subcritical bifurcation.
* Bimodal frequency distribution
For the bimodal frequency distribution, $`\alpha `$ may be zero. This means that there is a critical frequency $`\mathrm{\Omega }_0^c(m)`$ above which $`K^{}=\mathrm{}`$. In fact, we have shown in Fig. 4, looking to the branch corresponding to $`\lambda =0`$, that in case of a bimodal frequency distribution, it is possible to find some values of $`\mathrm{\Omega }_0`$ without its corresponding $`K`$ in the stability diagram of incoherence, due to the existence of an horizontal asymptote in the space parameter $`(K,\mathrm{\Omega }_0)`$. This means that there is not a finite coupling $`K^{}`$ where the stationary solution branches off. On the other hand, this never occurs in the bimodal Kuramoto model \[$`m=0`$, $`K_c/D=2(1+\mathrm{\Omega }_0^2/D^2)`$\], because in this case there is not such an horizontal asymptote.
Eq. (79) shows that the critical frequency at which $`\alpha =0`$ is
$$\frac{\mathrm{\Omega }_0^{\mathrm{}}}{D}=\sqrt{\frac{1}{mD}}.$$
(95)
Figure 12 compares the previous approximate expression to the exact critical frequency obtained from (78). Note that the approximation improves as $`m`$ decreases. Fig. 13 shows that $`K^{}`$ increases as $`\mathrm{\Omega }_0`$ does. $`K^{}`$ becomes infinity for $`\mathrm{\Omega }_0\mathrm{\Omega }_0^{\mathrm{}}`$.
For $`\alpha >0`$, there is another important critical frequency. In this case, the sign of $`\beta `$ decides whether the bifurcation is subcritical ($`\beta >0`$) or supercritical ($`\beta <0`$). The sign of $`\beta `$ depends on $`m`$ and $`\mathrm{\Omega }_0`$. For small mass, we can use the approximations (79) and (91) (ignoring terms which are quadratic in the mass). Then we find that the critical frequency at which $`\beta =0`$ is
$$\frac{\mathrm{\Omega }_0^c}{D}=\sqrt{\frac{1}{2+\frac{13}{8}mD}}.$$
(96)
We have solved numerically the system of equations (B22), and compared with the approximation (91) for $`\beta `$. In Figure 14 the coefficient $`\beta `$ is plotted as a function of $`\mathrm{\Omega }_0`$ for different masses. Similarly, Figure 15 shows how the critical frequency $`\mathrm{\Omega }_0^c`$ varies as a function of $`m`$. Note that the analytical approximation improves as $`m`$ goes to zero (as expected).
Notice that $`\mathrm{\Omega }_0^c<\mathrm{\Omega }_0^{\mathrm{}}`$. Then the branch of synchronized stationary states bifurcates subcritically from incoherence at $`K^{}=\mathrm{}`$, provided $`\mathrm{\Omega }_0>\mathrm{\Omega }_0^{\mathrm{}}`$.
Summarizing, inertia favors the subcritical character of bifurcations describing the transition from incoherence to the partial synchronized state. In fact in the bimodal case, the transition of incoherence to synchronization will most likely occur as a subcritical bifurcation as inertia increases. For a continuous unimodal Lorentzian frequency distribution, inertia may turn subcritical the supercritical bifurcation to stationary synchronized states, which is always found in the masslesss Kuramoto model.
## IV High frequency limit
The high frequency limit of Eq. (6) can be analyzed by means of a multiscale method. For the Kuramoto model , this method lead to the result that the probability density is (to leading order) a superposition of different components corresponding to the different peaks of the oscillator frequency distribution. To apply this method to the present model, we shall assume that the frequency distribution $`g(\mathrm{\Omega })`$ has $`m`$ maxima located at $`\mathrm{\Omega }_0\nu _l,l=1,\mathrm{}.,m`$, so that $`g(\mathrm{\Omega })d\mathrm{\Omega }`$ tends to the distribution
$$\mathrm{\Gamma }(\nu )\underset{l=1}{\overset{m}{}}\alpha _l\delta (\nu \nu _l)d\nu ,$$
(97)
as $`\mathrm{\Omega }_0\mathrm{}`$. In order to simplify the calculations, we change variables to a comoving frame:
$$\beta =\theta \mathrm{\Omega }t\theta \frac{\nu }{ϵ}t,$$
(98)
and let $`\omega ^{}=\omega \mathrm{\Omega }`$, where
$$ϵ=\frac{1}{\mathrm{\Omega }_0}1$$
(99)
Then we obtain the following equations:
$`{\displaystyle \frac{\rho _j}{t}}=D{\displaystyle \frac{^2\rho _j}{\omega ^2}}{\displaystyle \frac{1}{m}}{\displaystyle \frac{}{\omega ^{}}}(\rho _jU_j)\omega ^{}{\displaystyle \frac{\rho _j}{\beta }},`$ (100)
$`U_j=\omega ^{}+\text{Im}K\{{\displaystyle \underset{l=1}{\overset{m}{}}}\alpha _le^{i(\nu _l\nu _j)t/ϵ}`$ (101)
$`\times {\displaystyle }e^{i(\beta ^{}\beta )}\rho (\beta ^{},\omega ^{},t,\nu _l;ϵ)d\beta ^{}d\omega ^{}\},`$ (102)
$`{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle _0^{2\pi }}\rho _j(\beta ^{},t,\omega ^{};ϵ)𝑑\beta ^{}𝑑\omega ^{}=1,`$ (103)
where $`\rho _j=\rho (\beta ,t,\omega ^{},\nu _j;ϵ)`$, and $`\rho \alpha _j\rho _j\delta (\nu \nu _j)`$, . We now define fast and slow time scales, $`\tau =t/ϵ`$ and $`t`$, and make the Ansatz:
$$\rho =\underset{n=0}{\overset{2}{}}\rho ^{(n)}(\beta ,\omega ^{},\tau ,t,\nu )ϵ^n+O(ϵ^3).$$
(104)
Inserting this into the governing equations, we obtain the following hierarchy:
$`{\displaystyle \frac{\rho _j^{(0)}}{\tau }}=0,`$ (105)
$`{\displaystyle \frac{\rho _j^{(1)}}{\tau }}={\displaystyle \frac{\rho _j^{(0)}}{t}}+D{\displaystyle \frac{^2\rho _j^{(0)}}{\omega ^2}}+{\displaystyle \frac{1}{m}}{\displaystyle \frac{}{\omega ^{}}}(\omega ^{}\rho _j^{(0)})`$ (106)
$`\omega ^{}{\displaystyle \frac{\rho _j^{(0)}}{\beta }}{\displaystyle \frac{K}{m}}{\displaystyle \frac{}{\omega ^{}}}\{\rho _j^{(0)}[\text{Im}(\alpha _je^{i\beta }Z_j^{(0)}`$ (107)
$`+{\displaystyle \underset{lj}{}}\alpha _je^{i(\nu _l\nu _j)\tau }e^{i\beta }Z_l^{(0)})]\},`$ (108)
where
$$Z_j^{(0)}(t)=_{\mathrm{}}^{\mathrm{}}_0^{2\pi }e^{i\eta }\rho _j^{(0)}(\eta ,\omega ,t,\nu _j)𝑑\eta 𝑑\omega .$$
(109)
Elimination of secular terms yields the following condition:
$``$ $`{\displaystyle \frac{\rho _j^{(0)}}{t}}+D{\displaystyle \frac{^2\rho _j^{(0)}}{\omega ^2}}\omega ^{}{\displaystyle \frac{\rho _j^{(0)}}{\beta }}`$ (110)
$``$ $`{\displaystyle \frac{1}{m}}{\displaystyle \frac{}{\omega ^{}}}\left\{\rho _j^{(0)}\left[\omega ^{}+K\alpha _j\text{Im}(e^{i\beta }Z_j^{(0)})\right]\right\}=0.`$ (111)
Notice that this equation corresponds to the equation of a unimodal frequency distribution, already studied in . Its solution evolves towards the following stationary state as time evolves:
$`\rho _j^{(0)}(\beta ,\omega ^{})`$ $`=`$ $`\sqrt{{\displaystyle \frac{m}{2\pi D}}}e^{(m/2D)\omega ^2}`$ (112)
$`\times `$ $`{\displaystyle \frac{e^{(K\alpha _j/D)R_j\mathrm{cos}(\mathrm{\Psi }_j\beta )}}{_0^{2\pi }e^{(K\alpha _j/D)R_j\mathrm{cos}(\mathrm{\Psi }_j\beta ^{})}𝑑\beta ^{}}},`$ (113)
where
$`R_je^{i\mathrm{\Psi }_j}={\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle _0^{2\pi }}e^{i\eta }\rho _j^{(0)}(\eta ,\omega )𝑑\eta 𝑑\omega \underset{t\mathrm{}}{lim}Z_j^{(0)}(t).`$ (114)
Incoherence of a component corresponds to $`R_j=0`$. We know that as $`K\alpha _j`$ surpasses $`2D`$, a stable synchronized solution bifurcates from incoherence. In the particular case of a symmetric bimodal frequency distribution, the bifurcation value is $`K=4D`$, independently of the inertia $`m`$. This agrees with the results of the linear stability analysis for $`\mathrm{\Omega }_0\mathrm{}`$. All the results previously obtained for the Kuramoto model can be applied to the present model without any modification . Thus, both a stable standing wave solution (SW) and an unstable travelling wave solution (TW) bifurcate supercritically from incoherence at $`K=4D`$ . In Figure 16 illustrates the comparison between the asymptotic solution (113) and the result of direct numerical simulation for a large enough value of the frequency $`\mathrm{\Omega }_0`$. Finally, Fig. 17 shows the global bifurcation diagram of the bimodal case for the case of positive $`\alpha `$ and $`\beta `$. As explained in the previous Section, the branch of synchronized stationary states bifurcates subcritically from incoherence at $`K^{}=\mathrm{}`$, provided $`\mathrm{\Omega }_0>\mathrm{\Omega }_0^{\mathrm{}}`$. Thus $`K^{}`$ and the end of the SW and TW branches in Fig. 17 should extend to infinity as $`\mathrm{\Omega }_0+\mathrm{}`$.
## V Numerical results
Four different numerical methods have been used. Numerical simulations of the system of Langevin equations (3) were carried out for a large number of oscillators (N=20 000), using an Euler method. A standard finite difference method was used to solve numerically the Fokker-Planck equation (6), or the system of partial differential equations (69). In addition, we have used a simple spectral method, which generalizes the one proposed in . The idea is to solve a set of ordinary differential equations for moments of $`\rho `$:
$`(x_i^j)_k:={\displaystyle _0^{2\pi }}c_i(\theta ,\mathrm{\Omega }_k,t)\mathrm{cos}[j(\psi \theta )]𝑑\theta ,`$ (115)
$`(y_i^j)_k:={\displaystyle _0^{2\pi }}c_i(\theta ,\mathrm{\Omega }_k,t)\mathrm{sin}[j(\psi \theta )]𝑑\theta ,`$ (116)
$`r={\displaystyle _0^{2\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}\rho (\theta ,\omega ,\mathrm{\Omega },t)\mathrm{cos}(\psi \theta )`$ (117)
$`\times g(\mathrm{\Omega })d\mathrm{\Omega }d\omega d\theta .`$ (118)
The coefficient functions $`c_i(\theta ,\mathrm{\Omega },t)`$ are the same as in the parabolic cylinder expansion (62). The integral in (118) will be approximated by a suitable quadrature formula, as
$`r{\displaystyle \underset{q=1}{\overset{Q}{}}}\alpha _q{\displaystyle _0^{2\pi }}c_0(\mathrm{\Omega }_q,\theta ,t)\mathrm{cos}(\psi \theta )𝑑\theta ={\displaystyle \underset{q=1}{\overset{Q}{}}}\alpha _q(x_0^1)_q.`$ (119)
For instance, the Gauss-Laguerre quadrature has been chosen case for the case of a Lorentzian frequency distribution. The system of ordinary differential equations is given by:
$`(\dot{x}_i^j)_k=j\sqrt{i}\sqrt{{\displaystyle \frac{D}{m}}}(y_{i1}^j)_k+\sqrt{i}{\displaystyle \frac{1}{\sqrt{mD}}}\mathrm{\Omega }_k(x_{i1}^j)_k`$ (120)
$`+\sqrt{{\displaystyle \frac{i}{mD}}}{\displaystyle \frac{Kr}{2}}\left[(y_{i1}^{j+1})_k(y_{i1}^{j1})_k\right]{\displaystyle \frac{i}{m}}(x_i^j)_k`$ (121)
$`+\sqrt{i+1}\sqrt{D}mj(y_{i+1}^j)_kj\dot{\psi }(y_i^j)_k,`$ (122)
$`(\dot{y}_i^j)_k=j\sqrt{i}\sqrt{{\displaystyle \frac{D}{m}}}(x_{i1}^j)_k+\sqrt{i}{\displaystyle \frac{1}{\sqrt{mD}}}\mathrm{\Omega }_k(y_{i1}^j)_k`$ (123)
$`+\sqrt{{\displaystyle \frac{i}{mD}}}{\displaystyle \frac{Kr}{2}}\left[(x_{i1}^{j1})_k(x_{i1}^{j+1})_k\right]{\displaystyle \frac{i}{m}}(y_i^j)_k`$ (124)
$`\sqrt{i+1}\sqrt{D}mj(x_{i+1}^j)_k+j\dot{\psi }(x_i^j)_k,`$ (125)
$`i=1,\mathrm{},N,j=1,\mathrm{},M,`$ (126)
$`(\dot{x}_0^j)_k=j\sqrt{{\displaystyle \frac{D}{m}}}(y_1^j)_kj\dot{\psi }(y_0^j)_k,`$ (127)
$`(\dot{y}_0^j)_k=j\sqrt{{\displaystyle \frac{D}{m}}}(x_1^j)_k+j\dot{\psi }(x_0^j)_k,`$ (128)
$`i=0,j=1,\mathrm{},M.`$ (129)
In order to numerically simulate this system, it is necessary to truncate the hierarchy after a reasonable number of modes $`N`$, and $`M`$. These numbers will depend on the inertia $`m`$, the spread, and the coupling strength. They should be chosen large enough, so that the numerical results do not depend on $`N`$ and $`M`$.
Fig. 18 shows the ratio $`c_n^{max}/c_0^{max}`$ as a function of $`n`$, for the stationary solution and for various mass values. $`c_n^{max}`$ is the maximum value of the stationary coefficient $`c_n(\theta ,\mathrm{\Omega }_0=1)`$. It has been obtained from a finite-difference solution of Eq. (6). Note that this ratio increases as inertia does. Thus the truncation approximation (89) ceases to make sense for larger values of inertia. Whether truncating (at higher order) the moment equations (115) and (116) yields good numerical results, is tested in Fig. 19. This figure shows how closely the previous system (with moments of order 4 or 10) resembles the direct solution of the Langevin equations. We notice that the system containing moments of order 10 is rather close to the solution of the Langevin equations, but it does not contain the fluctuations unavoidable in stochastic methods. If we are interested in solving the nonlinear Fokker-Planck equation, the system of moment equations seems a good alternative to solving Langevin equations for a large number of oscillators.
## VI conclusions
We have investigated synchronization properties of a model of globally and nonlinearly coupled phase oscillators, where the effects of white noise, inertia and spread in the natural frequency distribution are all considered. The linear stability of the incoherent solution is rigorously analyzed. Stationary and time-dependent solutions of the standing and travelling wave type are obtained by a variety of perturbative methods. These include finding and approximately solving mode-coupling equations for expansions of the probability distribution in parabolic cylinder functions, finding amplitude equations near bifurcation points, and hierarchy-closure assumptions. Numerical simulations of the different equations and the original model favorably agree with the different perturbative results.
Inertia changes the stability boundaries of the incoherent solution in a non trivial way. In the case of an unimodal Lorentzian frequency distribution, incoherence is stabilized, but the effect of mass completely drops out if the frequency spread vanishes. For a discrete bimodal frequency distribution, both stability boundaries and the character of the transition from incoherence to synchronized states depend on the values of the natural frequency $`\mathrm{\Omega }_0`$ and on inertia. The effect of inertia on the stationary solution is dramatic in some cases. In general, inertia hardens the synchronization transition: it may render subcritical (hard) originally supercritical (soft) transitions (in unimodal Lorentzian frequency distributions), or it increases the region in parameter space where the transition is subcritical (discrete bimodal frequency distributions). Analytical as well as numerical calculations confirm these findings.
## VII Acknowledgments
This work was supported, in part, by the European Union TMR contract ERB FMBX-CT97-0157, by UNESCO under contract UVO–ROSTE 875.629.9, by the Spanish DGES grant PB98-0142-C04-C01, and by the GNFM of the Italian C.N.R. All numerical simulations were conducted at CASPUR–Rome.
## A Integrals of special functions and eigenvalue equation
In (31) there appear two integrals over $`w`$:
$`A={\displaystyle _{\mathrm{}}^+\mathrm{}}e^{\left(\frac{w}{2}i\sqrt{mD}\right)^2}D_p\rho _0^{}𝑑w,`$ (A1)
$`B={\displaystyle _{\mathrm{}}^+\mathrm{}}e^{\left(\frac{w}{2}i\sqrt{mD}\right)^2}D_p𝑑w`$ (A2)
The second integral is directly $`𝒜_p(\sqrt{mD})`$. The first integral equals
$`A={\displaystyle \frac{m}{(2\pi )^{\frac{3}{2}}D}}{\displaystyle _{\mathrm{}}^+\mathrm{}}(wi2\sqrt{mD})e^{\left(\frac{w}{2}i\sqrt{mD}\right)^2}D_p(w)𝑑w,`$ (A3)
because of (28). An equivalent expression is
$`A={\displaystyle \frac{2m}{(2\pi )^{\frac{3}{2}}D}}{\displaystyle _{\mathrm{}}^+\mathrm{}}D_p(w){\displaystyle \frac{}{w}}e^{\left(\frac{w}{2}i\sqrt{mD}\right)^2}𝑑w.`$ (A4)
This can be written as
$`A={\displaystyle \frac{2m}{(2\pi )^{\frac{3}{2}}D}}{\displaystyle _{\mathrm{}}^+\mathrm{}}D_p(w){\displaystyle \frac{}{w}}e^{\frac{1}{4}\left(wi2x\right)^2}𝑑w`$ (A5)
$`={\displaystyle \frac{m}{i(2\pi )^{\frac{3}{2}}D}}{\displaystyle _{\mathrm{}}^+\mathrm{}}D_p(w){\displaystyle \frac{}{x}}e^{\frac{1}{4}\left(wi2x\right)^2}𝑑w`$ (A6)
provided we set $`x=\sqrt{mD}`$ after performing the calculations. Thus we see that
$`A={\displaystyle \frac{m}{i(2\pi )^{\frac{3}{2}}D}}𝒜_p^{}(\sqrt{mD}).`$ (A7)
This expression and the previous one for $`B`$ yield (33).
Let us now obtain $`𝒜_p(x)`$ for integer $`p`$ and $`x>0`$. By using that $`D_p(x)`$ is even for even $`p`$ and odd for odd $`p`$, we can write
$`𝒜_p(x)=e^{x^2}{\displaystyle _{\mathrm{}}^+\mathrm{}}e^{\frac{w^2}{4}+iwx}D_p(w)𝑑w`$ (A8)
$`=e^{x^2}{\displaystyle _{\mathrm{}}^+\mathrm{}}e^{\frac{w^2}{4}}(\mathrm{cos}xw+i\mathrm{sin}xw)D_p(w)𝑑w`$ (A9)
$`=2e^{x^2}{\displaystyle _0^+\mathrm{}}D_p(w)e^{\frac{w^2}{4}}f(xw)𝑑w,`$ (A10)
where $`f(xw)=\mathrm{cos}xw`$ for $`p`$ even and $`f(xw)=i\mathrm{sin}xw`$ for $`p`$ odd. We can find (cf. §7.741 of ):
$`{\displaystyle _0^+\mathrm{}}D_p(w)e^{\frac{w^2}{4}}\mathrm{cos}xwdw=(1)^n\sqrt{{\displaystyle \frac{\pi }{2}}}e^{\frac{x^2}{2}}x^{2n},`$ (A11)
$`{\displaystyle _0^+\mathrm{}}D_p(w)e^{\frac{w^2}{4}}\mathrm{sin}xwdw=(1)^n\sqrt{{\displaystyle \frac{\pi }{2}}}e^{\frac{x^2}{2}}x^{2n+1},`$ (A12)
provided $`x>0`$. From these formulas, we obtain
$`𝒜_{2n}(x)=(1)^n\sqrt{2\pi }e^{\frac{x^2}{2}}x^{2n},`$ (A13)
$`𝒜_{2n+1}(x)=i(1)^n\sqrt{2\pi }e^{\frac{x^2}{2}}x^{2n+1}.`$ (A14)
These two expressions are equivalent to (35).
## B Calculation of the coefficients $`\alpha `$ and $`\beta `$
We will now calculate coefficients $`\alpha `$ and $`\beta `$ of Eq. (75). The idea is simply to expand the right hand side in the definition of the order parameter,
$$r=_0^{2\pi }_{\mathrm{}}^{\mathrm{}}_{\mathrm{}}^{\mathrm{}}\mathrm{cos}(\psi \theta )\rho (\theta ,\omega ,\mathrm{\Omega },t)g(\mathrm{\Omega })𝑑\mathrm{\Omega }𝑑\omega 𝑑\theta .$$
(B1)
as a power series in $`r`$. To this end, we fix $`r`$ and $`\psi `$ in Eq. (6), expand its stationary solution in powers of $`r`$ and insert the result in (B1). As we do not have a closed formula for the stationary solution of (6), we shall first derive a hierarchy of equations for its coefficients in the series in powers of $`r`$. Thus we have
$`\rho (\theta ,\omega ,\mathrm{\Omega };r)={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{G_n(\theta \psi ,\omega ,\mathrm{\Omega })}{n!}}r^n,`$ (B2)
$`r={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{r_n}{n!}}r^n.`$ (B3)
Then
$$r_n=_0^{2\pi }_{\mathrm{}}^{\mathrm{}}_{\mathrm{}}^{\mathrm{}}\mathrm{cos}\theta G_n(\theta ,\omega ,\mathrm{\Omega })g(\mathrm{\Omega })𝑑\mathrm{\Omega }𝑑\omega 𝑑\theta ,$$
(B4)
where we have absorbed $`\psi `$ in the definition of $`G_n`$. Comparison with eq. (75), establishes that $`r_0=r_2=0`$, $`r_1=K\alpha /2D`$, and $`r_3=K^3\beta `$. From (B2) and (6), we obtain
$$\frac{D}{m^2}\frac{^2G_n}{\omega ^2}+\frac{1}{m}\frac{}{\omega }[(\omega \mathrm{\Omega })G_n]\omega \frac{G_n}{\theta }=\frac{n}{m}K\mathrm{sin}\theta \frac{G_{n1}}{\omega },$$
(B5)
with $`G_10`$. Since we are trying to find solutions bifurcating from incoherence, we should have $`G_0=\rho _0(\omega ,\mathrm{\Omega })`$, i.e. the $`\theta `$-independent incoherent solution (9). This directly confirms that $`r_0=0`$.
The unknowns $`G_n(\theta ,\omega ,\mathrm{\Omega })`$ are $`2\pi `$-periodic in $`\theta `$, so that we may Fourier expand them as
$$G_n(\theta ,\omega ,\mathrm{\Omega })=\underset{l=\mathrm{}}{\overset{\mathrm{}}{}}Z_n^l(\omega ,\mathrm{\Omega })e^{il\theta },$$
(B6)
where
$$Z_n^l(\omega ,\mathrm{\Omega })=\frac{1}{2\pi }_0^{2\pi }e^{il\theta }G_n(\theta ,\omega ,\mathrm{\Omega })𝑑\theta .$$
(B7)
Here $`Z_0^l0`$ for $`l0`$ because $`G_0`$ is the $`\theta `$-independent incoherent solution. The condition that the probability density function be real yields $`\overline{G_n}=G_n`$, which in turn implies $`Z_n^l=\overline{Z_n^l}`$. Inserting these expressions into (B4), we find
$$r_n=2\pi _{\mathrm{}}^{\mathrm{}}_{\mathrm{}}^{\mathrm{}}Re(Z_n^1)g(\mathrm{\Omega })𝑑\mathrm{\Omega }𝑑\omega .$$
(B8)
The unknowns $`Z_n^l`$ satisfy the following hierarchy of equations:
$$\frac{D}{m^2}\frac{d^2Z_n^l}{d\omega ^2}+\frac{1}{m}\frac{d}{d\omega }[(\omega \mathrm{\Omega })Z_n^l]il\omega Z_n^l=\frac{n}{m}\frac{K}{2i}\frac{d}{d\omega }[Z_{n1}^{l1}Z_{n1}^{l+1}].$$
(B9)
The normalization condition for $`\rho `$ and the incoherent solution together with (B7) imply
$$_{\mathrm{}}^{\mathrm{}}Z_n^0g(\omega )𝑑\omega =\delta _{n0}.$$
(B10)
In order to obtain $`\alpha `$, and $`\beta `$, we should show that $`r_2=0`$ and calculate $`r_1`$ and $`r_3`$. We do this by means of (B9), with $`n=1,2,3`$.
(a) Case n=1. Eq. (B9) becomes
$`{\displaystyle \frac{D}{m^2}}{\displaystyle \frac{d^2Z_1^1}{d\omega ^2}}+{\displaystyle \frac{1}{m}}{\displaystyle \frac{d}{d\omega }}[(\omega \mathrm{\Omega })Z_1^1]i\omega Z_1^1`$ (B11)
$`={\displaystyle \frac{1}{m}}{\displaystyle \frac{K}{2i}}{\displaystyle \frac{dZ_0^0}{d\omega }},`$ (B12)
where $`Z_0^0=G_0`$ is the incoherent solution (9), and we have used $`Z_0^2=0`$. We can solve Eq. (B12) by means of the solution of Eq. (17), in which we set $`\lambda =0`$ and replace $`b_1(\omega ,\mathrm{\Omega })`$ with $`Z_1^1(\omega ,\mathrm{\Omega })`$, and the right-hand side with $`i\frac{mK}{2D}dZ_0^0/d\omega `$. Then we obtain
$`Z_1^1(\omega ;\mathrm{\Omega })`$ $`=`$ $`{\displaystyle \frac{iK}{2m}}e^{\frac{m}{4D}(\omega \mathrm{\Omega })^2}`$ (B13)
$`\times `$ $`{\displaystyle \underset{p=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{_{\mathrm{}}^{\mathrm{}}e^{\left(\frac{w}{2}i\sqrt{mD}\right)^2}D_p\frac{dZ_0^0}{d\omega }𝑑w}{\sqrt{2\pi }p!\left(\frac{p}{m}+i\mathrm{\Omega }+D\right)}}D_p(w),`$ (B14)
The previous analogy allows us to calculate $`r_1`$ from (B8) and Equations (31), (34) and (35) (with $`\lambda =0`$). The result is
$`r_1={\displaystyle \frac{KD}{2}}e^{mD}[{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{1}{D^2+\mathrm{\Omega }^2}}g(\mathrm{\Omega })d\mathrm{\Omega }`$ (B15)
$`+{\displaystyle \underset{p=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(mD)^p(1+\frac{p}{mD})^2}{p!}}`$ (B16)
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{1}{(D+\frac{p}{m})^2+\mathrm{\Omega }^2}}g(\mathrm{\Omega })d\mathrm{\Omega }].`$ (B17)
(b) Case n=2. We now show that $`r_2=0`$.
In order to analyze $`r_2`$, it is necessary to consider $`Z_2^1`$, which is the solution to the following equation:
$$\frac{D}{m^2}\frac{d^2Z_2^1}{d\omega ^2}+\frac{1}{m}\frac{d}{d\omega }[(\omega \mathrm{\Omega })Z_2^1]i\omega Z_2^1=\frac{2}{m}\frac{K}{2i}\frac{d}{d\omega }[Z_1^0Z_1^2].$$
(B18)
We shall show that $`Z_1^0=Z_1^2=0`$. If this is so, the resulting homogeneous equation can be transformed to the parabolic cylinder equation with a quadratic potential having complex coefficients. Its solution cannot decay to zero both as $`\omega \pm \mathrm{}`$ unless it is identically zero. The reason $`Z_1^0`$ and $`Z_1^2`$ are zero is similar. The equation for $`Z_1^0`$ is homogeneous ($`Z_0^l=0`$ for $`l0`$), and it has a solution $`C(\mathrm{\Omega })e^{m(\omega \mathrm{\Omega })^2/(2D)}`$. The normalization condition (B10) then implies $`C0`$. $`Z_1^2`$ again obeys a homogeneous equation which can be transformed into the parabolic cylinder equation with a complex potential. The only solution which decays to zero as $`\omega \pm \mathrm{}`$ is again $`Z_1^2=0`$.
(c) Case n=3.
The equations for $`Z_3^1`$ are
$`{\displaystyle \frac{D}{m^2}}{\displaystyle \frac{d^2Z_3^1}{d\omega ^2}}+{\displaystyle \frac{1}{m}}{\displaystyle \frac{d}{d\omega }}[(\omega \mathrm{\Omega })Z_3^1]i\omega Z_3^1={\displaystyle \frac{3}{m}}{\displaystyle \frac{K}{2i}}{\displaystyle \frac{d}{d\omega }}[Z_2^0Z_2^2],`$ (B19)
$`{\displaystyle \frac{D}{m^2}}{\displaystyle \frac{d^2Z_2^2}{d\omega ^2}}+{\displaystyle \frac{1}{m}}{\displaystyle \frac{d}{d\omega }}[(\omega \mathrm{\Omega })Z_2^2]2i\omega Z_2^2={\displaystyle \frac{2}{m}}{\displaystyle \frac{K}{2i}}{\displaystyle \frac{d}{d\omega }}Z_1^1,`$ (B20)
$`{\displaystyle \frac{D}{m^2}}{\displaystyle \frac{d^2Z_2^0}{d\omega ^2}}+{\displaystyle \frac{1}{m}}{\displaystyle \frac{d}{d\omega }}[(\omega \mathrm{\Omega })Z_2^0]={\displaystyle \frac{2}{m}}{\displaystyle \frac{K}{2i}}{\displaystyle \frac{d}{d\omega }}[Z_1^1Z_1^1]={\displaystyle \frac{2K}{m}}Im{\displaystyle \frac{d}{d\omega }}Z_1^1,`$ (B21)
$`{\displaystyle \frac{D}{m^2}}{\displaystyle \frac{d^2Z_1^1}{d\omega ^2}}+{\displaystyle \frac{1}{m}}{\displaystyle \frac{d}{d\omega }}[(\omega \mathrm{\Omega })Z_1^1]i\omega Z_1^1={\displaystyle \frac{1}{m}}{\displaystyle \frac{K}{2i}}{\displaystyle \frac{d}{d\omega }}Z_0^0.`$ (B22)
This system of equations has to be solved numerically in order to obtain the coefficient $`\beta `$.
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# CECS-PHY-00/01ULB-TH-00/01 Black Hole Scan
## I Introduction
Black holes are much more than a particular class of exact solutions of the Einstein Equations; they are an essential feature of the spacetime dynamics in almost any sensible theory of gravity. Within the framework of General Relativity, the singularity theorems of Hawking and Penrose show that singular configurations –such as the Schwarzschild black hole– are inevitable under quite generic initial conditions. Furthermore, the Schwarzschild solution describes the leading asymptotic behavior of the geometry for any localized distribution of matter. The existence of this solution at spacelike infinity is a central ingredient to prove the positivity of energy in General Relativity . On the other hand, black holes are also fundamental objects where the thermodynamics of the gravitational field and its connection with information theory is expected to shed light on the quantization problem.
In this paper, we survey the black hole solutions in a class of gravitation theories, selected by requiring that they have a unique anti-de Sitter vacuum with a fixed cosmological constant. For a given dimension $`d`$, the Lagrangians under consideration are labeled by an integer $`k=1,2,\mathrm{},\left[\frac{d1}{2}\right]`$, where the Einstein-Hilbert Lagrangian corresponds to $`k=1`$. For each of these theories we examine their static, spherically symmetric solutions. The existence of physical black holes is then used as a criterion to assess the validity of those theories, leading to a natural splitting between theories with even and odd $`k`$.
Coupling these gravity theories with the Maxwell action predicts the smallest size of a spherically symmetric electrically charged source, except for $`k=1`$.
An important aspect of the black holes under consideration is their thermodynamics, which is expected to be a reflection of the underlying quantum theory. The canonical ensemble for minisuperspaces containing the black holes found in these theories is well defined provided a negative cosmological constant exists. It is found that black holes are unstable against decay by Hawking radiation, unless their horizon radius is large, compared to the AdS radius.
Among all theories under consideration, there is only one representative in each odd dimension, given by a Chern-Simons action, having physical black holes whose spectrum has a mass gap separating them from AdS spacetime. These black holes always reach thermal equilibrium with a heat bath, and have positive specific heat, which guarantees their stability under thermal fluctuations.
### A Higher Dimensional Gravity Revisited
The standard higher dimensional extension of the four-dimensional Einstein-Hilbert (EH) action reads
$$I_{EH}=\frac{1}{2(d2)\mathrm{\Omega }_{d2}G}d^dx\sqrt{g}(R2\mathrm{\Lambda }).$$
(1)
String and $`M`$-theory corrections to this action would bring in higher powers of curvature –see, e.g. Refs. . This may be a source of inconsistencies because higher powers of curvature could give rise to fourth order differential equations for the metric. This not only complicates the causal evolution, but in general would introduce ghosts and violate unitarity. However, Zwiebach and Zumino observed that ghosts are avoided if stringy corrections would only consist of the dimensional continuations of the Euler densities, so that the resulting field equations remain second order.
These theories are far from exotic. Indeed, they are described by the most general Lagrangians constructed with the same principles as General Relativity, that is, general covariance and second order field equations for the metric. These theories were first discussed by Lanczos for $`d=5`$ in $`1938`$ and more recently by Lovelock for $`d3`$ .
The Lanczos-Lovelock (LL) action is a polynomial of degree $`[d/2]`$ in curvature<sup>*</sup><sup>*</sup>*Here $`[x]`$ is the integer part of $`x`$., which can also be written in terms of the Riemann curvature $`R^{ab}=d\omega ^{ab}+\omega _c^a\omega ^{cb}`$ and the vielbein $`e^a`$ asWedge product between forms is understood throughout.
$$I_G=\kappa \underset{p=0}{\overset{[d/2]}{}}\alpha _pL^{(p)},$$
(2)
where $`\alpha _p`$ are arbitrary constants, and $`L^{(p)}`$ is given by
$$L^{(p)}=ϵ_{a_1\mathrm{}a_d}R^{a_1a_2}R^{a_{2p1}a_{2p}}e^{a_{2p+1}}e^{a_d}.$$
(3)
In first order formalism the action (2) is regarded as a functional of the vielbein and the spin connection, and the corresponding field equations obtained varying with respect to $`e^a`$ and $`\omega ^{ab}`$read
$`{\displaystyle \underset{p=0}{\overset{[\frac{d1}{2}]}{}}}\alpha _p(d2p)_a^p`$ $`=`$ $`0,`$ (4)
$`{\displaystyle \underset{p=1}{\overset{[\frac{d1}{2}]}{}}}\alpha _pp(d2p)_{ab}^p`$ $`=`$ $`0,`$ (5)
where we have defined
$`_a^p:=`$ $`ϵ_{ab_1\mathrm{}b_{d1}}R^{b_1b_2}\mathrm{}R^{b_{2p1}b_{2p}}e^{b_{2p+1}}\mathrm{}e^{b_{d1}},`$
$`_{ab}^p:=`$ $`ϵ_{aba_3\mathrm{}a_d}R^{a_3a_4}\mathrm{}R^{a_{2p1}a_{2p}}T^{a_{2p+1}}e^{a_{2p+2}}\mathrm{}e^{a_d}.`$
Here $`T^a=de^a+\omega _b^ae^b`$ is the torsion $`2`$-form.
Note that in even dimensions, the term $`L^{(d/2)}`$ is the Euler density and therefore does not contribute to the field equations. However, the presence of this term in the action –with a fixed weight factor– guarantees the existence of a well defined variational principle for asymptotically locally AdS spacetimes . Moreover, the Euler density should assign different weights to non-homeomorphic geometries in the quantum theory.
The first two terms in the LL action (2) are the cosmological and kinetic terms of the EH action (1) respectively, and therefore General Relativity is contained in the LL theory as a particular case.
The linearized approximation of the LL and EH actions around a flat, torsionless background are classically equivalent . However, beyond perturbation theory the presence of higher powers of curvature in the Lagrangian makes both theories radically different. In particular, black holes and big-bang solutions of (2), have different asymptotic behaviors from their EH counterparts in general. Hence, a generic solution of the LL action cannot be approximated by a solution of Einstein’s theory.
### B Drawbacks
For a given dimension and an arbitrary choice of coefficients $`\alpha _p`$’s, higher dimensional LL theories have some drawbacks. One difficulty is the fact that the dynamical evolution can become unpredictable because the Hessian matrix cannot be inverted for a generic field configuration. Thus, the velocities are multivalued functions of the momenta and therefore the passage from the Lagrangian to the Hamiltonian is ill defined .
A reflection of this problem can be viewed in the static, spherically symmetric solutions of (4) and (5). For arbitrary $`\alpha _p`$’s there are negative energy solutions with horizons and positive energy solutions with naked singularities .
These problems can be curbed if the coefficients $`\alpha _p`$’s are chosen in a suitable way. The aim of the next section is to show that requiring the theories to possess a unique cosmological constant, strongly restricts the coefficients $`\alpha _p`$’s. As a consequence, one obtains a set of theories labelled by an integer $`k`$ which lead to well defined black hole configurations.
## II Selecting Sensible Theories
The field equations of LL theory (4) can be rearranged as a polynomial of $`k`$th degree in the curvature
$$ϵ_{ab_1\mathrm{}b_{d1}}\beta _0\overline{R}_{\beta _1}^{b_1b_2}\mathrm{}\overline{R}_{\beta _k}^{b_{2k1}b_{2k}}e^{b_{2k+1}}\mathrm{}e^{b_{d1}}=0$$
(6)
where $`\overline{R}_{\beta _i}^{ab}:=R^{ab}+\beta _ie^ae^b`$, and the coefficients $`\beta _i`$’s are related to the $`\alpha _p`$’s through
$$\underset{p}{\overset{[\frac{d1}{2}]}{}}(d2p)\alpha _px^p=\beta _0\underset{i}{\overset{k}{}}(x\beta _i).$$
(7)
Equation (6) can possess in general, several constant curvature solutions with different radii $`r_i=|\beta _i|^{1/2}`$, making the value of the cosmological constant ambiguous. In fact, the cosmological constant could change in different regions of a spatial section, or it could jump arbitrarily as the system evolves in time .
On the other hand, solving (6) for a given global isometry leads in general to several solutions with different asymptotic behaviors. Some of these solutions are “spurious” in the sense that perturbations around them yield ghosts. For instance, if $`\alpha _1`$ and $`\alpha _2`$ were the only nonvanishing coefficients in the LL action (3), two different static, spherically symmetric solutions would be obtained, which are asymptotically (A)dS and flat respectively. The perturbations around the latter solution are gravitons, while those on the former are spurious in the sense described above .
These problems are overcome demanding the theory to have a unique cosmological constant.
Requiring the existence of a unique cosmological constant implies that locally maximally symmetric solutions possess only one fixed radius, that is $`R^{ab}=\beta e^ae^b`$. This in turn means that the polynomial (7) must have only one real root. Hence, the coefficients $`\alpha _p`$’s are fixed through equation (7), so that the real $`\beta `$’s in (6) are all equal, allowing –for $`d7`$– an arbitrary number of distinct imaginary $`\beta `$’s which must come in conjugate pairs. Under this assumption, solutions representing localized sources of matter approach a constant curvature spacetime with a fixed radius in the asymptotic region.
In what follows, we consider the simplest class of such theories, namely, we assume the field equations to be of the form (6) with only one real $`\beta :=\frac{1}{l^2}`$, and no complex rootsA negative cosmological constant is assumed for later convenience, but this analysis does not depend on its sign.. These theories are described by the action
$$I_k=\kappa \underset{p=0}{\overset{k}{}}c_p^kL^{(p)},$$
(8)
which is obtained from (2) with the choice
$$\alpha _p:=c_p^k=\{\begin{array}{cc}\frac{l^{2(pk)}}{(d2p)}\left(\begin{array}{c}k\\ p\end{array}\right)\hfill & ,\text{ }pk\hfill \\ 0\hfill & ,\text{ }p>k\hfill \end{array}$$
(9)
where $`1k[\frac{d1}{2}]`$.
For a given dimension $`d`$, the coefficients $`c_p^k`$ give rise to a family of inequivalent theories, labeled by the integer $`k\{1,\mathrm{},[\frac{d1}{2}]\}`$ which represents the highest power of curvature in the Lagrangian. This set of theories possesses only two fundamental constants, $`\kappa `$ and $`l`$, related to the gravitational constant $`G_k`$ and the cosmological constant $`\mathrm{\Lambda }`$ through<sup>§</sup><sup>§</sup>§Here the gravitational constant has natural units given by $`[G_k]=(`$length$`)^{d2k}`$.
$`\kappa `$ $`=`$ $`{\displaystyle \frac{1}{2(d2)!\mathrm{\Omega }_{d2}G_k}},`$ (10)
$`\mathrm{\Lambda }`$ $`=`$ $`{\displaystyle \frac{(d1)(d2)}{2l^2}}.`$ (11)
The field equations for the action $`I_k`$ in (8), read
$`ϵ_{ba_1\mathrm{}a_{d1}}\overline{R}^{a_1a_2}\overline{R}^{a_{2k1}a_{2k}}e^{a_{2k+1}}e^{a_{d1}}`$ $`=`$ $`0,`$ (12)
$`ϵ_{aba_3\mathrm{}a_d}\overline{R}^{a_3a_4}\overline{R}^{a_{2k1}a_{2k}}T^{a_{2k+1}}e^{a_{2k+2}}e^{a_{d1}}`$ $`=`$ $`0,`$ (13)
with $`\overline{R}^{ab}:=R^{ab}+\frac{1}{l^2}e^ae^b`$.
### A Examples
There are special cases of interest which are obtained for particular values of the integer $`k`$.
* The Einstein-Hilbert action in $`d`$ dimensions (1) is recovered setting $`k=1`$ in (8).
* At the other end of the range, $`k=[\frac{d1}{2}]`$, even and odd dimensions must be distinguished. These cases are exceptional in that they are the only ones which allow sectors with non-trivial torsion , as discussed in Appendix A. When $`d=2n1`$, the maximum value of $`k`$ is $`n1`$, and the corresponding Lagrangian is a Chern-Simons (CS) $`2n1`$-form defined through (()VII A). For $`d=2n`$ and $`k=n1`$, the action can be written as the Pfaffian of the $`2`$-form $`\overline{R}^{ab}=R^{ab}+\frac{1}{l^2}e^ae^b`$ and, in this sense, it has a Born-Infeld-like (BI-like) form given by (95)Strictly speaking one must add the Euler density to the Lagrangian in (8) with the coefficient $`\alpha _n=c_n^{n1}:=\frac{l^2}{2n}`$, which does not modify the field equations. Therefore, the same BI Lagrangian (95) is recovered from (9) but now the index $`p`$ ranges from $`0`$ to $`n`$..
* In three and four dimensions equation (9) defines only one possible theory which corresponds to EH. As is well known, the EH action is equivalent to CS theory in three dimensions , and for $`d=4`$ the EH action coincides with the BI action up to the Euler density.
* In five and six dimensions, there are only two inequivalent theories which correspond to $`k=1,2`$. In five dimensions, $`k=1`$ represents EH and $`k=2`$ leads to CS. For $`d=6`$, one obtain EH and BI respectively.
* For $`d7`$ there exist other interesting possibilities which are neither EH, BI nor CS. For instance, consider the theory given by the action $`I_k`$ in (8) with $`k=2`$, which exists only for dimensions greater than 4. In this case the Lagrangian reads
$$L=\kappa \left(\frac{l^4}{d}L^{(0)}+\frac{2l^2}{d2}L^{(1)}+\frac{1}{d4}L^{(2)}\right),$$
(14)
with
$`L^{(0)}`$ $`=`$ $`ϵ_{a_1\mathrm{}a_d}e^{a_1}e^{a_d},`$ (15)
$`L^{(1)}`$ $`=`$ $`ϵ_{a_1\mathrm{}a_d}R^{a_1a_2}e^{a_3}e^{a_d},`$ (16)
$`L^{(2)}`$ $`=`$ $`ϵ_{a_1\mathrm{}a_d}R^{a_1a_2}R^{a_3a_4}e^{a_5}e^{a_d}.`$ (17)
Here $`L^{(0)}`$ and $`L^{(1)}`$ are proportional to the standard cosmological and kinetic terms for the EH action, and $`L^{(2)}`$ is proportional to the four dimensional Gauss-Bonnet density ,
$$^2:=(R_{\mu \nu \alpha \beta }R^{\mu \nu \alpha \beta }4R_{\mu \nu }R^{\mu \nu }+R^2),$$
(18)
where $`R^{\mu \nu \alpha \beta }`$, $`R^{\mu \nu }`$ and $`R`$ are the Riemann, Ricci and scalar curvatures, respectively. The action in standard tensor components reads
$$I_2=\frac{2(d3)!\kappa }{l^2}\underset{M}{}d^dx\sqrt{g}\left[\frac{l^2^2}{2(d3)(d4)}+R\mathrm{\Lambda }\right],$$
(19)
with $`\mathrm{\Lambda }`$ given by (11). In sum, the theory with $`k=2`$ is described by a Lagrangian which is a linear combination of Gauss-Bonnet density, the EH Lagrangian and the volume term with fixed weights.
Each of the theories described by $`I_k`$ for all $`k`$ possesses a unique cosmological constant. In fact, as is apparent from equations (12) and (13), spacetimes satisfying $`\overline{R}^{ab}=0`$ are the only locally maximally symmetric solutions. This ensures that localized matter fields give rise to solutions which are asymptotically AdS spacetimes.
## III Static and Spherically Symmetric Solutions
In this section, we test the theories described by $`I_k`$ analyzing their static, spherically symmetric solutions including their electrically charged extensions. It is shown that they possess well behaved black holes, resembling the Schwarzschild-AdS and Reissner-Nordstrom-AdS solutions. The subset of theories with $`k`$ odd differ from their even counterparts, because in the first case there is a unique black hole solution, whereas in the latter, an additional solution with a naked singularity exists.
### A Pure Gravity
Consider static and spherically symmetric solutions of equations (12) and (13) for a fixed value of the label $`k`$. In Schwarzschild-like coordinates, the metric can be written as
$$ds^2=N^2(r)f^2(r)dt^2+\frac{dr^2}{f^2(r)}+r^2d\mathrm{\Omega }_{d2}^2.$$
(20)
Replacing this ansatz in the field equations (12) and (13) leads to the following equations for $`N`$ and $`f^2`$
$`{\displaystyle \frac{dN}{dr}}`$ $`=`$ $`0,`$ (21)
$`{\displaystyle \frac{d}{dr}}\left(r^{d1}\left[F(r)+{\displaystyle \frac{1}{l^2}}\right]^k\right)`$ $`=`$ $`0,`$ (22)
where the function $`F(r)`$ is given by
$$F(r)=\frac{1f^2(r)}{r^2}.$$
(23)
Integrating equations (22) yields
$`N`$ $`=`$ $`N_{\mathrm{}},`$ (24)
$`f^2(r)`$ $`=`$ $`1+{\displaystyle \frac{r^2}{l^2}}\sigma \left({\displaystyle \frac{C_1}{r^{d2k1}}}\right)^{1/k},`$ (25)
where the integration constant $`N_{\mathrm{}}`$ relates coordinate time to the proper time of an observer at spatial infinity and in what follows is chosen equal to one. Here $`\sigma =(\pm 1)^{(k+1)}`$, and the integration constant $`C_1`$ is identified as
$`C_1=2G_k(M+C_0),`$
where $`M`$ stands for the mass, as is discussed in detail in section III.C.
For even $`k`$, the ambiguity of sign expressed through $`\sigma `$ in (25) implies that there are two possible solutions provided $`C_1>0`$. The solution with $`\sigma =1`$ describes a real black hole with a unique event horizon surrounding the singularity at the origin. The solution with $`\sigma =1`$ has a naked singularity with positive mass.
If $`k`$ is odd, there is no ambiguity of sign because $`\sigma `$ cannot be different from unity, therefore in that case there exists a unique static, spherically symmetric solution, which corresponds to a black hole with positive mass.
The black hole mass for any value of $`k`$ is a monotonically increasing function of the horizon radius $`r_+`$, which reads
$$M(r_+)=\frac{r_+^{d2k1}}{2G_k}\left(1+\frac{r_+^2}{l^2}\right)^kC_0.$$
(26)
The additive constant $`C_0`$ is chosen so that the horizon shrinks to a point for $`M0`$, hence
$$C_0=\frac{1}{2G_k}\delta _{d2k,1},$$
(27)
which vanishes in all cases except for CS theory.
Summarizing, for a given dimension $`d3`$ the full set of $`[\frac{d1}{2}]`$ inequivalent theories given by the action $`I_k`$ in (8), possess asymptotically AdS black hole solutions whose line elements read
$`ds^2`$ $`=`$ $`\left(1+{\displaystyle \frac{r^2}{l^2}}\left({\displaystyle \frac{2G_kM+\delta _{d2k,1}}{r^{d2k1}}}\right)^{1/k}\right)dt^2+`$ (29)
$`{\displaystyle \frac{dr^2}{1+\frac{r^2}{l^2}\left(\frac{2G_kM+\delta _{d2k,1}}{r^{d2k1}}\right)^{1/k}}}+r^2d\mathrm{\Omega }_{d2}^2.`$
One can see from (29) that for $`k=1`$, the three dimensional black hole and Schwarzschild-AdS solutions of the $`d`$-dimensional Einstein-Hilbert action with negative cosmological constant are recovered. The black hole solutions corresponding to BI and CS theories are obtained also from (29) setting $`k=[\frac{d1}{2}]`$.
The whole set of black hole metrics given by (29) share a common causal structure when $`M>0`$, which coincides with the familiar one described by the Penrose diagram of the four dimensional Schwarzschild-AdS solution. Nevertheless, the presence of the Kronecker delta within the metrics (29) signals the existence of two possible black hole vacua ($`M=0`$) with different causal structures. The generic case holds for the whole set of theories except CS, whose line elements are described by (29) with $`d2k1`$, that is
$`ds^2`$ $`=`$ $`\left(1+{\displaystyle \frac{r^2}{l^2}}\left({\displaystyle \frac{2G_kM}{r^{d2k1}}}\right)^{1/k}\right)dt^2+`$ (31)
$`{\displaystyle \frac{dr^2}{1+\frac{r^2}{l^2}\left(\frac{2G_kM}{r^{d2k1}}\right)^{1/k}}}+r^2d\mathrm{\Omega }_{d2}^2.`$
Analogously with the Schwarzschild-AdS metric, this set possesses a continuous mass spectrum, whose vacuum state is the AdS spacetime. The other case is obtained only for $`d=2n1`$ dimensions, and it is a peculiarity of CS theories, whose black hole solutions are recovered from (29) with $`k=n1`$, which read
$`ds^2`$ $`=`$ $`\left(1+{\displaystyle \frac{r^2}{l^2}}\left(2G_{n1}M+1\right)^{\frac{1}{n1}}\right)dt^2+`$ (33)
$`{\displaystyle \frac{dr^2}{1+\frac{r^2}{l^2}\left(2G_{n1}M+1\right)^{\frac{1}{n1}}}}+r^2d\mathrm{\Omega }_{d2}^2.`$
In that case, the black hole vacuum ($`M=0`$) differs from AdS spacetime. Although this configuration has no constant curvature for $`d>3`$, it possesses the same causal structure as the three-dimensional zero mass black hole. Another common feature with $`2+1`$ dimensions is the existence of a mass gap between the zero mass black hole and AdS spacetime, where the later is obtained for $`M=\frac{1}{2G_{n1}}`$.
### B Coupling to the Electromagnetic Field
The standard coupling with the electromagnetic field is obtained adding to the gravitational action $`I_k`$ in Eq. (8) the Maxwell termThe constant $`ϵ`$ is related with the “vacuum permeability” through $`ϵ=\frac{1}{\mathrm{\Omega }_{d2}ϵ_0}`$. Its natural units are $`[ϵ]=(`$length$`)^{d4}`$.
$$I_M=\frac{1}{4ϵ\mathrm{\Omega }_{d2}}\sqrt{g}F^{\mu \nu }F_{\mu \nu }\text{ }d^dx.$$
(34)
Electrically charged solutions which are static and spherically symmetric can be found through the ansatz (20), and requiring that and the only non vanishing component of the electromagnetic field strength be
$$F_{0r}=_rA_0(r).$$
(35)
The field equations for $`N`$, $`f^2`$and $`A_0`$ read
$`{\displaystyle \frac{dN}{dr}}`$ $`=`$ $`0,`$ (36)
$`{\displaystyle \frac{d}{dr}}(r^{d2}p)`$ $`=`$ $`0,`$ (37)
$`{\displaystyle \frac{dA_0}{dr}}+Np`$ $`=`$ $`0,`$ (38)
$`{\displaystyle \frac{d}{dr}}\left(r^{d1}\left[F(r)+{\displaystyle \frac{1}{l^2}}\right]^k\right)`$ $`=`$ $`{\displaystyle \frac{G_k}{ϵ}}r^{d2}p^2,`$ (39)
where $`F(r)`$ is defined in equation (23), and $`p(r)`$ is a redefinition of the electric field,
$$p=\frac{1}{N}F_{0r}.$$
(40)
Integrating these equations yields
$`N`$ $`=`$ $`N_{\mathrm{}}=1,`$ (41)
$`p(r)`$ $`=`$ $`ϵ{\displaystyle \frac{Q}{r^{d2}}},`$ (42)
$`A_0(r)`$ $`=`$ $`\varphi _{\mathrm{}}+{\displaystyle \frac{ϵ}{(d3)}}{\displaystyle \frac{Q}{r^{d3}}}\text{ },`$ (43)
$`f^2(r)`$ $`=`$ $`1+{\displaystyle \frac{r^2}{l^2}}\sigma g_k(r),`$ (44)
with $`\sigma =(\pm 1)^{(k+1)}`$ and
$$g_k(r)=\left(\frac{2G_kM+\delta _{d2k,1}}{r^{d2k1}}\frac{ϵG_k}{\left(d3\right)}\frac{Q^2}{r^{2(dk2)}}\right)^{\frac{1}{k}}.$$
(45)
The integration constants $`M`$ and $`Q`$ in (45) are the mass and the electric charge of the black hole respectively, as is shown in the next subsection.
Equations (44) provide the electrically charged extension of the vacuum solution (25)<sup>\**</sup><sup>\**</sup>\**The expression (45) is valid for $`d>3`$. The three dimensional case is discussed in Refs. .. The presence of $`\sigma `$ in (44) leads to a similar picture as in the uncharged case. When $`k`$ is odd, there is a unique electrically charged black hole solution because $`\sigma `$ is always equal to one, but when $`k`$ is even, the solution with $`\sigma =1`$ represents a black hole, and the solution with $`\sigma =1`$ possess a naked singularity.
Therefore, electrically charged asymptotically AdS black hole solutions are obtained from (44) with $`\sigma =1`$, whose line element read –for $`d>3`$– as
$`ds^2`$ $`=`$ $`\left(1+{\displaystyle \frac{r^2}{l^2}}g_k(r)\right)dt^2+`$ (47)
$`{\displaystyle \frac{dr^2}{1+\frac{r^2}{l^2}g_k(r)}}+r^2d\mathrm{\Omega }_{d2}^2,`$
where $`g_k(r)`$ is given by (45).
As is naturally expected, the set of black holes described by (47), reduce to the $`d`$-dimensional Reissner-Nordstrom-AdS solution for $`k=1`$. The electrically charged black hole solutions corresponding to BI and CS theories are also recovered for $`d=2n`$ and $`d=2n1`$ respectively, as it can be seen replacing $`k=n1`$ in (47).
For a generic value of the label $`k`$, in analogy with standard Reissner-Nordström-AdS geometry, the black hole solutions given by (47) possess in general two horizons located at the roots of $`f^2(r)`$. They satisfy $`0<r_{}<r_+`$ provided the mass is bounded from below as $`Mh_k(Q)`$, where $`h_k`$ is a monotonically increasing function of the electric charge. Both horizons merge when the bound is saturated, corresponding to the extreme case, that is $`r_+=r_{}`$ for $`M=h_k(Q)`$. Solutions with $`M<h_k(Q)`$ possess naked singularities which should be considered unphysical. Thus, for a given electric charge, the existence of a lower bound on $`M`$ is in agreement with the cosmic censorship principle.
An important difference with the Reissner-Nordström-AdS case ($`k=1`$) is shared by all electrically charged black hole solutions with $`k1`$, as can be inferred evaluating the scalar curvature for the metrics (47), given by
$$R=\frac{1}{r^{d2}}\frac{d^2}{dr^2}\left[r^{d2}\left(g_k(r)\frac{r^2}{l^2}\right)\right].$$
(48)
For any $`k1`$, equation (48) has a branch point unbounded singularity at the zero of the function $`g_k(r)`$. This is a real timelike singularity located at
$$r_e=\left(\frac{ϵ}{2(d3)}\frac{Q^2}{(M+\frac{1}{2G_k}\delta _{d2k,1})}\right)^{\frac{1}{d3}},$$
(49)
which can be reached in a finite proper time. However, an external observer is protected from it because it is surrounded by both horizons, i.e. $`0<r_e<r_{}<r_+`$.
When $`k`$ is even, spacetime cannot be extended to $`r<r_e`$, because in that case the metric (47) would become complex. This means that the manifold possesses a real boundary at $`r=r_e`$, and therefore, $`r_e`$ is the smallest possible size of a spherical body endowed of electric charge $`Q`$ and mass $`M`$.
For odd values of $`k1`$ there is no obstruction to define spacetime within the region $`r<r_e`$. However, as it can be seen from (48), there is an additional timelike singularity located at $`r=0`$. In that case, a spherical source with electric charge $`Q`$ and mass $`M`$, whose radius is smaller than $`r_e`$ possesses an exterior geometry described by (47) which cannot be empty, since it has a singularity at $`r=r_e`$. This means that the original source generates “a shield”, which acts as the effective source of the external geometry. Hence again, $`r_e`$ is the smallest size for the source.
This means that the presence of electric charge brings in a new length scale into the system, except when one deals with the EH action. For CS theory ($`d=2k+1`$), the radius $`r_e`$ depends on the gravitational constant. However, in the generic case, which is given by the set of theories which are neither EH or CS, the radius $`r_e`$ depends only on intrinsic features of the source and it is completely independent from gravity. That is, $`r_e`$ is independent of the label $`k`$, the gravitational constant $`G_k`$ and the cosmological constant – or equivalently the AdS radius $`l`$ –, that is
$$r_e=\left(\frac{ϵ}{2(d3)}\frac{Q^2}{M}\right)^{\frac{1}{d3}},$$
(50)
which has the same expression as the classical radius of the electron in $`d`$ dimensions. It is noteworthy that $`r_e`$ is encoded in the geometry.
Remarkably, the only theory within the family discussed here, which is unable to predict a minimum size for the source is General Relativity.
### C Mass and Electric Charge from Boundary Terms
In order to identify the integration constants appearing in the black hole solutions (29) and (47) with the mass and electric charge, it is convenient to carry out the canonical analysis . The total action can be written in Hamiltonian form as
$$I_T=I_G+I_M+B,$$
(51)
where $`I_G`$ and $`I_M`$ are the canonical actions for gravity and electromagnetism, respectively
$`I_G`$ $`=`$ $`{\displaystyle d^dx(\pi ^{ij}\dot{g}_{ij}N^{}H_GN^iH_{Gi})},`$ (52)
$`I_M`$ $`=`$ $`{\displaystyle d^dx(p^i\dot{A}_iN^{}H_MN^iH_{Mi}A_0_ip^i)},`$ (53)
and $`B`$ stands for a boundary term which is needed so that the action attains an extremum on the classical solution. Here $`H_{G\mu }`$ and $`H_{M\mu }`$ are the Hamiltonian generators of diffeomorphisms on the gravitational and electromagnetic phase spaces, respectively (see Ref. ).
In case of static, spherically symmetric spacetimes, a general theorem implies that the extremum of the action can be found through a minisuperspace model, which is obtained replacing the Ansätze (20) and (35) into the action, as well. Hence, one deals with a simple one-dimensional model which allows fixing the boundary term $`B`$as a function of the integration constants requiring the total action (51) to have an extremum on the classical solutions. The minisuperspace action takes the form
$`I_T`$ $`=`$ $`\mathrm{\Delta }t{\displaystyle \frac{N}{2}\left[\frac{d}{dr}\left\{\frac{r^{d1}}{G_k}\left[F(r)+\frac{1}{l^2}\right]^k\right\}\frac{1}{ϵ}r^{d2}p^2\right]𝑑r}`$ (55)
$`+{\displaystyle \frac{1}{ϵ}}\mathrm{\Delta }t{\displaystyle A_0\frac{d}{dr}\left(r^{d2}p\right)𝑑r}+B,`$
where $`N:=N^{}(r)f^2(r)`$, and $`p`$ is a redefinition of the canonical momentum $`p^r`$, conjugate to $`A_r`$,
$$p=\frac{1}{N}F_{0r}=\frac{ϵ\mathrm{\Omega }_{d2}}{r^{d2}\sqrt{\gamma }}p^r,$$
(56)
and $`\gamma `$ is the determinant of the angular metric.
The action (55) is a functional of the fields $`N`$, $`f^2`$, $`A_0`$ and $`p`$, whose variation leads to a bulk term which vanishes on the field equations (39). Thus, the variation of the action (55) on shell is a boundary term given by
$`\delta I_T`$ $`=`$ $`\mathrm{\Delta }t{\displaystyle \frac{d}{dr}\left(N\frac{r^{d1}}{2G_k}\delta \left[F(r)+\frac{1}{l^2}\right]^k\right)𝑑r}`$ (58)
$`+{\displaystyle \frac{1}{ϵ}}\mathrm{\Delta }t{\displaystyle \frac{d}{dr}\left(A_0r^{d2}\delta p\right)𝑑r}+\delta B,`$
which means that the action is stationary on the black hole solution provided
$$\delta B=\mathrm{\Delta }t(N_{\mathrm{}}\delta M+\varphi _{\mathrm{}}\delta Q).$$
(59)
Since $`\delta M`$ is multiplied by the proper time separation at infinity, one identifies $`M`$ and $`Q`$ as the mass and the electric charge up to additive constants. The additive constant related with the mass is called $`C_0`$ and it is fixed in (27), requiring that the horizon shrink to a point for $`M0`$. The additive constant related with the electric charge vanishes demanding that the electrically charged solution (47) reduces to the uncharged one (29) for $`Q=0`$. Therefore, the boundary term that must be added to the action is
$$B=\mathrm{\Delta }t(M+\varphi _{\mathrm{}}Q)+B_0,$$
(60)
where $`N_{\mathrm{}}`$ has been chosen equal to $`1`$, and $`B_0`$ is an arbitrary constant without variation. This proves that the integration constants $`M`$ and $`Q`$ appearing in the black hole metrics (47) and (29) are the mass and the electric charge respectively.
These results are confirmed also through an alternative method which holds for even dimensions, as is discussed in Appendix B.
### D Asymptotically flat limit $`(l\mathrm{})`$
The black hole metrics (29) and (47) tend asymptotically to an AdS spacetime with radius $`l`$, whose curvature satisfies $`R^{ab}l^2e^ae^b`$ at the boundary. Then, their asymptotically flat limit is obtained by taking $`l\mathrm{}`$. Thus, instead of taking the vanishing limit of the volume term $`(\alpha _00)`$, the vanishing cosmological constant limit of the action $`I_k`$ is obtained setting $`l\mathrm{}`$ in (9). This procedure is consistent with taking the same limit in the field equations (12) and (13).
When $`l\mathrm{}`$ the only non-vanishing term in (9) is the $`k`$th one, consequently the action is obtained from (2) with the following choice of coefficients:
$$\alpha _p:=\stackrel{~}{c}_p^k=\frac{1}{(d2k)}\delta _p^k.$$
(61)
Therefore, replacing (61) in (2), a new family of Lagrangians labeled by the integer $`k\{1,2,\mathrm{},[\frac{d1}{2}]\}`$, is obtained. For a fixed value of $`k`$, the Lagrangian is given just by $`L^{(k)}`$ defined in (3), so that the action reads
$$\stackrel{~}{I}_k=\frac{\kappa }{(d2k)}ϵ_{a_1\mathrm{}a_d}R^{a_1a_2}R^{a_{2k1}a_{2k}}e^{a_{2k+1}}e^{a_d},$$
(62)
where $`\kappa `$ is defined in (10). The field equations coincide with the $`l\mathrm{}`$ limit of (12), (13), which merely amounts to replacing $`\overline{R}^{ab}`$ by $`R^{ab}`$.
Note that for $`k=1`$, the standard EH action without cosmological constant is recovered, while for $`k=2`$ the Lagrangian is the Gauss-Bonnet density (18).
Static and spherically symmetric solutions of (62) lead to a similar picture as in the electrically (un)charged asymptotically AdS case: when $`k`$ is odd, one obtains only one solution describing a black hole, but for even values of $`k`$, two different solutions exist, one of them describes a black hole, while the other possesses naked singularities even when the mass bound holds.
It is simple to verify that black hole solutions of the action (62) correspond to the vanishing cosmological constant limit of the solutions for pure gravity (29). This also holds for the electrically charged solutions (47).
#### 1 $`Q=0:`$
The asymptotically flat solutions without electric charge are given by
$`ds^2`$ $`=`$ $`\left(1\left({\displaystyle \frac{2G_kM}{r^{d2k1}}}\right)^{1/k}\right)dt^2+`$ (64)
$`{\displaystyle \frac{dr^2}{1\left(\frac{2G_kM}{r^{d2k1}}\right)^{1/k}}}+r^2d\mathrm{\Omega }_{d2}^2.`$
The generic cases correspond to $`d2k10`$, for which the metrics (64) represent black hole solutions with an event horizon located at $`r_+=(2G_kM)^{1/(d2k1)}`$. As usual, their common vacuum geometry is the flat Minkowski spacetime, and their causal structure is described through the standard Penrose diagram of the Schwarzschild solution. In case of $`k=1`$ (EH), the Schwarzschild solution is recovered from (64) for $`d>3`$. Exceptional cases occur when $`d=2k+1`$, for which the action (62) correspond to a CS theory for the Poincaré group $`ISO(d1,1)`$. Their static, spherically symmetric solutions (64) do not describe black holes because they have a naked singularity at the origin. This can be inferred from (33) because when $`l\mathrm{}`$ the horizon recedes to infinity. For instance, in three dimensions, the solution (64) represent a conical spacetime .
#### 2 $`Q0`$:
The electrically charged asymptotically flat black hole solutions can be obtained for $`d>3`$ from (47) in the limit $`l\mathrm{}`$. As for the uncharged solutions, the generic case holds for $`d2k10`$, whose line elements read
$$ds^2=\left(1g_k(r)\right)dt^2+\frac{dr^2}{1g_k(r)}+r^2d\mathrm{\Omega }_{d2}^2,$$
(65)
with $`g_k(r)`$ given by
$$g_k(r)=\left(\frac{2G_kM}{r^{d2k1}}\frac{ϵG_k}{(d3)}\frac{Q^2}{r^{2(dk2)}}\right)^{\frac{1}{k}}.$$
(66)
For different generic values of the label $`k`$, the black hole solutions given by (65) resemble the Reissner-Nordström one, possessing two horizons which are found solving $`g_k(r)=1`$. As usual, these horizons satisfy $`0<r_{}<r_+`$ provided the mass is bounded from below by
$$Q^2\frac{(d2k1)}{ϵG_k}\left(\frac{(d3)G_kM}{dk2}\right)^{\frac{2d2k4}{d2k1}}.$$
(67)
The extreme case occurs when both horizons coalesce, that is
$$r_+=r_{}=\left(\frac{(d3)G_kM}{dk2}\right)^{\frac{1}{d2k1}},$$
(68)
so that the bound (67) is saturated.
The $`d`$-dimensional Reissner-Nordstrom solution is obtained from (65) setting $`k=1`$. Equation (67) reproduces the well known four-dimensional bound given by
$$Q_{EH}^2\frac{GM^2}{ϵ},$$
(69)
which is saturated when $`r_+=r_{}=G_kM`$, as can be seen from (68) for $`d=4`$ and $`k=1`$.
A further example corresponds to the electrically charged black hole in the vanishing cosmological constant limit of the BI action. The bound and the extreme radius are obtained in that case from (67) and (68) for $`d=2n`$ and $`k=n1`$:
$`Q_{BI}^2`$ $``$ $`{\displaystyle \frac{1}{ϵG_{n1}}}\left[{\displaystyle \frac{(2n3)G_{n1}M}{n1}}\right]^{2(n1)}`$ (70)
$`r_+`$ $`=`$ $`r_{}={\displaystyle \frac{(2n3)G_{n1}M}{n1}}.`$ (71)
The full set of asymptotically flat electrically charged black hole solutions (65) share a common feature with its asymptotically AdS counterparts given by (47) in the generic case $`(d2k10)`$. That is the existence of a timelike singularity for $`k1`$ located at the zero of $`g_k(r)`$ in (66) given by
$$r_e=\left(\frac{ϵ}{2(d3)}\frac{Q^2}{M)}\right)^{\frac{1}{d3}},$$
(72)
which satisfies $`0<r_e<r_{}<r_+`$ and is again interpreted as the smallest possible size of a spherical body with electric charge $`Q`$ and mass $`M`$. Then one concludes that this feature is absent only when one deals with the EH action with or without cosmological constant.
## IV Thermodynamics
### A Temperature
As usual, we define the black hole temperature by the condition that in the Euclidean sector, the solution be well defined (smooth) at the horizon. This means that the Euclidean time is a periodic coordinate with period
$$\tau =4\pi \left(\frac{df^2}{dr}|_{r_+}\right)^1,$$
(73)
which is identified with $`\beta =\frac{1}{\kappa _BT}`$, where $`\kappa _B`$ is the Boltzmann constant. Thus, the Hawking temperature is given by
$$T=\frac{1}{4\pi \kappa _B}\frac{df^2}{dr}|_{r_+}.$$
(74)
For the electrically uncharged cases, the black hole temperature for the set of metrics (29) is
$$T=\frac{1}{4\pi \kappa _Bk}\left((d1)\frac{r_+}{l^2}+\frac{(d2k1)}{r_+}\right).$$
(75)
For all $`k`$ such that $`d2k10`$, the function $`T(r_+)`$ exhibits the same behavior as the standard Schwarzschild-AdS black hole (which is obtained for $`k=1`$), that is: the temperature diverges at $`r_+=0`$. It has a minimum at $`r_c`$ given by
$$r_c=l\sqrt{\frac{d2k1}{d1}},$$
(76)
and grows linearly for large $`r_+`$. Considering $`k=n1`$, formula (75) reproduces the known results for BI ($`d=2n`$) and CS ($`d=2n1`$) black holes . The temperature (75) reaches an absolute minimum at $`r_c`$ equal to
$$T_c=\frac{\sqrt{\left(d2k1\right)\left(d1\right)}}{2\pi \kappa _Bkl},$$
(77)
provided the existence of a nonvanishing cosmological constant ($`l\mathrm{}`$).
In case of CS theory, that is when $`d2k1=0`$, $`T(r_+)`$ is not divergent at all, its absolute minimum is at $`r_c=0`$ and $`T_c=0`$. Thus, CS black holes are the only exceptional cases among all the possibilities considered here. Both, CS and generic cases are depicted in Figure $`1`$.
### B Specific Heat and Thermal Equilibrium
As seen in Section III.A, the black hole mass is a monotonically increasing function of $`r_+`$, therefore the behavior of $`T(M)`$ is qualitatively similar to that of $`T(r_+)`$.
Using (75) and (26), the specific heat $`C_k=\frac{M}{T}`$, can be expressed as a function of $`r_+`$,
$$C_k=k\frac{2\pi \kappa _B}{G_k}r_+^{d2k}\left(\frac{r_+^2+r_c^2}{r_+^2r_c^2}\right)\left(1+\frac{r_+^2}{l^2}\right)^{k1},$$
(78)
In case of $`d2k10`$, the specific heat (78) possesses an unbounded discontinuity at $`r_+=r_c`$ (see Figure $`1`$), signaling a phase transition. The specific heat $`C`$ is positive for $`r_+>r_c`$, and has the opposite sign for $`r_+<r_c`$.
Again, the CS case is exceptional. Setting $`d=2n1`$ and $`k=n1`$ in (78), the specific heat is found as
$$C_{CS}=\left(n1\right)\frac{2\pi \kappa _B}{G_{n1}}r_+\left(1+\frac{r_+^2}{l^2}\right)^{n2},$$
(79)
which is a continuous monotonically increasing positive function of $`r_+`$ and does not diverge for any finite value of $`r_+`$ .
The presence of a negative cosmological constant makes it possible for the family of black hole solutions (29) to reach thermal equilibrium, as is possible for the Schwarzschild-AdS<sub>4</sub> spacetime and for the three-dimensional black hole. Let us assume that any black hole described by (29) is immersed in a thermal bath of temperature $`T_0>T_c`$. If $`d2k10`$, the thermal behavior splits in two branches: for $`r_+<r_c`$, the specific heat is negative and therefore black hole state is driven away from that with temperature $`T_0`$; for $`r_+>r_c`$, the black hole state is attracted towards the equilibrium configuration at temperature $`T_0`$ (see Figure $`3`$). Thus, the temperature $`T_0`$ corresponds to two equilibrium states of radii $`r_u`$ (unstable) and $`r_s`$ (locally stable), with $`r_u<r_c<r_s`$. Neglecting quantum tunneling processes, there are two possible scenarios: if the initial black hole state has $`r_+<r_u`$, the black hole cannot reach the equilibrium because it evaporates until its final stage. Otherwise, for $`r_+>r_u`$, the black hole evolves towards an equilibrium configuration at $`r_+=r_s`$.
If the heat bath has temperature below $`T_c`$, the black hole cannot reach a stable equilibrium state and must evaporate, as depicted in Figure $`4`$.
None of the above arguments hold for the Chern-Simons case. When $`d2k=1`$, the specific heat (79) is always positive, therefore the equilibrium configuration is always reached, independently from the initial black hole state and for any finite temperature of the heat bath.
### C Entropy
It is well known that the partition function which describes the black hole thermodynamics is obtained through the Euclidean path integral in the saddle point approximation around the black hole solution . That is,
$`Ze^{I_E},`$
which means that the Euclidean action evaluated on the black hole configuration is identified with $`\beta `$ times the free energy of the system
$$I_E=\beta M\frac{S}{\kappa _B}+\beta \underset{i}{}\mu _iQ_i.$$
(80)
where the $`\mu _i`$’s are the chemical potentials corresponding to the charges $`Q_i`$. The Euclidean minisuperspace action is given by the Wick-rotated form of (55), that is
$`I_E`$ $`=`$ $`\beta {\displaystyle \underset{r_+}{\overset{\mathrm{}}{}}}{\displaystyle \frac{N}{2}}\left[{\displaystyle \frac{d}{dr}}\left\{{\displaystyle \frac{r^{d1}}{G_k}}\left[F(r)+{\displaystyle \frac{1}{l^2}}\right]^k\right\}{\displaystyle \frac{1}{ϵ}}r^{d2}p^2\right]𝑑r+`$ (82)
$`{\displaystyle \frac{1}{ϵ}}\beta {\displaystyle \underset{r_+}{\overset{\mathrm{}}{}}}A_0{\displaystyle \frac{d}{dr}}\left(r^{d2}p\right)𝑑r+B_E,`$
In what follows we shall consider the electrically uncharged cases only. The bulk part of the Euclidean action is a linear combination of the constrains and therefore, its on-shell value is given by the boundary term $`B_E`$. This boundary piece is determined by the requirement that $`I_E`$ be stationary on the black hole geometry. Varying (82) leads to
$$\delta I_E=\frac{\beta N_{\mathrm{}}}{2G_k}\underset{r_+}{\overset{\mathrm{}}{}}\frac{d}{dr}\left\{r^{d1}\delta \left[F(r)+\frac{1}{l^2}\right]^k\right\}𝑑r+\delta B_E,$$
(83)
on shell. From this expression, one finds
$`\delta B_E=\beta \delta M{\displaystyle \frac{2\pi k}{G_k}}r_+^{d2k1}\left(1+{\displaystyle \frac{r_+^2}{l^2}}\right)^{k1}\delta r_+,`$
where $`N_{\mathrm{}}`$ has been set equal to one and we have used $`\frac{df^2}{dr}|_{r_+}=4\pi \beta ^1`$. From (80) one identifies
$$\delta S=k\frac{2\pi \kappa _B}{G_k}r_+^{d2k1}\left(1+\frac{r_+^2}{l^2}\right)^{k1}\delta r_+,$$
(84)
which is integrated into
$$S_k=k\frac{2\pi \kappa _B}{G_k}\underset{0}{\overset{r_+}{}}r^{(d2k1)}\left(1+\frac{r^2}{l^2}\right)^{k1}𝑑r.$$
(85)
This is a monotonically increasing function of $`r_+`$, in agreement with the second law of thermodynamics. In (85) the lower limit in the integral has been fixed by the condition $`S_k(r_+=0)=0`$ for the whole set of black holes given by (29).
For the EH action (that is for $`k=1`$), expression (85) readily reproduces, for the Schwarzschild-AdS solution
$`S_{EH}={\displaystyle \frac{2\pi \kappa _B}{(d2)G}}r_+^{d2},`$
which in standard units is the celebrated “area law”
$`S_{EH}={\displaystyle \frac{\kappa _B}{\stackrel{~}{G}}}{\displaystyle \frac{A}{4}}.`$
For $`k=\left[\frac{d1}{2}\right]`$ (BI and CS), formula (85) reduces to the known results . The theory described by $`I_2`$ in (19) is an intrinsically higher dimensional one, and the corresponding black hole entropy is given by
$$S_2=\frac{4\pi \kappa _B}{G_\mathrm{𝟐}}r_+^{d4}\left[\frac{1}{(d4)}+\frac{r_+^2}{(d2)l^2}\right].$$
(86)
Hence, the area law is a peculiarity of the Einstein-Hilbert theory ($`k=1`$), while for $`k1`$ the entropy (85) becomes proportional to the area in the large $`r_+`$ limit, that is
$$S_kk\frac{2\pi \kappa _B}{(d2)G_kl^{2(k1)}}r_+^{d2}=k\frac{G}{G_kl^{2(k1)}}S_{EH},$$
(87)
with $`r_+>>l`$.
### D Asymptotically flat limit
In the limit $`l\mathrm{}`$, the geometry of the uncharged black hole is given by (64) whose corresponding temperature is
$$T^0=\frac{1}{4\pi \kappa _Bk}\frac{(d2k1)}{r_+}.$$
(88)
This gives a vanishing value for CS theory $`(d2k1=0)`$, which is consistent with the fact that in that case, the geometry possesses a singularity which is not surrounded by a horizon in the limit $`l\mathrm{}`$, so that no temperature can be associated with it. For all the other cases $`(d2k10)`$, the horizon is located at $`r_+=(2G_kM)^{1/(d2k1)}`$, so that the black hole temperature (88) is a monotonically decreasing function of the mass. Therefore, thermal equilibrium can never be reached, consistently with the fact that the specific heat is always negative
$$C^0=k\frac{2\pi \kappa _B}{G_k}r_+^{d2k}.$$
(89)
The entropy is also an increasing function of $`r_+`$,
$$S_k^0=k\frac{2\pi \kappa _B}{G_k}\frac{r_+^{(d2k)}}{(d2k)},$$
(90)
which is consistent with the second law of thermodynamics. Note that formula (90) is proportional to the area of the horizon only for $`k=1`$ (EH). Thus, in the $`l\mathrm{}`$ limit, the area law cannot be recovered even as an approximation in the cases with $`k1`$.
### E Canonical Ensemble
In four dimensions, Hawking and Page have shown that in the presence of a negative cosmological constant, the partition function in the canonical ensemble is well defined, unlike in case of a vanishing $`\mathrm{\Lambda }`$ . The same argument can be extended for higher dimensions for the whole set of theories (8) labelled by $`k`$.
The partition function in the canonical ensemble reads
$$Z(\beta )=_0^{\mathrm{}}e^{\beta M}\rho (M)𝑑M,$$
(91)
where $`\rho (M)=\mathrm{exp}\left(\frac{S_k}{\kappa _B}\right)`$ is the density of states as a function of the energy. The convergence of this integral depends on the asymptotic behavior of $`S_k`$ for large $`M`$,
$`S_ka_{d,k}M^{\left(\frac{d2}{d1}\right)},`$
where $`a_{d,k}`$ is a positive constant. Thus, the integrand of (91) goes as $`\mathrm{exp}(\beta M+\kappa _B^1a_{d,k}M^{\left(\frac{d2}{d1}\right)})`$ and therefore the partition function converges.
This argument breaks down in the $`l\mathrm{}`$ limit: in that case, the entropy is
$`S_k^0=a_{d,k}^0M^{\left(\frac{d2k}{d2k1}\right)},`$
with $`a_{d,k}^0`$ a different positive constant, which yields a divergent partition function.
The lesson one can draw from this exercise is that the presence of a negative cosmological constant is sufficient to render the canonical ensemble well defined for all the theories described here.
## V Summary and Discussion
### A Theories described by the action $`I_k`$
We have examined a family of gravitation theories in dimension $`d`$, whose common feature is to possess vacuum solutions with maximal symmetry. This means that the theories –described by the action $`I_k`$– have a unique cosmological constant. For a given $`d`$ there exist $`[\frac{d1}{2}]`$ different theories labeled by the integer $`k`$, which is the highest power of curvature in the Lagrangian. For $`k=1`$, the EH action is recovered, while for the largest value of $`k`$, that is $`k=[\frac{d1}{2}]`$, BI and CS theories are obtained. These three cases exhaust the different possibilities up to six dimensions, and new interesting cases arise for $`d7`$. For instance, the case with $`k=2`$, which is described by the action (19), exists only for $`d>4`$: In five dimensions this theory is equivalent to CS, for $`d=6`$ it is equivalent to BI, and for $`d=7`$ and up, it defines a new class of theories.
### B Special cases selected from cosmic censorship
A first distinction between the different theories mentioned above comes from the study of their spherically symmetric, static solutions. It is found that for odd $`k`$, physical black holes satisfying the cosmic censorship criterion exist. For even $`k`$, however, both physical black holes and solutions with naked singularities with positive mass exist. This already casts doubt on the soundness of this subset of theories. Moreover, the absence of a cosmic censorship principle would be in conflict with the existence of a positive energy theorem obtained from supersymmetry. This means that the supersymmetric extensions of the theories considered here can be expected to be very different for odd and even $`k`$. In fact, as it has been shown in , CS theories with even $`k`$ –defined for $`d=5,9,\mathrm{}`$–, have a supersymmetric extension based on superunitary groups, whereas for odd $`k`$ $`(d=3,7,11,\mathrm{})`$ the corresponding supergravities are based on the orthosymplectic groups.
The different theories considered here are summarized in the scheme shown in Fig. 5.
Here we have highlighted the odd $`k`$ columns as they would represent better candidates for physical theories based on the criterion of cosmic censorship versus supersymmetry.
Note that CS theories are the representatives of the lowest possible dimension for a given $`k`$. Moreover, CS gravity theories exhibit local AdS symmetry whereas all other gravitation theories of the same dimension only have local Lorentz invariance (see Appendix A).
Over the years, 11-dimensional spacetime has been believed to be the arena for the ultimate unified theory. From the present analysis, it follows that in $`d=11`$, the cases $`k=1,3,5`$ are of special interest. The supersymmetric extension for $`k=1`$ is the famous Cremmer-Julia-Scherk supergravity , which only exists if the cosmological constant vanishes . The supersymmetric extension for $`k=5`$ with a finite $`\mathrm{\Lambda }`$ is also known , whose vanishing cosmological constant version is described in . The corresponding supersymmetric extension of the gravity theory with $`k=3`$ is an open problem.
### C Black Holes
For all dimensions and for any $`k`$, there exist well behaved black hole solutions, in the sense that the singularities are hidden by an event horizon. For $`d2k1`$, the causal structure of these black holes is the same as that of Schwarzschild-AdS and Reissner-Nordström-AdS spacetimes. However, this set of black holes differs from standard $`d`$-dimensional Schwarzschild and Reissner-Nordström solutions in that their asymptotic behavior, with respect to the vacuum, is given by $`g_{00}\overline{g}_{00}r^{\left(\frac{d2k1}{k}\right)}`$. Again, the CS case stands separate from the rest, in that the causal structure of the vacuum is the same as that of $`2+1`$ dimensions, and analogously, there is a mass gap between the $`M=0`$ black hole and AdS spacetime $`\left(M=\frac{1}{2G_{n1}}\right)`$. Furthermore, in the vanishing cosmological constant limit, the CS theory supports no static, spherically symmetric black holes.
In the electrically charged case, the black holes for $`k1`$ predict a minimum size for a physical source. It is noteworthy that the geometry encodes this restriction for all cases, except for the EH action.
### D Thermodynamics
The presence of a negative cosmological constant for the entire set of theories described by the action $`I_k`$ makes it possible for black holes to reach thermal equilibrium with a heat bath. The AdS radius $`l`$ acts as a regulator allowing the canonical ensemble to be well defined, unlike the case of zero cosmological constant. The black hole entropy obeys the area law only in the case $`k=1`$. For other values of $`k`$, the entropy respects the second law of thermodynamics, because $`\frac{dS}{dr_+}>0`$, but the area law is recovered only in the limit $`\frac{r_+}{l}\mathrm{}`$.
In the limit $`\mathrm{\Lambda }0`$, the area law never holds, except for $`k=1`$. In that limit, the temperature has no minimum and consequently the thermodynamic equilibrium cannot be reached.
The thermodynamic behavior is qualitatively the same as the Schwarzschild-AdS<sub>4</sub> black hole in the generic cases $`d2k1`$. On the other hand, Chern-Simons black holes for odd dimensions behave like the $`d=3`$ case.
In the generic cases, black holes have a minimum temperature $`T_c`$ at $`r_+=r_c=l\sqrt{\frac{d2k1}{d1}}`$, so that –as is depicted in Figure $`3`$– those whose horizon radius exceed the unstable equilibrium position $`r_u`$ can reach equilibrium with a heat bath of temperature higher than $`T_c`$. If the heat bath has a temperature below $`T_c`$, or $`r_+<r_u`$, the black holes evaporate.
In the CS case, the temperature grows linearly with $`r_+`$, hence there is no critical temperature and the thermal equilibrium is always attained.
In an equilibrium configuration, the free energy $`F=MTS`$ can be expressed as a function of $`r_+`$. For fixed $`k`$ the behavior of $`F`$ can be found from (26), (75) and (85) as
$`F(r_+`$ $``$ $`0){\displaystyle \frac{r_+^{d2k1}}{2(d2k)G_k}},`$ (93)
$`F(r_+`$ $``$ $`\mathrm{}){\displaystyle \frac{r_+^{d1}}{2(d2)G_kl^{2k}}}.`$ (94)
This change in sign has been interpreted as an indication that, for small $`r_+`$ the black hole would be unstable for decay into AdS spacetime, while for large $`r_+`$ the black hole would be stable . This suggests that a phase transition would occur at $`F(r_+)=0`$. This conclusion, however contradicts the fact that the phase transition actually occurs at the critical value $`r_c`$, where the specific heat $`C`$ changes sign, and which does not coincide with the zero of $`F(r_+)`$. In particular, considering the EH action $`(k=1)`$, the change of sign in $`F`$ occurs at $`r_+=l`$ while $`r_c=l\sqrt{\frac{d3}{d1}}<l`$. Moreover, for the CS case, $`d2k=1`$, there is no phase transition at all, although $`F`$ still has a change in sign. The source of the disagreement lies in that the canonical ensemble is defined keeping $`T`$ fixed, while the limits in (93) and (94) do not respect this condition.
¿From all the evidence presented here, it is apparent that CS theories form an exceptional class: They are genuine gauge theories whose supersymmetric extension is known; their black hole spectrum has a mass gap separating it from AdS spacetime, and these black holes possess remarkable thermodynamical properties. CS black holes can reach thermal equilibrium with a heat bath at any temperature, and the positivity of the specific heat guarantees their stability under thermal fluctuations.
In contrast with the generic case, a small CS black hole is stable against decay by Hawking radiation. This suggests that, as in the three dimensional case, CS (super)gravities could have a well defined quantum theory.
## VI Acknowledgments
The authors are grateful to R. Aros, M. Bañados, M. Contreras, M. Henneaux, C. Martínez, F. Méndez, R. Olea, M. Plyushchay, J. Saavedra and C. Teitelboim for many enlightening discussions and helpful comments. This work was supported in part through grants 1990189, 1980788 from FONDECYT, and by the “Actions de Recherche Concertées” of the “Direction de la Recherche Scientifique - Communauté Française de Belgique”, by IISN - Belgium (convention 4.4505.86). The institutional support of Fuerza Aérea de Chile, I. Municipalidad de Las Condes, and a group of Chilean companies (AFP Provida, CODELCO, Empresas CMPC, and Telefónica del Sur) is also recognized. CECS is a Millenium Science Institute. J. Z. wishes to thank the organizers of the $`1999`$ ICTP Summer Workshop on Black Hole Physics for hospitality in Trieste. J. C. and J.Z thank the organizers of the V La Hechicera School, Mérida.
## VII Appendix
### A CS & BI Theories
Requiring that the integrability conditions of equation (4) do not impose further algebraic constraints on the curvature or the torsion beyond Eq. (5) implies that the coefficients $`\alpha _p`$’s in Eq. (2) satisfy a recursive equation, whose solution fixes them in terms of the gravitational and cosmological constants . An equivalent way to express this is that the $`\alpha _p`$’s become fixed as in equation (9) with $`k=[\frac{d1}{2}]`$, just requiring the existence of a sector in the theory with propagating torsion. Thus, in $`d=2n`$ dimensions, the Lagrangian reads
$$L=\frac{\kappa l^2}{2n}ϵ_{a_1\mathrm{}a_d}\overline{R}^{a_1a_2}\mathrm{}\overline{R}^{a_{d1}a_d},$$
(95)
where<sup>††</sup><sup>††</sup>††A positive cosmological constant is obtained making $`l^2l^2`$. $`\overline{R}^{ab}:=R^{ab}+\frac{1}{l^2}e^ae^b`$.
The expression (95) is proportional to the Pfaffian of the $`2`$-form $`\overline{R}^{ab}`$ and, in this sense, it has a Born-Infeld-like form :
$$L=2^{n1}(n1)!\kappa l^2\sqrt{det\left(R^{ab}+\frac{1}{l^2}e^ae^b\right)}.$$
(95)
For $`d=2n1`$ dimensions, the Lagrangian is given by the Euler-Chern-Simons form for the AdS group, whose exterior derivative is proportional to the Euler density in $`2n`$ dimensions,
$`dL_{G\mathrm{\hspace{0.33em}2}n1}^{AdS}`$ $`=`$ $`{\displaystyle \frac{\kappa l}{2n}}ϵ_{A_1\mathrm{}A_{2n}}\overline{R}^{A_1A_2}\mathrm{}\overline{R}^{A_{2n1}A_{2n}}`$ (96)
$`=`$ $`\overline{\kappa }_{2n},`$ ()
where $`\overline{R}^{AB}`$ stands for the AdS curvature. This Lagrangian was discussed in and also in for torsion-free manifolds.
Additional terms which depend explicitly on the torsion are required by local supersymmetry and they can be consistently added to the Lagrangian only for $`d=4m1`$ .
These torsional Lagrangians are odd under parity and are obtained from the Chern characters associated with the AdS curvature in $`4m`$ dimensions. Furthermore, the coefficients in front of the different terms in these torsional Lagrangians are necessarily quantized. The odd dimensional action, with or without torsional terms, has a larger local symmetry given by $`SO(d1,2)`$, so that beyond standard local Lorentz symmetry ($`\delta e^a=\lambda _b^ae^b`$ and $`\delta \omega ^{ab}=D\lambda ^{ab}`$), these theories are invariant also under local “AdS-translations:”
$`\delta e^a`$ $`=`$ $`D\lambda ^a`$ (97)
$`\delta \omega ^{ab}`$ $`=`$ $`{\displaystyle \frac{1}{l^2}}(\lambda ^ae^b\lambda ^be^a).`$ ()
### B Conserved Charges from a Background-Independent Surface Integral
If one deals with more general solutions possessing different isometries, the identification of the integration constants with the conserved charges through the minisuperspace trick does not work, because in general the reduced action does not lead to the true extremum of the original action. The Hamiltonian method provides a way to express the mass as a surface integral . However, this procedure requires the invertibility of the symplectic matrix associated with the action $`I_k`$. This is impossible to perform globally in phase space, because there are field configurations for which the symplectic form degenerates. Therefore, no general formula could be found for an arbitrary field configuration.
A way to circumvent this problem is carried out in $`d=2n`$ following a recently proposed method which is appropriate to deal with asymptotically AdS spacetimes.
Consider the action $`I_k`$ defined in (8). In first order formalism, the existence of an extremum of $`I_k`$ for asymptotically locally AdS spacetimes fixes the boundary term that must be added to the action as being proportional to the Euler density multiplied by a fixed weight factor. Hence, in order to cancel the boundary term coming from the variation of $`I_k`$, the total action including the boundary term –up to a constant–, is given by
$$I_T=I_k+\kappa \alpha _n_{2n},$$
(98)
with
$$\alpha _n=c_n^k:=\frac{(1)^{n+k+1}l^{2(nk)}}{2n\left(\genfrac{}{}{0pt}{}{n1}{k}\right)}.$$
(99)
The total action $`I_T`$ is invariant under diffeomorphisms by construction, because $`I_k`$ is written in terms of differential forms. Thus, Noether’s theorem provides a conserved current $`(dJ=0)`$ associated with this invariance, which can be locally written as $`J=dQ`$. Assuming the topology of the manifold to be of the form $`=R\times \mathrm{\Sigma }`$, this procedure yields a regularized and background-independent expression for the conserved charges associated with a Killing vector $`\xi `$, which is globally defined on the boundary of the spatial section $`\mathrm{\Sigma }`$. The surface integral reads
$$Q(\xi )=\underset{\mathrm{\Sigma }}{}\xi ^\mu \omega _\mu ^{ab}𝒯_{ab},$$
(100)
where, $`𝒯_{ab}`$ is the variation of the total Lagrangian with respect to the curvature
$$𝒯_{ab}:=\frac{\delta L_T}{\delta R^{ab}}=\underset{p=1}{\overset{n}{}}c_p^kp𝒯_{ab}^p,$$
(101)
with
$$𝒯_{ab}^p=\kappa ϵ_{aba_3\mathrm{}a_d}R^{a_3a_4}\mathrm{}R^{a_{2p1}a_{2p}}e^{a_{2p+1}}\mathrm{}e^{a_d},$$
(102)
and where the coefficients $`c_p^k`$ are defined through equations (9) and (99).
The mass is obtained from (100) when $`\xi =_t`$, without making further assumptions about the matching with a background geometry nor with its topology.
One way to check this result is evaluating the mass for the black hole metrics (29), which leads to the expected result
$$Q(_t)=M.$$
(103)
It is a simple exercise to check that formula (100) vanishes when evaluated on any constant curvature spacetime – satisfying $`\overline{R}^{ab}=R^{ab}+l^2e^ae^b=0`$ – which admits at least one Killing vector. This means that spaces which are locally AdS have vanishing Noether charges for the whole set of theories defined by $`I_k`$ in even dimensions. These spaces in general possess non-trivial topologies and could be regarded as different possible vacua. Hence one can find massive solutions which correspond to excitation of the corresponding vacuum in the same topological sector.
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# Observations of Eddington-limited type-I X-ray bursts from 4U 1812–12
## 1 Introduction
Since its first Uhuru detections (Forman, Jones, & Tananbaum (1976); Forman et al. (1978)), 4U 1812$``$12 was observed by several satellite X-ray experiments: OSO 7 (1M 1812-121, Markert et al. (1979)), Ariel V (3A 1812-121, Warwick et al. (1981)), HEAO 1 (1H 1815$``$121, Wood et al. (1984)), and EXOSAT (GPS 1812-120, Warwick et al. (1988)). From these observations, it is clear that 4U 1812$``$12 is a persistent, though variable, source. Uhuru found a 2–10 keV maximum intensity of $`5\times 10^{10}\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$ and a variability of a factor of at least 2 (Forman et al. (1978)). Similar variability characteristics were observed by Ariel V, as the source varied in the range $`36\times 10^{10}\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$ in the same energy band (Warwick et al. (1981)), while lower intensities of $`2\times 10^{10}\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$ and $`3\times 10^{10}\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$ were measured by HEAO 1 (2–10 keV, Wood et al. (1984)) and EXOSAT (2–6 keV, Warwick et al. (1988)). The 3–10 keV source spectrum as obtained by the EXOSAT GSPC was best fitted by a power law (Gottwald et al. (1995)). 4U 1812$``$12 is being monitored by RXTE-ASM since February 1996 <sup>1</sup><sup>1</sup>1the ASM measurements are publicly available at URL http://www.space.mit.edu/XTE, confirming its previously reported characteristics. The source is always detected, with an average 2–10 keV flux of $`3.8\times 10^{10}\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$ ($`20`$ mCrab) and a variability of a factor $`3`$ on $`1`$ week time scale.
Three X-ray bursts were detected from this source in 1982 by Hakucho (Murakami et al. (1983)). Two of the events showed clear evidence for photospheric radius expansion, and reached a maximum 1–22 keV intensity of $`1.7\times 10^7\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$. The burst spectra were consistent with a $`2.5`$ keV blackbody emission, and showed evidence for softening during the exponential decay (e-folding time $`\tau 20`$ s). This indicated the bursts to be type-I, i.e. thermonuclear flashes originating on the hot surface of a neutron star, and the source to be likely located in a low-mass X-ray binary. During the observation no persistent emission was detected above 20 Uhuru flux units ($`5\times 10^{10}\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$ in 2-10 keV). Anyway this is not in disagreement with the identification of the Hakucho burster with the persistent source 4U 1812$``$12. Murakami et al. (1983) also proposed the association of the burst source with the transient Ser X-2, observed once in 1965 (Friedman, Byram & Chubb (1967)). More recently, a single-peaked burst from 4U 1812$``$12 was observed by the BeppoSAX-WFC instrument on 1997 Mar 12.2209 UT (Burderi et al. (1997)). The event had a peak intensity of 1.2 and 0.6 Crab in the 1.5–10 keV and 10$``$26 keV band respectively, and a decay time of $`20`$ s.
4U 1812$``$12 is classified as an atoll source, which is common among the type-I X-ray bursters, and shows band-limited noise and a $`0.8`$ Hz QPO (Wijnands & van der Klis (1999)).
In this paper we investigate the burst characteristics of 4U 1812$``$12, about 15 years after the Hakucho observations of its type-I bursts. The event observed by Burderi et al. (1997) is also re-analysed. Through a homogeneus sample of Eddington-limited type-I X-ray bursts, we accurately estimate the source distance and test the reliabilty of near-Eddington bursts as a standard candle. In the next section we briefly introduce the Wide Field Cameras telescopes and report on the observations of 4U 1812$``$12. Time-resolved spectroscopy of the bursts is presented in Section 3, while the scientific implications of our results are discussed in Section 4.
## 2 Observations
The Wide Field Cameras (WFC) on board the BeppoSAX satellite consist of two identical coded aperture multi-wire proportional counters (Jager et al. (1997)) pointing in opposite directions. Each camera covers a $`40\mathrm{°}\times 40\mathrm{°}`$ full width to zero-response field of view, the largest ever flown for an arcminute resolution X-ray telescope. With their source location accuracy in the range $`1\mathrm{}`$$`3\mathrm{}`$ (99% confidence), a time resolution of 0.488 ms, and an energy resolution of 18% at 6 keV, the WFCs are effective in studying X-ray transient phenomena in the 2–28 keV bandpass. The imaging capability and the good instrument sensitivity (5-10 mCrab on-axis in $`10^4`$ s, depending on the number of sources in the field) allow an accurate monitoring of complex sky regions, like the Galactic centre.
One of the main scientific objectives of the WFCs is the study of the timing and spectral behaviour of both transient and persistent sources of the Galactic Bulge region on time scales ranging from seconds to years. To this end, an observation program of systematic wide field monitoring of the Sgr A sky region is being carried out since August 1996 (see e.g. Heise (1998); Heise et al. (1999); Ubertini et al. (1999)). This program consists of a series of observations, each lasting $`60`$ ks, nearly weekly spaced throughout the two visibility periods (August-October and February-April) of the Galactic Centre region.
The WFC Galactic Bulge monitoring program is significantly contributing in the study of X-ray bursting sources. Up to now, a total of 15 new objects were discovered in $`3.2`$ years observing time, thus enlarging the population of the bursters by $`35\%`$ (Heise et al. (1999); Ubertini et al. (1999)). The data of the two cameras are systematically searched for bursts and flares by analyzing the time profiles of the detectors in the 2–11 keV energy range with a time resolution down to 1 s. Reconstructed sky images are generated for any statistically meaningful event and the accuracy of the reconstructed position, which of course depends on the burst intensity, is typically better than $`5\mathrm{}`$. This analysis procedure has led to the identification of $`950`$ X-ray bursts (156 of which from the Bursting Pulsar GRO J1744$``$28) in a total of about $`4\times 10^6`$ s net observing time (e.g. Cocchi et al. 1998a ).
Whenever the WFCs point at the Galactic Centre region, 4U 1812$``$12 is in the field of view, though at a rather offset position, being $`18\mathrm{°}`$ away from Sgr A ($`l_{\mathrm{II}}=18.0\mathrm{°},b_{\mathrm{II}}=2.4\mathrm{°}`$). Due to the source’s relatively low intensity and to the unfavourable pointing, the WFC data is not sensitive enough to study the persistent emission with sufficient accuracy. A total of 8 X-ray bursts were detected at a position consistent with that of 4U 1812$``$12 during all the time spent on the source by the WFCs both in primary and in secondary observing mode ($`4`$ Ms net time in 3.2 years). None of the observed bursts can be associated with other known sources.
We analysed all the observed bursts, including the one already investigated by Burderi et al. (1997). The main characteristics of the observed bursts (hereafter burst A,B,…,H, chronologically) are summarised in Table 1. One of the bursts, namely burst C, was observed with much better statistics than the others, as it was in a more favourable position ($`10\mathrm{°}`$ offset). For this reason, burst C was analysed with a higher time resolution than the others.
## 3 Data Analysis and Results
Energy-resolved time histories of the bursts were constructed by accumulating the detector counts associated with the shadowgram obtained for the sky position of the analysed source, thus improving the signal-to-noise ratio of the profile. For a given source, the background is the sum of (part of) the diffuse X-ray background, the particles background and the contamination of other sources in the field of view. Source contamination is the dominating background component for crowded sky fields like the Galactic Bulge. Nevertheless, the probability of source confusion during a short time-scale event (10–100 s) like an X-ray burst is negligible. The time profiles of burst C were accumulated in three different energy bands (2–5, 5–10, and 10–28 keV, see Fig. 1). The time histories of the other bursts, due to their limited counting statistics, were obtained for the bands 2-8 and 8-28 keV only (Fig.2). The 2-28 keV time profiles of all the bursts are characterized by fast rise times (within a few seconds) and longer exponential decay with e-folding times of $`15`$ s. The high energy time histories of almost all the bursts show clear evidence for double-peaked profiles (see Fig.1 and Fig.2) and their e-folding times are significantly shorter ($`47`$ s) than the ones of the low energy profiles ($`1520`$ s), as reported in Table 1.
The integrated spectra of the eight bursts are all consistent with absorbed blackbody radiation with average colour temperatures of $`2`$ keV and an average blackbody radius of the emitting sphere of $`20`$ km assuming a standard 10 kpc source distance (Table 1). The spectra were subtracted for the source persistent emission, which accounts for only $`0.5\%`$ of the burst peak intensity. The $`N_\mathrm{H}`$ parameter could not be satisfactorily constrained for any of the bursts, so we kept its value fixed according to the interpolated value computed at the source position, namely $`N_\mathrm{H}=7.3\times 10^{21}\mathrm{cm}^2`$ (Dickey & Lockman (1990)).
Time-resolved spectra were accumulated for all the bursts, in order to study the time evolution of their spectral parameters. Thanks to the good counting statistics, the spectral analysis of burst C could be performed with a time resolution of 1 s. Conversely, for each of the other bursts only four time resolved spectra could be obtained with a poorer time resolution (4 s). The time intervals of the four spectra were chosen to match the first-peak, interpeak, second-peak and decay phases in the corresponding 8-28 keV burst time profiles (Fig.2). Since all the obtained spectra are consistent with absorbed blackbody emission, the time histories of the colour temperature and the emitting sphere radius can be determined (Table 2 and Fig.3). The blackbody radii were calculated assuming isotropic emission at a source distance of 10 kpc and not correcting for gravitational redshift and conversion to effective blackbody temperature from color temperature (see Lewin, van Paradijs, & Taam (1993) for details). A radius expansion by a factor of $`5`$ is observed in burst C while, probably due to the larger time bins used, a lower expansion factor ($`2`$) is obtained for all the other bursts but burst H, which does not show evidence for radius variations.
## 4 Discussion
Following the classification proposed by Hoffman, Marshall, & Lewin (1978), X-ray bursts are classified in two main types (type-I, type-II, see Lewin, van Paradijs & Taam (1993) for a comprehensive review).
On the basis of the spectral and timing properties of the eight bursts observed by the WFCs, it is apparent that 4U 1812$``$12 is a type-I burster. In fact, the blackbody emission and the measured colour temperatures of $`2`$ keV are consistent with type-I bursting. Moreover, spectral softening is observed in the time resolved spectra of the bursts (Table 2 and Figure 3), and the bursts time profiles can be fitted with exponential decays whose characteristic times are significantly shorter at higher energies (see e.g. Table 1). These results confirm the measurements obtained 15 years earlier by Hakucho (Murakami et al. (1983)), indicating 4U 1812$``$12 is a neutron-star low-mass X-ray binary.
The photospheric radius expansion derived from the time resolved spectral analysis of most of the observed bursts can be interpreted as adiabatic expansion during a high luminosity (Eddington-limit) type-I burst. Actually, the double-peaked profiles observed in the high energy (above 8 keV) time histories of the bursts (with the only exception of burst H) are typical of super-Eddington events (e.g. Lewin, van Paradijs, & Taam (1995)). Even though burst H is not double-peaked, its peak luminosity is consistent with those of the other observed events. Moreover, its 8–28 keV time profile could be associated with the flat-top profiles of some observed Eddington-limited type-I bursts (Lewin, van Paradijs, & Taam (1995)), so burst H too can be regarded as an event with peak luminosity close to the Eddington limit. We also notice that the two bursts which show less clear evidence for photospheric radius expansion, namely bursts E and H, are the less energetic ones, their fluences being the lowest observed (Table 1). We can regard their total energy release ($`5\times 10^{39}`$ erg assuming a 4 kpc distance, see later) as the minimum needed to drive the expansion of the photosphere of the neutron star in 4U 1812$``$12.
Eddington-luminosity X-ray bursts can lead to an estimate of the source distance, assuming the burst emission to be isotropic and the peak flux to be very close to the Eddington luminosity. Actually, the peak intensities observed for the eight events are all consistent with a constant value of $`4.53\pm 0.09`$ Crab (2–28 keV); the associated reduced $`\chi ^2`$ is 1.03 for 7 d.o.f.. This average peak intensity extrapolates to an unabsorbed bolometric flux of $`(15.03\pm 0.29)\times 10^8\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$. The consistency of the peak luminosities of all the bursts with a constant value supports the adoption of the peak bolometric intensities of super-Eddington bursts as a standard candle. An average luminosity of Eddington-limited bursts was empirically calculated by Lewin, van Paradijs and Taam (1995) on a sample of bursters whose distance was estimated with other methods: a luminosity value of ($`3.0\pm 0.6\times 10^{38}\mathrm{erg}\mathrm{s}^1`$) was obtained. The adoption of this standard luminosity leads to a distance value $`d=4.1\pm 0.5`$ kpc for 4U 1812$``$12. On the other hand, assuming the theoretical Eddington luminosity for a typical $`1.4\mathrm{M}_{\mathrm{}}`$ neutron star ($`2\times 10^{38}\mathrm{erg}\mathrm{s}^1`$) we obtain $`d3.3`$ kpc. For the calculated distance of $`4`$ kpc, and with the simple assumptions on the burst emission made in Section 3, an average radius of $`8\pm 1`$ km for the blackbody emitting region during the bursts is obtained. This value, which supports the neutron-star nature of the collapsed object, should indeed be regarded as a lower limit for the actual neutron star radius, according to Ebisuzaki, Hanawa & Sugimoto (1984).
Assuming the source’s persistent spectrum to be consistent with the one suggested by Barret et al. (2000) for the X-ray bursters in low state, i.e. a Comptonized spectrum with electron temperature $`k\mathrm{T}_e25`$ keV and $`\tau 3`$, the bolometric luminosity of 4U 1812$``$12 can be extrapolated. For a distance of $`4`$ kpc and an average 2–10 keV persistent intensity of $`4\times 10^{10}\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$, we obtain $`L_{bol}5.6\times 10^{36}\mathrm{erg}\mathrm{s}^1`$. We also derive, for a canonical $`1.4\mathrm{M}_{}`$ neutron star with a radius of 10 km, an average accretion rate of $`5\times 10^{10}\mathrm{M}_{}\mathrm{y}^1`$. These values are common among low-mass X-ray binaries.
Due to the non-continuous WFCs coverage of the Galactic Centre region and to the BeppoSAX orbit characteristics, we cannot accurately establish the burst occurrence rate of 4U 1812$``$12. Anyway the minimum observed intervals are 6.53 d and 6.16 d for bursts A-B (August 1996) and F-G (October 1998) respectively, and such intervals are of the same order of magnitude than the one measured by Hakucho in 1982 (4.61 d). Under the hypothesis that the above values are the actual wait times for bursts B and G, we can calculate the ratio $`\alpha =E_p/E_b`$, where $`E_b`$ and $`E_p`$ are the bolometric fluences of the burst and of the persistent emission between two contiguous bursts, respectively. Average 2–10 keV intensities of $`(4.4\pm 0.4)\times 10^{10}\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$ and $`(3.5\pm 0.9)\times 10^{10}\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$ were measured by RXTE-ASM between bursts A-B and F-G respectively. Again, with the former assumptions on the spectrum of the persistent emission, we determine $`\alpha =(6.4\pm 1.0)\times 10^2`$ and $`\alpha =(5.8\pm 1.9)\times 10^2`$ for the two events. These values are consistent with each other and are within the observed range ($`1010^3`$, distribution peaking at $`10^2`$) for the $`\alpha `$ parameter of known X-ray bursters, even if on the higher side. This is suggestive of helium-burning with no spare fuel left for the next burst, and possibly of steady burning of part of the accreted matter (Lewin, van Paradijs, & Taam (1993); van Paradijs et al. (1988)).
As pointed out above, the eight bursts we analysed show very similar features. Moreover, their characteristics are also consistent with those of the bursts detected by Hakucho (Murakami et al. (1983)). For the Hakucho bursts, 1–22 keV peak intensities of $`4.6`$ Crab were measured. Also the event Burderi et al. (1997) reported to have very different burst parameters, i.e. single peak profile and much lower luminosity (1.2 Crab in 1.5–10 keV), is actually very similar to the others, according to the results of our re-analysis (burst D). The above consistencies, together with the $`56`$ d burst wait times observed by both Hakucho and BeppoSAX, suggest the burst characteristics of the binary 4U 1812$``$12 to be remarkably stable in observations spanning $`15`$ years apart.
###### Acknowledgements.
We thank the staff of the BeppoSAX Science Operation Centre and Science Data Centre for their help in carrying out and processing the WFC Galactic Centre observations. The BeppoSAX satellite is a joint Italian and Dutch program. M.C., A.B., L.N. and P.U. thank Agenzia Spaziale Nazionale (ASI) for grant support. M.C. also thanks M. Federici (IAS) for technical support.
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# 1) Introduction
## 1) Introduction
The description of baryonic matter on the basis of QCD remains a theoretical challenge, especially since lattice gauge calculations have so far been of little help for this problem. Nuclear physics, relativistic heavy-ion physics and astrophysics are some of the fields which would greatly benefit from any progress on this issue. Recently, a new surge of interest has been triggered by the suggestion that at high density, the novel phenomenon of color superconductivity might set in . This development has highlighted how little is known reliably about dense, strongly interacting matter.
Here, we address a much simpler finite density problem where a full analytic solution can be found: We consider two-dimensional model field theories with interacting fermions at or near the chiral limit. Specifically, we have in mind the two-dimensional version of the Nambu–Jona-Lasinio model , i.e., the chiral Gross-Neveu model , and QCD<sub>2</sub> with fundamental quarks, the ’t Hooft model . In both cases, one considers a large number $`N`$ of fermion species and investigates the limit $`N\mathrm{}`$, keeping $`Ng^2`$ fixed . These two models are quite similar as far as their chiral properties are concerned but differ with respect to confinement of quarks which is only exhibited by the ’t Hooft model. So far, the phase diagram of the Gross-Neveu model has been studied extensively as function of temperature and chemical potential , and the results seem to be uncontroversial. The ’t Hooft model on the other hand has been investigated only sporadically at finite temperature , most recently in Ref. , but hardly anything is known yet about its properties at finite density .
Our point of departure is the following observation: Both of these models possess light baryons whose mass vanishes in the chiral limit (by light, we mean that $`M_B/N`$ is small on the relevant physical scale). This is of course no accident, but a generic feature of models with broken chiral symmetry in 1+1 dimension — the baryons are topologically non-trivial excitations of the Goldstone boson (“pion”) field. By contrast, the Gross-Neveu model with discrete chiral symmetry can only accommodate baryons whose mass scales with the physical fermion mass and hence stays finite in the chiral limit. Nevertheless, it has been argued that both variants of the Gross-Neveu model have identical phase diagrams . Since we find it rather perplexing that the structure of the single baryon should have no influence on the structure of baryonic matter, we have reinvestigated this issue. We have found that a combination of large $`N`$ techniques with the strong constraints arising from broken chiral symmetry is powerful enough to allow for a simple, analytic solution of this problem in the vicinity of the chiral limit. The results of our analysis differ qualitatively from the conventional wisdom about the Gross-Neveu model and carry over to the (confining) ’t Hooft model as well. They seem to confirm a number of investigations of other field theories at finite fermion density, where strikingly similar behaviour was found. These include exact studies of two-dimensional models like the massive and massless Schwinger model or finite $`N_c`$ and $`N_f`$ massless QCD<sub>2</sub> as well as more approximate treatments of 4-dimensional large $`N`$ QCD and effective chiral quark models .
A word of caution is in order here: Throughout this paper, we shall constantly deal with spontaneous breaking of continuous symmetries and Goldstone bosons in 1+1 dimensions, in seeming conflict with the Coleman-Mermin-Wagner theorem . It is well understood by now that the large $`N`$ limit enables one to circumvent this no-go theorem. As clearly explained by Witten , the bad infrared behaviour of the boson propagator, when exponentiated, gives rise to a power law correlator $`|xy|^{1/N}`$ which becomes constant in the limit $`N\mathrm{}`$. This kind of almost long-range order is also familiar from the two dimensional $`XY`$-model . Alternatively, one may argue that the mean field approximation predicts symmetry breakdown and that this result is protected against fluctuations (which would otherwise restore the symmetry in two dimensions) by $`1/N`$ suppression factors. In this sense, low dimensional large $`N`$ theories are not only more tractable, but also physically more appealing than their finite $`N`$ counterparts. They bear more resemblance to the real, 3+1 dimensional world.
This paper is organized as follows. In Sect. 2 we briefly review the conventional analytical treatment of the chiral Gross-Neveu model at finite density and point out a certain deficiency of this approach. In Sect. 3 we repeat a similar analysis for the ’t Hooft model, supplementing the analytical methods by numerical computations where necessary. In Sect. 4 the Skyrme type of approach to the light baryons in both models is recalled and generalized to the case of baryonic matter in the strict chiral limit. In Sect. 5 we then allow for a small symmetry breaking mass term and make contact with the sine-Gordon kink chain, the two dimensional analog of the Skyrme crystal . This is followed by a short summary and conclusions in Sect. 6.
## 2) Chiral Gross-Neveu model at finite density: Conventional approach
Let us first recall the standard treatment of the Gross-Neveu model at finite density. In the large $`N`$ limit mean-field techniques become exact. Technically, they may be phrased in a variety of ways. We choose the language of relativistic many-body theory, following Refs. , which we find particularly intuitive for the problem at hand. Then the vacuum, the baryon and baryonic matter are all described by a relativistic Hartree-Fock approach (for baryons in the large $`N`$ limit this was first recognized in Ref. ). “Conventional approach” in the title of this section refers to translational invariance — we shall assume that the system is described by an interacting Fermi gas with prescribed, homogeneous density. We shall first deal with the Gross-Neveu model with continuous chiral symmetry ($`\psi \mathrm{e}^{\mathrm{i}\alpha \gamma _5}\psi `$) and Lagrangian density
$$=\overline{\psi }\mathrm{i}/\psi +\frac{1}{2}g^2\left((\overline{\psi }\psi )^2+(\overline{\psi }\mathrm{i}\gamma _5\psi )^2\right).$$
(1)
As a matter of fact, the corresponding calculation would be identical for the model with discrete chiral symmetry only ($`\psi \gamma _5\psi `$), where the $`\gamma _5`$-term in Eq. (1) is omitted. The results presented here are well known, but our aim is to criticize them in a novel way.
We denote the fermion density per color (or baryon density) by $`\rho _B=p_f/\pi `$ ($`p_f`$: Fermi momentum). At the mean field level, the fermions acquire a physical mass $`m`$ which has to be determined self-consistently. The ground state energy density per color is given by
$$\frac{}{N}=2_{p_f}^{\mathrm{\Lambda }/2}\frac{\mathrm{d}k}{2\pi }\sqrt{m^2+k^2}+\frac{m^2}{2Ng^2}$$
(2)
where $`\mathrm{\Lambda }`$ is an ultra-violet cutoff. The first term is just the sum over single particle energies for all occupied states (the Dirac sea plus all positive energy states with $`|p|<p_f`$), the second term the usual correction for double counting of interaction effects familiar from the Hartree-Fock approximation. To renormalize the theory, let us first consider the limit $`p_f0`$, denoting the physical fermion mass in the vacuum by $`m_0`$,
$$\frac{}{N}=2_0^{\mathrm{\Lambda }/2}\frac{\mathrm{d}k}{2\pi }\sqrt{m_0^2+k^2}+\frac{m_0^2}{2Ng^2}.$$
(3)
Minimizing $``$ with respect to $`m_0`$ yields the relativistic Hartree-Fock equation
$$m_0\left(1+\frac{Ng^2}{\pi }\mathrm{ln}\frac{m_0}{\mathrm{\Lambda }}\right)=0.$$
(4)
Due to the similarity in structure between the relativistic Hartree-Fock approach and BCS theory , this is often referred to as the “gap equation”. The non-trivial solution (which has always lower vacuum energy) yields the relation
$$\frac{Ng^2}{\pi }\mathrm{ln}\frac{\mathrm{\Lambda }}{m_0}=1$$
(5)
which teaches us how the bare coupling constant depends on the cutoff parameter, given $`m_0`$. Recall that the Gross-Neveu model shares with real QCD both asymptotic freedom and dimensional transmutation; these properties are contained in Eq. (5). Using this relation to renormalize the matter ground state energy density, Eq. (2), we find (dropping an irrelevant term $`\mathrm{\Lambda }^2/8\pi `$)
$$\frac{}{N}=\frac{m^2}{4\pi }+\frac{1}{2\pi }p_f\sqrt{p_f^2+m^2}+\frac{1}{2\pi }m^2\mathrm{ln}\left(\frac{p_f+\sqrt{m^2+p_f^2}}{m_0}\right).$$
(6)
The energy is minimal provided $`m`$ satisfies
$$m\mathrm{ln}\left(\frac{p_f+\sqrt{m^2+p_f^2}}{m_0}\right)=0,$$
(7)
i.e., for
$$m=0\text{or}m=m_0\sqrt{1\frac{2p_f}{m_0}}\left(p_f<\frac{m_0}{2}\right).$$
(8)
The corresponding energy densities are
$`{\displaystyle \frac{}{N}}|_{m=0}`$ $`=`$ $`{\displaystyle \frac{p_f^2}{2\pi }},`$
$`{\displaystyle \frac{}{N}}|_{m0}`$ $`=`$ $`{\displaystyle \frac{m_0^2}{4\pi }}+{\displaystyle \frac{p_fm_0}{\pi }}{\displaystyle \frac{p_f^2}{2\pi }}\left(p_f<{\displaystyle \frac{m_0}{2}}\right).`$ (9)
The physical quark masses (8) and the energy densities (9) are plotted in Figs. 1 and 2. From these figures one might be tempted to conclude that chiral symmetry is broken at low densities and gets restored in a second order phase transition at $`p_f=m_0/2`$. As is well known, this does not occur, rather there is a first order chiral phase transition at $`p_f=m_0/\sqrt{2}`$. This can easily be inferred by inspection of the thermodynamic potential of the Gross-Neveu model . For our purpose, the following physical reasoning is perhaps more instructive: Let us compare the energy densities (9) with the energy density for a system of size $`L`$ divided into two homogeneous regions I (size $`\mathrm{}`$) and II (size $`L\mathrm{}`$). In region I chiral symmetry is restored; it contains the extra fermions needed to get the prescribed average density (the “MIT bag” ). Region II consists of the physical vacuum with broken chiral symmetry, void of excess fermions. The mean energy density obtained in this way is
$$\frac{}{N}=\left(\frac{L\mathrm{}}{L}\right)\frac{m_0^2}{4\pi }+\frac{Lp_f^2}{2\pi \mathrm{}}.$$
(10)
Minimization with respect to $`\mathrm{}`$ yields
$$\mathrm{}=\frac{\sqrt{2}p_fL}{m_0}$$
(11)
valid for $`p_f<m_0/\sqrt{2}`$, and hence the optimal energy density
$$\frac{}{N}=\frac{m_0^2}{4\pi }+\frac{p_fm_0}{\sqrt{2}\pi }\left(p_f<\frac{m_0}{\sqrt{2}}\right).$$
(12)
As shown in Fig. 2, this solution is lower in energy than the homogeneous one; moreover, it yields the convex hull of $``$. It ends exactly at the first order phase transition point $`p_f=m_0/\sqrt{2}`$ where all space is filled with one big bag. This should be contrasted to the scenario underlying Fig. 1 where the fermion mass decreases continuously. We thus recover the generally accepted mixed phase interpretation of the Gross-Neveu model at finite density. Notice also that only the total size of regions I and II matters, not how they are subdivided; there could be baryon “droplets” as well. Alternatively, the mixed phase curve in Fig. 2 with its linear dependence on $`p_f`$ could have been inferred from a standard Maxwell construction. It is interesting that a very similar qualitative behaviour was found recently in 3+1 dimensions, where the close relationship with the bag model was also stressed .
One important point to which we would like to draw the attention of the reader is the behaviour of $``$ near $`\rho _B=p_f/\pi =0`$. Since ultimately, at very low density, the fermionic matter problem must reduce to the problem of a single baryon, one would expect
$$\frac{}{\rho _B}|_{\rho _B=0}=M_B$$
(13)
where $`M_B`$ is the baryon mass. In the present calculation, $`M_B`$ is not the physical baryon mass, but the mass of an alleged “delocalized” baryon. This is inherent in the translationally invariant Hartree-Fock approach, i.e., the assumption that the single particle orbitals are momentum eigenstates. Using Eq. (13) we obtain in the homogeneous, single phase calculation, Eq. (9), $`M_B=Nm_0`$, consistent with a short range force and a delocalized baryon. The (physically more viable) mixed phase approach, Eq. (12), predicts a baryon mass lower by a factor of $`1/\sqrt{2}`$. This factor can readily be understood in terms of the bag model. Indeed, it follows from Eq. (10) that ($`E=L`$)
$$E_BE_0=N\mathrm{}\left(\frac{m_0^2}{4\pi }+\frac{q_f^2}{2\pi }\right),q_f=\frac{\pi B}{\mathrm{}}.$$
(14)
For $`B=1`$, this expression can be interpreted as energy of a single baryon, $`\mathrm{}`$ being its diameter. The first term is just the bag pressure (the difference between the energy density of the physical vacuum and that of the perturbative one, cf. Eq. (9) for $`p_f=0`$), the second the kinetic energy of $`N`$ massless quarks. The bag size $`\mathrm{}`$ is found through minimization of the energy (for $`B=1`$) to be
$$\mathrm{}=\frac{\sqrt{2}\pi }{m_0}.$$
(15)
Inserting this result into Eq. (14), one finds that the bag pressure and the quark kinetic energy contributions are exactly equal in this model and that $`M_B=Nm_0/\sqrt{2}`$.
However, the Gross-Neveu model possesses bound baryons with lowest mass $`Nm_0/\pi `$ (kink solution for the model with discrete chiral symmetry ), or even massless baryons (model with continuous chiral symmetry ). These binding effects are not $`1/N`$ suppressed and should be correctly reproduced in a Hartree-Fock approach, in the low density limit. They have obviously been missed here due to our tacit assumption of translational invariance. There is no good reason why such effects should not play a role at higher densities as well. Moreover, differences between the continuous and discrete chirally symmetric Gross-Neveu models based on their different baryon structure and masses are not at all captured by the “conventional” approach. Below, we shall present a cure for this disease. Before that however, let us first repeat the naive calculation for the ’t Hooft model, where the corresponding results are not yet available in the literature.
## 3) ’t Hooft model at finite density, assuming translational invariance
The ’t Hooft model is defined as the large $`N`$ limit of 1+1 dimensional SU($`N`$) gauge theory with quarks in the fundamental representation ,
$$=\overline{\psi }\mathrm{i}D/\psi \frac{1}{2}\mathrm{tr}F_{\mu \nu }F^{\mu \nu }.$$
(16)
Since the light-cone approach originally used by ’t Hooft to determine the meson spectrum seems to be less convenient for the vacuum, baryon and baryonic matter problems, we shall work in normal coordinates. This approach was pioneered by Bars and Green and further developed in Refs. . Common to all of these works is the fact that the gluons are gauged away (axial gauge), leaving behind a theory of fermions interacting via a linear Coulomb potential. We refer the reader to the detailed derivation of the Hartree-Fock approach in Refs. and immediately proceed to the formulae which are relevant for our purpose. Let us first summarize the treatment of the vacuum. A central quantity is the single particle density matrix in momentum space,
$$\rho (p)=\frac{1}{2}+\gamma ^0\rho _0(p)\mathrm{i}\gamma ^1\rho _1(p)+\gamma ^5\rho _5(p).$$
(17)
Its precise definition in terms of the quark fields is
$$\rho _{\alpha \beta }(p)=dx\mathrm{e}^{\mathrm{i}px}0|\frac{1}{N}\underset{i}{}\psi _{i\beta }^{}(0)\psi _{i\alpha }(x)|0.$$
(18)
The Slater determinant condition characteristic for Hartree-Fock,
$$\rho ^2(p)=\rho (p)\rho _0^2(p)+\rho _1^2(p)+\rho _5^2(p)=1/4$$
(19)
holds manifestly in the parametrization
$$\left(\begin{array}{c}\rho _0(p)\\ \rho _1(p)\\ \rho _5(p)\end{array}\right)=\frac{1}{2}\left(\begin{array}{c}\mathrm{sin}\theta (p)\mathrm{cos}\phi \\ \mathrm{sin}\theta (p)\mathrm{sin}\phi \\ \mathrm{cos}\theta (p)\end{array}\right).$$
(20)
$`\theta (p)`$ is the Bogoliubov angle, $`\phi `$ the global angle which locates the broken symmetry vacuum on the chiral circle (hence it has no $`p`$-dependence). The Bogoliubov angles are the variational parameters; this terminology stems once more from BCS theory, which has the same formal structure as relativistic Hartree-Fock. We choose $`\phi =0`$, the value reached if one lets the bare quark mass approach zero starting from a finite value. Then,
$$\rho (p)=v(p)v^{}(p)$$
(21)
with the (positive and negative energy) Hartree-Fock spinors
$$u(p)=\left(\begin{array}{c}\hfill \mathrm{cos}\theta (p)/2\\ \hfill \mathrm{sin}\theta (p)/2\end{array}\right),v(p)=\left(\begin{array}{c}\hfill \mathrm{sin}\theta (p)/2\\ \hfill \mathrm{cos}\theta (p)/2\end{array}\right).$$
(22)
The vacuum expectation value of the Hamiltonian density reads
$$\frac{}{N}=\frac{\mathrm{d}p}{2\pi }p\mathrm{cos}\theta (p)\frac{Ng^2}{8}\frac{\mathrm{d}p}{2\pi }\frac{\mathrm{d}p^{}}{2\pi }\frac{\mathrm{cos}(\theta (p)\theta (p^{}))1}{(pp^{})^2}$$
(23)
where the first term is the kinetic energy, the second the Coulomb interaction of the quarks. Varying with respect to the Bogoliubov angles $`\theta (p)`$, the gap equation is obtained in the form
$$p\mathrm{sin}\theta (p)+\frac{Ng^2}{4}\frac{\mathrm{d}p^{}}{2\pi }\frac{\mathrm{sin}\left(\theta (p)\theta (p^{})\right)}{(pp^{})^2}=0.$$
(24)
The integral has to be defined as principal value integral, cf. Refs. . We shall also need the expression for the quark condensate in the vacuum,
$$\overline{\psi }\psi _\mathrm{v}=N\frac{\mathrm{d}p}{2\pi }\mathrm{sin}\theta (p),$$
(25)
and the quark single particle energies,
$$\omega (p)=p\mathrm{cos}\theta (p)+\frac{Ng^2}{4}\frac{\mathrm{d}p^{}}{2\pi }\frac{\mathrm{cos}\left(\theta (p)\theta (p^{})\right)}{(pp^{})^2}.$$
(26)
Although the gap equation (24) for the ’t Hooft model must be solved numerically, the value of the quark condensate (25) is known analytically, owing to an indirect determination via sum rules and the ’t Hooft equation for mesons ; it is
$$\overline{\psi }\psi _\mathrm{v}=\frac{N}{\sqrt{12}}\left(\frac{Ng^2}{2\pi }\right)^{1/2}.$$
(27)
The single particle energies (26) are badly infrared divergent, a source of a long and ongoing debate in the literature . To exhibit the divergence, we follow Ref. , isolate the divergent part of the integral and regularize it by using a finite box of length $`L`$,
$$\omega (p)=p\mathrm{cos}\theta (p)+\frac{Ng^2}{4}\frac{\mathrm{d}p^{}}{2\pi }\frac{\mathrm{cos}\left(\theta (p)\theta (p^{})\right)1}{(pp^{})^2}+\frac{Ng^2L}{48}.$$
(28)
This last constant diverges for $`L\mathrm{}`$ but seems to be essential to account for confinement in such an independent particle picture: The isolated quarks behave roughly as if they had infinite mass. If one simply throws the infinite constant away, as is often done, one gets an awkward sign change in $`\omega (p)`$ at some low momentum $`p`$, cf. Refs. , and runs into serious inconsistencies in finite temperature Hartree-Fock calculations . Fortunately, the constant drops out of the calculation of color singlet mesons, as already noticed by ’t Hooft (his IR cutoff parameter $`\lambda `$ is related to our box size $`L`$ by $`\lambda =12/\pi L`$, cf. Ref. ): The infinite self-energy term is cancelled by an equally infinite piece in the Coulomb interaction. We also note in passing that if one employs a finite box as infrared regulator, one is unambiguously led to ’t Hooft’s treatment of the quark self-energies rather than to Wu’s alternative regularization prescription . Since the emergence of the constant $`Ng^2L/48`$ in the single particle energy (28), but not in the vacuum energy (23), is rather important for our discussion and somewhat hidden in Ref. , we have included a simplified version of the arguments underlying Eqs. (23) and (28) in the appendix.
After this review of the treatment of the vacuum, we are in a position to include a finite baryon density, assuming translational invariance. If $`p_f`$ denotes the Fermi momentum, we have to replace the density matrix (21) by
$`\rho (p)`$ $`=`$ $`\mathrm{\Theta }(p_f|p|)u(p)u^{}(p)+v(p)v^{}(p)`$ (29)
$`=`$ $`\mathrm{\Theta }(p_f|p|)+\mathrm{\Theta }(|p|p_f)v(p)v^{}(p)`$
where we have used the completeness relation for the spinors in the second step. In the expression for the Hartree-Fock ground state energy density (23), according to the second line of Eq. (29), we must exclude the region $`[p_f,p_f]`$ from the momentum integrations and pick up an additional term due to the change in the baryon density $`\mathrm{tr}\rho `$,
$`{\displaystyle \frac{}{N}}`$ $`=`$ $`{\displaystyle \frac{\mathrm{d}p}{2\pi }\mathrm{\Theta }(|p|p_f)p\mathrm{cos}\theta (p)}`$ (30)
$`{\displaystyle \frac{Ng^2}{8}}{\displaystyle \frac{\mathrm{d}p}{2\pi }\frac{\mathrm{d}p^{}}{2\pi }\mathrm{\Theta }(|p|p_f)\mathrm{\Theta }(|p^{}|p_f)\frac{\mathrm{cos}(\theta (p)\theta (p^{}))1}{(pp^{})^2}}`$
$`+{\displaystyle \frac{Ng^2}{4}}{\displaystyle \frac{\mathrm{d}p}{2\pi }\frac{\mathrm{d}p^{}}{2\pi }\mathrm{\Theta }(p_f|p|)\mathrm{\Theta }(|p^{}|p_f)\frac{1}{(pp^{})^2}}.`$
This yields at once the following finite density generalization of the gap equation,
$$p\mathrm{sin}\theta (p)+\frac{Ng^2}{4}\frac{\mathrm{d}p^{}}{2\pi }\mathrm{\Theta }(|p^{}|p_f)\frac{\mathrm{sin}\left(\theta (p)\theta (p^{})\right)}{(pp^{})^2}=0,(|p|>p_f),$$
(31)
whereas the condensate now becomes
$$\overline{\psi }\psi =N\frac{\mathrm{d}p}{2\pi }\theta (|p|p_f)\mathrm{sin}\theta (p).$$
(32)
The gap equation (31) can easily be solved numerically for various $`p_f`$. The resulting condensate is shown in Fig. 3. We find that it decreases monotonically with increasing density, disappearing at a critical Fermi momentum
$$p_f^c0.117\left(\frac{Ng^2}{2\pi }\right)^{1/2}.$$
(33)
This behaviour is strikingly similar to the corresponding result for the Gross-Neveu model depicted in Fig. 1, again suggesting some phase transition with restoration of chiral symmetry at high density. Here we are not able to go on and discuss whether we are dealing with a first or second order phase transition. The reason lies in the following problem: If we compute the energy density (30) for the ’t Hooft model we discover that subtraction of the value at $`p_f=0`$ is not sufficient to give a finite result. Unlike in the Gross-Neveu model, the difference is still IR divergent. To be able to proceed, we enclose the system once more in a box of length $`L`$. We then find that the divergence is due to the last term in Eq. (30) (the one which does not involve the Bogoliubov angles) which now contributes the following double sum to the energy per color,
$$\frac{E}{N}|_{\mathrm{div}}=\frac{Ng^2L}{16\pi ^2}\underset{pI}{}\underset{n0,(pn)I}{}\frac{1}{n^2}.$$
(34)
Here antiperiodic boundary conditions for fermions have been employed in the box regularization, and correspondingly the interval $`I`$ is defined in the following way,
$$I=[n_f,n_f]\text{for}B=2n_f+1\text{odd},I=[n_f1,n_f]\text{for}B=2n_f+2\text{even}.$$
(35)
The result (34) is even more alarming than the non-convex behaviour of $``$ in the Gross-Neveu model, Fig. 2, due to its $`L`$-dependence. Adding quarks to the vacuum causes the energy to increase by an infinite amount in the limit $`L\mathrm{}`$. Evaluating the double sums in Eq. (34) for low values of $`B`$, we obtain information on the origin of this divergent behaviour. For $`B=1`$ ($`I=[0,0]`$) in particular, the calculated baryon mass (to leading order in $`L`$) is
$$M_B=N\left(\frac{Ng^2L}{48}\right).$$
(36)
This is the same relation as $`M_B=Nm_0`$ in the Gross-Neveu model except that the physical fermion mass is replaced by the infinite constant $`Ng^2L/48`$ characteristic of confinement, cf. Eq. (28). For larger values of $`B`$ Eq. (34) does not simply yield multiples of the baryon mass (36), but one finds “interaction effects” of the same order of magnitude as the mass. As far as $`N`$-counting is concerned, this is still in agreement with Witten’s analysis of baryons at large $`N`$ . However, since these delocalized baryons are presumably not very physical in the ’t Hooft model, we refrain from further discussing these effects.
Summarizing, the problems encountered in the Gross-Neveu model with translationally invariant baryonic matter again show up in the ’t Hooft model, although in a much more severe form. The physics reason is clear: In the Gross-Neveu model the cost of distributing $`N`$ fermions over the whole space is governed by their physical mass; in the ’t Hooft model, due to confinement of quarks, the corresponding quark effective mass diverges with the volume. On the other hand, it is known that both models do possess massless, delocalized baryons in the chiral limit. Evidently, this has to be accounted for, and we conclude that the naive, translationally invariant Hartree-Fock approximation fails miserably in describing the properties of baryonic matter.
## 4) Massless baryons and baryonic matter in the Skyrme picture
The existence of massless baryons in the chiral limit of the ’t Hooft model has been demonstrated by bosonization , variational , and light-cone techniques. These exotic objects are characteristic for 1+1 dimensional models with broken chiral symmetry and, as such, also present in the chiral Gross-Neveu model. A particularly illuminating derivation is due to Salcedo et al. . These authors point out that the potential energy in such models is invariant under local chiral transformations, unlike the kinetic term which is only invariant under global ones. This led them to the following variational ansatz for the one-body density matrix of the baryon,
$$\rho (x,y)=\mathrm{e}^{\mathrm{i}\chi (x)\gamma _5}\rho _\mathrm{v}(xy)\mathrm{e}^{\mathrm{i}\chi (y)\gamma _5}.$$
(37)
Here $`\rho _\mathrm{v}(xy)`$ is the vacuum density matrix. If the vacuum breaks chiral symmetry, one can generate with expression (37) a new (exact or approximate) Hartree-Fock solution which breaks translational invariance but can carry non-zero baryon number. As shown in , the baryon density is given by
$$\rho _B(x)=\mathrm{tr}(\rho (x,x)\rho _\mathrm{v}(0))=\frac{1}{\pi }_x\chi (x),$$
(38)
so that the baryon number coincides with the winding number of the chiral phase $`\chi (x)`$,
$$B=_0^Ldx\frac{1}{\pi }_x\chi =\frac{\chi (L)\chi (0)}{\pi }\mathrm{𝖹𝖹}.$$
(39)
(Notice that $`\chi (L)\chi (0)`$ must be an integer multiple of $`\pi `$ since otherwise bilinear fermion observables would no longer be periodic.) For this topological reasoning it is again recommendable to work in a finite box of size $`L`$. The topological interpretation of the baryon number also agrees with exact results of Ref. which were not restricted to the large $`N`$ limit. In the absence of an explicit quark mass term, the ground state energy obtained from Eq. (37) is
$$E[\rho ]=E[\rho _\mathrm{v}]+N_0^Ldx\frac{1}{2\pi }(_x\chi )^2.$$
(40)
This result holds independently of the specific model, since the potential energy does not contribute to $`E[\rho ]E[\rho _\mathrm{v}]`$. Differences between various models are of course still present in the vacuum density matrix $`\rho _\mathrm{v}`$ in Eq. (37) but do not manifest themselves in the baryon energy. Minimizing $`E[\rho ]`$ with respect to $`\chi `$ yields the free (static) bosonic equation
$$_x^2\chi (x)=0,\chi (L)=\chi (0)+\pi B$$
(41)
with the solution
$$\chi (x)=\pi B\left(\frac{xx_0}{L}\right).$$
(42)
Here $`x_0`$ is a parameter which reflects the breakdown of translational invariance. The baryon density is $`x`$-independent ($`\rho _B(x)=B/L`$) as follows more generally from axial current conservation in the chiral limit . However, the scalar and pseudoscalar condensates acquire a non-trivial $`x`$-dependence,
$`\overline{\psi }\psi `$ $`=`$ $`\overline{\psi }\psi _\mathrm{v}\mathrm{cos}\left(2\pi B(xx_0)/L\right),`$
$`\overline{\psi }\mathrm{i}\gamma _5\psi `$ $`=`$ $`\overline{\psi }\psi _\mathrm{v}\mathrm{sin}\left(2\pi B(xx_0)/L\right).`$ (43)
Since fluctuations of $`\chi (x)`$ describe the massless Goldstone boson field, the baryon picture emerging here is very similar in spirit to the Skyrme model . The fact that the baryon is a topological soliton will become somewhat clearer once we include a small bare quark mass (see Sect. 5) but this solitonic character also holds in the strict chiral limit considered here.
We can now discuss the baryon as well as baryonic matter from this point of view. The single baryon ($`B=1`$) is spread out over the whole space, the chiral phase $`\chi (x)`$ making one turn with constant speed to minimize the kinetic energy (Fig. 4). The baryon energy is, using Eqs. (40-42),
$$E_B=N\frac{\pi }{2L}.$$
(44)
This confirms that indeed the baryon becomes massless in the limit $`L\mathrm{}`$. Incidentally, expression (44) is identical to the kinetic energy of $`N`$ non-interacting, massless quarks in the lowest momentum state available for antiperiodic boundary conditions. Nevertheless, we are not dealing with the free, chirally symmetric theory, but with the broken phase of an interacting theory where the quarks are massive or even confined.
It may be worthwhile to contemplate the structure of the baryon for a moment from the point of view of the relativistic Hartree-Fock approximation. In chirally non-invariant models one would suspect that the baryon comprises the filled Dirac sea plus one filled, positive energy valence level. This is exactly what one finds analytically in the non-chiral Gross-Neveu model , or numerically in QCD<sub>2</sub> with heavy quarks . The picture implied by the ansatz (37) in the chiral limit is rather different though. Denoting the negative energy single particle orbitals in the Dirac sea by $`\phi _k^{()}(x)`$ (solutions of the first quantized Dirac equation with Hartree-Fock potential), the Skyrme type baryon (37) admits the density matrix
$$\rho (x,y)=\mathrm{e}^{\mathrm{i}\pi x\gamma _5/L}\underset{k}{}\phi _k^{()}(x)\phi _k^{()}(y)\mathrm{e}^{\mathrm{i}\pi y\gamma _5/L}$$
(45)
where
$$\phi _k^{()}(x)=\mathrm{e}^{\mathrm{i}kx}v(k),$$
(46)
and the $`k`$ are discrete momenta appropriate to the interval of length $`L`$. We first observe that the chiral phase factor splits the momenta of the right- and left-handed components into $`k\pm \pi /L`$. Since the transformed single particle wave functions are no longer momentum eigenstates, translational invariance is lost. Secondly, we note that the presence of an infinite Dirac sea is crucial for getting the extra baryon charge, rather than a single valence state. If the sum over occupied states $`k`$ in Eq. (45) was finite, we would trivially conclude that $`\rho (x,y)`$ and $`\rho _\mathrm{v}(xy)`$ belong to the same baryon density ($`\rho _B(x)=\mathrm{tr}\rho (x,x)`$). Due to the infinite number of occupied states however, $`\rho _\mathrm{v}(xy)`$ develops a singularity at $`x=y`$, and one has to do a more careful point splitting in order to compute the baryon density. The divergence is due to the UV region and therefore determined by the free theory (for more details, cf. Ref. ),
$$\underset{xy}{lim}\mathrm{tr}(\rho (x,y)\rho _\mathrm{v}(xy))=\underset{z0}{lim}\mathrm{tr}\left\{\mathrm{e}^{\mathrm{i}\pi z\gamma _5/L}\left(\frac{1}{2}\delta (z)\frac{\mathrm{i}\gamma _5}{2\pi z}\right)\frac{1}{2}\delta (z)\right\}=\frac{1}{L}.$$
(47)
The result $`1/L`$ is the baryon density for $`B=1`$. This mechanism is strongly reminiscent of the calculation of anomalous current commutators, for instance in the Schwinger model. The extra baryon number does not reside in a valence level added on top of the Dirac sea but somehow emerges from the bottom of the Dirac sea if one modifies all the levels slightly — it is a vacuum polarization effect.
Equipped with this exotic kind of baryon, we can now easily find the ground state of the system for any baryon density. As discussed above and illustrated in Fig. 4, the single baryon consists of one turn of a “chiral spiral” (parametrized by $`\chi (x)`$) over the total spatial length $`L`$ of the system — admittedly a somewhat elusive object in the thermodynamic limit. A finite density $`\rho _B=B/L=p_f/\pi `$ on the other hand implies that
$$\chi (x)=p_f(xx_0),$$
(48)
i.e., one full rotation over a physical distance which has a well defined limit for $`L\mathrm{}`$, namely 2/$`\rho _B`$. The baryon density remains constant in space, but the condensates are modulated as
$`\overline{\psi }\psi `$ $`=`$ $`\overline{\psi }\psi _\mathrm{v}\mathrm{cos}2p_f(xx_0),`$
$`\overline{\psi }\mathrm{i}\gamma _5\psi `$ $`=`$ $`\overline{\psi }\psi _\mathrm{v}\mathrm{sin}2p_f(xx_0).`$ (49)
They can be viewed as projections of a “chiral spiral” of radius $`|\overline{\psi }\psi _\mathrm{v}|`$ onto two orthogonal planes, see Fig. 5. This state breaks translational symmetry; it is a crystal. In fact, it may be viewed as the simplest possible realization of the old idea of a Skyrme crystal , here in the context of large $`N`$ two-dimensional field theories. One cannot tell where one baryon begins and ends — each full turn of the spiral contains baryon number 1. Only the condensates reveal that translational symmetry has been broken down to a discrete subgroup. The energy density of this unusual kind of “nuclear matter” is simply (after subtracting the vacuum energy density)
$$\frac{}{N}=\frac{p_f^2}{2\pi }.$$
(50)
Surprisingly, this is exactly what one would expect for a free Fermi gas of massless quarks although Eq. (50) holds for interacting theories where the vacuum has lower energy due to chiral symmetry breaking. In Fig. 6 we compare the energy density for this state to the ones discussed above for the Gross-Neveu model, where translational symmetry had been assumed. The crystal is always energetically favored, the dependence on $`p_f`$ is now convex, and there is no trace of a phase transition, neither first nor second order, at any density. The horizontal slope at $`p_f=0`$ correctly signals the presence of massless baryons and eliminates the above mentioned problems with the spurious massive, delocalized baryons. We cannot even draw the corresponding picture for the ’t Hooft model, simply because in this case the quark Fermi gas is infinitely higher in energy than the Skyrme crystal for $`L\mathrm{}`$. Nevertheless, all the results for baryonic matter discussed in this section apply to the ’t Hooft model as well.
In the high density limit the oscillations of the condensates become more and more rapid. If we are interested only in length scales large as compared to $`1/p_f`$, the condensates average to zero. In this sense, one might argue that chiral symmetry gets restored at high density, although not in the naive way suggested by Fig. 2.
Finally, we remark that the “chiral spiral” ground state for fixed baryon density still preserves one continuous, unbroken symmetry, namely the combination of translation and chiral rotation generated by $`P+p_fQ_5`$ ($`P`$: momentum operator, $`Q_5`$: axial charge). One would therefore predict that RPA excitations on this ground state (or mesons in nuclear matter) will have only one collective, gapless mode, a hybrid of a “phonon” and a “pion”.
## 5) Non-vanishing bare quark masses
In Ref. the Skyrme picture of the baryons in the ’t Hooft model and chiral Gross-Neveu model was developed for small, finite bare quark masses, using the expression in Eq. (37) as a variational ansatz. For a single baryon, these authors have tested the accuracy of their procedure against the full, numerical Hartree-Fock calculation on a lattice. The results agreed perfectly at $`m_q=0.05`$ and were still surprisingly good at $`m_q=0.20`$, in units of $`\sqrt{Ng^2/2\pi }`$. This makes it very tempting to speculate that the corresponding variational calculation can also give us a reliable picture of baryonic matter at finite density, away from the chiral limit. As compared to the formulae in the preceding section, the only change is the fact that the bare mass term now also contributes to the energy functional Eq. (40),
$$E[\rho ]=E[\rho _\mathrm{v}]+N_0^Ldx\left\{\frac{1}{2\pi }\left(_x\chi \right)^2+\frac{m_q\overline{\psi }\psi _\mathrm{v}}{N}\left(\mathrm{cos}2\chi 1\right)\right\}.$$
(51)
Here the condensate $`\overline{\psi }\psi _\mathrm{v}`$ refers to the vacuum in the chiral limit. Varying with respect to $`\chi (x)`$ then gives the static sine-Gordon equation ,
$$_x^2\chi +\frac{2\pi m_q\overline{\psi }\psi _\mathrm{v}}{N}\mathrm{sin}2\chi =0,$$
(52)
from which one reads off the “pion” mass (two dimensional version of the Gell Mann-Oakes-Renner relation )
$$m_\pi ^2=4\pi m_q\frac{\overline{\psi }\psi _\mathrm{v}}{N}.$$
(53)
The $`B=1`$ baryon can be identified with the familiar kink solution of the sine-Gordon equation,
$$\chi (x)=2\mathrm{arctan}\left(\mathrm{e}^{m_\pi (xx_0)}\right),$$
(54)
with mass
$$M_B=N\frac{2m_\pi }{\pi }.$$
(55)
Since the single baryon has been discussed in detail in Ref. , let us immediately turn to multi-kink solutions as candidates for baryonic matter. Luckily, the sine-Gordon kink crystal has already been studied thoroughly in the literature, first in solid state physics and more recently as a toy model for the Skyrme crystal , in terms of Jacobi elliptic functions and elliptic integrals . We take over the results from Ref. which is close in spirit to the present study although the authors did not have in mind two-dimensional large $`N`$ field theories. Adapting the formulae of this work to our notation, the following steps allow us to generalize the Skyrme crystal of the previous section to small, finite bare quark masses: Let $`m_\pi `$ denote the mass of the Goldstone boson, Eq. (53), and $`\overline{\rho }_B=p_f/\pi `$ the average baryon density (this is our definition of $`p_f`$ for the case of broken translational symmetry). We then first have to solve the transcendental equation
$$\frac{\pi m_\pi }{p_f}=2k𝐊(k)$$
(56)
for $`k`$ where $`𝐊(k)`$ is the complete elliptic integral of the first kind. The sine-Gordon kink crystal is then given by the following solution of Eq. (52),
$$\chi (x)=\frac{\pi }{2}+\mathrm{am}(\xi ,k),\xi =\frac{m_\pi }{k}(xx_0),$$
(57)
($`\mathrm{am}(\xi ,k)`$ is the Jacobian elliptic amplitude function). From this, we can express the baryon density and the various condensates in terms of further Jacobian elliptic functions ($`\mathrm{dn},\mathrm{sn},\mathrm{cn}`$) as follows,
$`\rho _B(x)`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}_x\chi (x)={\displaystyle \frac{m_\pi }{\pi k}}\mathrm{dn}(\xi ,k),`$
$`\overline{\psi }\psi `$ $`=`$ $`\overline{\psi }\psi _\mathrm{v}\mathrm{cos}2\chi (x)=\overline{\psi }\psi _\mathrm{v}\left(\mathrm{cn}^2(\xi ,k)\mathrm{sn}^2(\xi ,k)\right),`$
$`\overline{\psi }\mathrm{i}\gamma _5\psi `$ $`=`$ $`\overline{\psi }\psi _\mathrm{v}\mathrm{sin}2\chi (x)=\overline{\psi }\psi _\mathrm{v}2\mathrm{s}\mathrm{n}(\xi ,k)\mathrm{cn}(\xi ,k).`$ (58)
Here $`\xi `$ is as defined in Eq. (57). Finally, the energy divided by the volume of this kind of matter is given by
$$\frac{}{N}=\frac{m_\pi p_f}{4\pi ^2}\left\{\frac{8}{k}𝐄(k)+4k\left(1\frac{1}{k^2}\right)𝐊(k)\right\},$$
(59)
$`𝐄(k)`$ denoting the complete elliptic integral of the second kind.
Let us now illustrate these results in two regimes of interest, namely at low and high density. At low density ($`p_fm_\pi `$) $`k`$ in Eq. (56) approaches 1 exponentially, and the baryon density features a chain of well resolved lumps whose shape is determined by the single kink solution (Fig. 7). Likewise, the condensates behave like those of a single baryon: $`\overline{\psi }\psi `$ changes from the vacuum value outside the baryons to its negative in their center whereas $`\overline{\psi }\mathrm{i}\gamma _5\psi `$ is peaked in the surface region of each baryon (Figs. 8, 9). These condensates are projections of the distorted “chiral spiral” shown in Fig. 10. The energy (59) for low densities behaves as
$$N\frac{2m_\pi p_f}{\pi ^2}=M_B\rho _B,$$
(60)
showing the expected connection to the baryon mass. At high densities ($`p_fm_\pi `$), $`k`$ approaches 0 like
$$k\frac{m_\pi }{p_f}.$$
(61)
Thus $`\xi `$ in Eq. (57) becomes $`p_f(xx_0)`$. Moreover, for $`k0`$, the Jacobian elliptic functions $`\mathrm{am}(\xi ,k),\mathrm{sn}(\xi ,k),\mathrm{cn}(\xi ,k)`$ are known to reduce to the argument $`\xi `$ and the ordinary trigonometric functions $`\mathrm{sin}\xi `$ and $`\mathrm{cos}\xi `$, respectively. We thus recover the results for the simple chiral spiral in Sect. 4 (the parameter $`x_0`$ has to be readjusted to take care of the shift by $`\pi /2`$ in Eq. (57)). The energy in this case is approximately
$$\frac{}{N}\frac{p_f^2}{2\pi }+\frac{m_\pi ^2}{8\pi }.$$
(62)
The condensates look very much like the sin- and cos-functions of the massless case and need not be plotted. The baryon density wiggles around a constant value, reflecting the strong overlap of the baryons, and can be approximated at high density by
$$\rho _B(x)\frac{p_f}{\pi }\left(1\frac{1}{2}\left(\frac{m_\pi }{p_f}\right)^2\mathrm{sin}^2p_f(xx_0)\right).$$
(63)
The behaviour of the baryon density $`\rho _B(x)`$ as one increases $`p_f`$ (i.e., the mean density) is illustrated in Fig. 11. In the chiral- or high-density limit ($`m_\pi /p_f0`$) $`\rho _B(x)`$ eventually becomes $`x`$-independent. This provides us with another way of understanding the structure of matter described in the previous section, namely as arising from a chain of very extended, strongly overlapping lumps.
Finally, let us come back to the question of validity of the variational calculation based on the chirally modulated vacuum density matrix (37), which we have left open so far. At very low densities when the individual baryons are far apart, we can presumably rely on the numerical results of Ref. since the interaction effects between the baryons become small (as discussed in Refs. and , the baryon-baryon interaction is repulsive and falls off exponentially with the pion Compton wavelength). At high densities, on the other hand, one would expect that a finite quark mass cannot make much difference as long as $`p_fm_\pi `$. This takes us back to the massless case discussed in Sect. 4. In this limit, in turn, it is easy to convince oneself that the calculation becomes exact in the sense that one gets a true solution of the Hartree-Fock equation. Thus for instance for the ’t Hooft model, the massless Hartree-Fock equation reads
$$\omega _n\phi _\alpha ^{(n)}(x)=\mathrm{i}(\gamma _5)_{\alpha \beta }\frac{}{x}\phi _\beta ^{(n)}(x)+\frac{Ng^2}{4}dy|xy|\rho _{\alpha \beta }(x,y)\phi _\beta ^{(n)}(y).$$
(64)
Upon substituting
$$\phi _\alpha ^{(n)}(x)=\left(\mathrm{e}^{\mathrm{i}p_fx\gamma _5}\right)_{\alpha \beta }\stackrel{~}{\phi }_\beta ^{(n)}(x)$$
(65)
as we are instructed to do by the ansatz (37), we discover that $`\stackrel{~}{\phi }^{(n)}`$ does indeed solve the Hartree-Fock equation, the only change being that the single particle energy $`\omega _n`$ gets replaced by $`\omega _n+p_f`$. The same argument goes through in the chiral Gross-Neveu model, or in any field theory where the interaction term has a local chiral invariance. This proves that the result becomes exact in the chiral limit (to leading order in the $`1/N`$ expansion, of course) and makes plausible the hypothesis that it also correctly describes the high density regime for finite quark masses as long as $`p_fm_\pi `$.
## 6) Summary and conclusions
In this paper, we have addressed the problem of baryonic matter in a certain class of exactly soluble field theoretic models, namely chirally invariant, large $`N`$, interacting fermion theories. We started out from a seemingly innocuous and well understood problem, the chiral Gross-Neveu model at finite density and identified one remaining weak spot: The energy density for baryonic matter in the standard Hartree-Fock approach does not have the correct low density limit which can be predicted from the known baryon spectrum of the theory. The origin of this problem which is not cured by a mixed phase approach conceptionally related to the bag model, is evidently the assumption of translational invariance. Whereas this inconsistency can perhaps be ignored in the Gross-Neveu model (as it has been so far, to the best of our knowledge), in QCD<sub>2</sub>, it becomes fatal: Due to confinement, the analogous calculation yields an infinite energy for delocalized baryons or quark matter. This is unavoidable if one is careful in treating the infrared behaviour of the quark single particle energies. These findings have prompted us to think more thoroughly about the structure of baryons in such models and possible implications for the matter problem. We found that it takes only very little effort to generalize a previous Skyrme type treatment of the single, massless baryon to the case of baryonic matter. In the chiral limit, an extremely simple, yet non-trivial, picture emerges: Both the baryon and dense matter are described by a spatially varying chiral angle which is best characterized as a “chiral spiral” with constant helix angle. The number of windings within the full space of length $`L`$ measures the number of baryons in the box. Since, by construction, this kind of state does not cost any potential energy in addition to what is already stored in the vacuum, one does not have to pay the expected high price for delocalizing quarks. The energy density of baryonic matter is identical to that of a free Fermi gas of massless quarks, in spite of the presence of interaction effects. The same picture applies to the chiral Gross-Neveu as well as to the massless ’t Hooft model and should be generic for all chiral large $`N`$ models. The baryon density is constant in space as a consequence of axial current conservation, and we have verified that the whole scenario is exact to leading order in the $`1/N`$ expansion. In a slightly more speculative vein we then investigated modifications due to a small bare quark mass. Here our task was greatly facilitated by the fact that we only needed to pull together two independent investigations, the one of Ref. of the single baryon in field theoretic models with the one of Ref. of the sine-Gordon kink crystal, both inspired in some way by Skyrme’s original ideas. As a result, we have arrived at a rather comprehensive picture of matter at low and high density on the scale of the pion Compton wavelength. The crystal structure now becomes more conspicuous since also the baryon density displays a lattice of individual lumps. As an additional bonus, we have obtained a purely classical, mechanical model of what is going on (the sine-Gordon equation describes a chain of coupled pendulums, the quark mass playing the role of gravity). Given our starting point, namely the problem of baryonic matter in two dimensional quantum field theories like the Nambu–Jona-Lasinio model or QCD, this is rather amusing.
It is noteworthy that a similar chiral structure of fermionic matter has been reported previously in a variety of models different from the present ones. This indicates that the basic results are more generally valid than our derivation might suggest. We mention here in particular the early work on the massive Schwinger model and the more recent work on the massless Schwinger model with inert background charge and QCD<sub>2</sub> with a finite number of colors and flavours . Even more surprising are perhaps quite a number of speculations about spatially inhomogeneous chiral condensates with the same wave number as in our case, but in 3+1 dimensions, in the context of pion condensation , large $`N`$ QCD , or effective chiral models . In some of these works, the analogy with the Overhauser effect and spin-density waves (pairing of particle holes on opposite sides of the Fermi sphere) has been stressed. Although the language used is quite different from ours, there is no doubt that we are dealing with the same physical phenomenon.
As a last remark, we wish to comment on the original Gross-Neveu model with only discrete chiral symmetry (pure $`(\overline{\psi }\psi )^2`$-interaction). Most of the studies of the phase diagram for the Gross-Neveu model have in fact been performed for this model, and one might think that our analysis does not have anything to say about it. However, the criticism of Sect. 2 also applies here. Since the non-chiral Gross-Neveu model has (massive) bound baryons, the low density behaviour of the energy obtained in standard Hartree-Fock approximation cannot be correct, and the phase diagram may also have to be reconsidered.
We thank A.C. Kalloniatis for a critical reading of the manuscript. Partial support of this work by the Bundesministerium für Bildung, Wissenschaft, Forschung und Technologie is gratefully acknowledged.
## Appendix: Hartree-Fock single particle energies in the ’t Hooft model
The Hamiltonian for the massless ’t Hooft model in the axial gauge has the form
$$H=\underset{p,i}{}\frac{2\pi }{L}(p+1/2)\left(a_i^{}(p)a_i(p)b_i^{}(p)b_i(p)\right)+\frac{g^2L}{16\pi ^2}\underset{ij,n0}{}\frac{j_{ij}(n)j_{ji}(n)}{n^2}.$$
(66)
Here we have regularized the theory by enclosing it in a box of length $`L`$ with antiperiodic boundary conditions for the fermions. Note that this form is only valid in the limit $`L\mathrm{}`$ . The $`a_i(p),b_i(p)`$ denote right- and left-handed quark operators, respectively, and the currents $`j_{ij}(n)`$ can be taken in the U($`N`$) form at large $`N`$,
$$j_{ij}(n)=\underset{p}{}\left(a_j^{}(p)a_i(p+n)+b_j^{}(p)b_i(p+n)\right).$$
(67)
It is important to understand that the Coulomb term still contains one- and two-body operators which can be disentangled by normal-ordering (up to $`1/N`$ corrections) as follows,
$`{\displaystyle \underset{ij}{}}j_{ij}(n)j_{ji}(n)`$ $`=`$ $`N{\displaystyle \underset{i,p}{}}\left(a_i^{}(p)a_i(p)+b_i^{}(p)b_i(p)\right)`$
$``$ $`{\displaystyle \underset{ij,pq}{}}(a_j^{}(p)a_j(q)a_i^{}(q+n)a_i(p+n)+a_j^{}(p)b_j(q)b_i^{}(q+n)a_i(p+n)`$
$`+b_j^{}(p)a_j(q)a_i^{}(q+n)b_i(p+n)+b_j^{}(p)b_j(q)b_i^{}(q+n)b_i(p+n))`$
The Hamiltonian can be decomposed correspondingly into one- and two-body operators. Using the basic vacuum expectation values
$`{\displaystyle \underset{i}{}}0|a_i^{}(p)a_i(q)|0`$ $`=`$ $`{\displaystyle \frac{N}{2}}\delta _{pq}(1\mathrm{cos}\theta (p)),`$
$`{\displaystyle \underset{i}{}}0|b_i^{}(p)b_i(q)|0`$ $`=`$ $`{\displaystyle \frac{N}{2}}\delta _{pq}(1+\mathrm{cos}\theta (p)),`$
$`{\displaystyle \underset{i}{}}0|a_i^{}(p)b_i(q)|0`$ $`=`$ $`{\displaystyle \underset{i}{}}0|b_i^{}(p)a_i(q)|0={\displaystyle \frac{N}{2}}\delta _{pq}\mathrm{sin}\theta (p),`$ (69)
the vacuum expectation value of these two contributions are found to be
$`0|H^{(1)}|0`$ $`=`$ $`N{\displaystyle \underset{p}{}}\left({\displaystyle \frac{2\pi }{L}}\left(p+{\displaystyle \frac{1}{2}}\right)\mathrm{cos}\theta (p)+{\displaystyle \frac{Ng^2L}{48}}\right),`$
$`0|H^{(2)}|0`$ $`=`$ $`N{\displaystyle \underset{p}{}}\left({\displaystyle \frac{Ng^2L}{32\pi ^2}}{\displaystyle \underset{n0}{}}{\displaystyle \frac{1}{n^2}}\left[1+\mathrm{cos}(\theta (p)\theta (p+n))\right]\right).`$ (70)
The $`Ng^2L/48`$ term in the 1-body part is just minus twice the “1” term in the 2-body part (remember that $`_{n0}1/n^2=\pi ^2/3`$). Hence, in the sum of both terms, the Coulomb energy involves the combination
$$\frac{1}{n^2}[\mathrm{cos}(\theta (p)\theta (p+n))1]$$
(71)
where the infrared divergence has been tamed since denominator and numerator both vanish at $`n=0`$. This cancellation between quark self-energy and Coulomb potential is similar to what happens in the meson equation of the ’t Hooft model. The continuum limit then yields Eq. (23). It is important to distinguish between one- and two-body operators here, because they enter with different relative weights in the single particle energies and in the total energy. Indeed, in the Hartree-Fock approach, if the single particle energies are decomposed according to their 1- and 2-body contributions as
$$\omega (p)=\omega ^{(1)}(p)+\omega ^{(2)}(p),$$
(72)
then the ground state energy is
$$0|H|0=N\underset{p}{}\left(\omega ^{(1)}(p)+\frac{1}{2}\omega ^{(2)}(p)\right).$$
(73)
The factor 1/2 is necessary to avoid double counting of the 2-body interaction term. By comparison with Eq. (70), we can turn this observation around and simply read off the single particle energies. We find in this way
$`\omega ^{(1)}(p)`$ $`=`$ $`{\displaystyle \frac{2\pi }{L}}\left(p+{\displaystyle \frac{1}{2}}\right)\mathrm{cos}\theta (p)+{\displaystyle \frac{Ng^2L}{48}},`$
$`\omega ^{(2)}(p)`$ $`=`$ $`{\displaystyle \frac{Ng^2L}{16\pi ^2}}{\displaystyle \underset{n0}{}}{\displaystyle \frac{1}{n^2}}\left(1+\mathrm{cos}(\theta (p)\theta (p+n))\right).`$ (74)
Adding up the two contributions to $`\omega (p)`$, the “1” term is now cancelled instead of changing sign. This is the reason why in the continuum limit we get the badly infrared divergent expression (26) for the quark energies. This short-cut derivation gives the same result as the more elaborate approach of Ref. , where a single particle Hartree-Fock Hamiltonian was first identified by commuting $`H`$ with the quark operators and subsequently diagonalized.
### Figure captions
1. Physical fermion mass as a function of the Fermi momentum in the Gross-Neveu model, in units of $`m_0`$.
2. Energy density per color as a function of the Fermi momentum in the Gross-Neveu model. Solid curve: Chirally symmetric solution ($`m=0`$); diamonds: Broken chiral symmetry ($`m`$ according to Fig. 1); dashed straight line: Mixed phase. Units of $`m_0`$.
3. Quark condensate as a function of the Fermi momentum in the ’t Hooft model, in units of $`(Ng^2/2\pi )^{1/2}`$.
4. The complex condensate $`\overline{\psi }\psi +\mathrm{i}\overline{\psi }\mathrm{i}\gamma _5\psi `$ for the single baryon in units of $`\overline{\psi }\psi _\mathrm{v}`$, as a function of $`x`$ in units of $`L`$ (chiral limit).
5. Same as Fig. 4, but for baryonic matter. Each full turn of the spiral increases the baryon number by one unit.
6. Same as Fig. 2 (Gross-Neveu model). Here we have included the energy density of the Skyrme crystal type of state (crosses), the true ground state.
7. Solid curve: Spatial oscillation of the baryon density in the regime $`p_fm_\pi `$; circles: Baryon density for a single baryon.
8. Solid curve: Spatial oscillations of the scalar chiral condensate in the regime $`p_fm_\pi `$; circles: Scalar chiral condensate for a single baryon.
9. Same as Fig. 8, but for the pseudoscalar chiral condensate.
10. Illustration of the distorted “chiral spiral” for baryonic matter at non-zero bare quark mass.
11. Spatial dependence of baryon density as it evolves with increasing average density (or Fermi momentum), in units of $`m_\pi `$.
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# A Spacetime Foam approach to the cosmological constant and entropy
## I Introduction
The problem of merging General Relativity with Quantum Field Theory is known as Quantum Gravity. Many efforts to give meaning to a quantum theory of gravity have been done. Unfortunately until now, such a theory does not exist, even if other approaches like string theory and the still unknown M-theory seem to receive a wide agreement in this direction. Nevertheless a direct investigation of Quantum Gravity shows a lot of interesting aspects. One of these is that at the Planck scale, spacetime could be subjected to topology and metric fluctuations. Such a fluctuating spacetime is known under the name of “spacetime foam” which can be taken as a model for the quantum gravitational vacuum<sup>*</sup><sup>*</sup>*It is interesting to note that there are also indications on how a foamy spacetime can be tested experimentally.. At this scale of lengths (or energies) quantum processes like black hole pair creation could become relevant. To establish if a foamy spacetime could be considered as a candidate for a Quantum Gravitational vacuum, we can examine the structure of the effective potential for such a spacetime. It has been shown that flat space is the classical minimum of the energy for General Relativity. However there are indications that flat space is not the true ground state when a temperature is introduced, at least for the Schwarzschild space in absence of matter fields. It is also argued that when gravity is coupled to $`N`$ conformally invariant scalar fields the evidence that the ground-state expectation value of the metric is flat space is false. Moreover it is also believed that in a foamy spacetime, general relativity can be renormalized when a density of virtual black holes is taken under consideration coupled to $`N`$ fermion fields in a $`1/N`$ expansion. With these examples at hand, we have investigated the possibility of having a ground state different from flat space even at zero temperature and what we have discovered is that there exists an imaginary contribution to the effective potential (more precisely effective energy) at one loop in a Schwarzschild background, that it means that flat space is unstable. What is the physical interpretation associated to this instability. We can begin by observing that the “simplest” quantum process approximating a foamy spacetime, in absence of matter fields, could be a black hole pair creation of the neutral type. One possibility of describing such a process is represented by the Schwarzschild-de Sitter metric (SdS) which asymptotically approaches the de Sitter metric. Its degenerate or extreme version is best known as the Nariai metric. Here we have an external background, the cosmological constant $`\mathrm{\Lambda }`$, which gives a nonzero probability of having a neutral black hole pair produced with its components accelerating away from each other. Nevertheless this process is believed to be highly suppressed, at least for $`\mathrm{\Lambda }1`$ in Planck’s units. In any case, metrics with a cosmological constant have different boundary conditions compared to flat space. The Schwarzschild metric is the only case available. Here the whole spacetime can be regarded as a black hole-anti-black hole pair formed up by a black hole with positive mass $`M`$ in the coordinate system of the observer and an anti black-hole with negative mass $`M`$ in the system where the observer is not present. In this way we have an energy preserving mechanism, because flat space has zero energy and the pair has zero energy too. However, in this case we have not a cosmological force producing the pair: we have only pure gravitational fluctuations. The black hole-anti-black hole pair has also a relevant pictorial interpretation: the black hole with positive mass $`M`$ and the anti black-hole with negative mass $`M`$ can be considered the components of a virtual dipole with zero total energy created by a large quantum gravitational fluctuation. Note that this is the only physical process compatible with the energy conservation. The importance of having the same energy behaviour (asymptotic) is related to the possibility of having a spontaneous transition from one spacetime to another one with the same boundary condition . This transition is a decay from the false vacuum to the true one. However, if we take account of a pair of neutral black holes living in different universes, there is no decay and more important no temperature is involved to change from flat to curved space. To see if this process is realizable we need to compute quantum corrections to the energy stored in the boundaries. These quantum corrections are pure gravitational vacuum excitations which can be measured by the Casimir energy, formally defined as
$$E_{Casimir}\left[\right]=E_0\left[\right]E_0\left[0\right],$$
(1)
where $`E_0`$ is the zero-point energy and $``$ is a boundary. For zero temperature, the idea underlying the Casimir effect is to compare vacuum energies in two physical distinct configurations. For gravitons embedded in flat space, the one-loop contribution to the zero point energy (ZPE) is
$$2\frac{1}{2}\frac{d^3k}{\left(2\pi \right)^3}\sqrt{k^2}.$$
(2)
This term is UV divergent, and its effect is equivalent to inducing a divergent “cosmological constant”
$$\mathrm{\Lambda }_{ZPE}=8\pi G\rho _{ZPE}=4\pi G\frac{d^3k}{\left(2\pi \right)^3}\sqrt{k^2}.$$
(3)
However, it is likely that a natural cutoff of the order of the inverse Planck length comes into play. If this is the case we expect
$$\mathrm{\Lambda }_{ZPE}=c\frac{8\pi }{l_P^2},$$
(4)
where $`c`$ is a dimensionless constant. There are some observational data coming from the Friedmann-Robertson-Walker cosmology constraining the cosmological constant. An estimate of these ones is
$$\begin{array}{c}\mathrm{\Lambda }10^{122}l_P^2\\ c10^{123}\end{array}.$$
(5)
What is the relation between $`\mathrm{\Lambda }_{ZPE}`$ and the neutral black hole pair production? Since the pair produced is mediated by a three-dimensional wormhole with its own ZPE showing a Casimir-like energy and an imaginary part, if we enlarge the number of wormholes (or equivalently the pair creation number) to a certain value $`N`$. In section II, we will prove that the imaginary part disappears and the system from unstable turns to be stable. In Ref. a spacetime foam model made by $`N`$ coherent wormholes has been proposed. In that paper, we have computed the energy density of gravitational fluctuations reproducing the same behavior conjectured by Wheeler during the sixties on dimensional grounds. As an application of the model proposed in Ref., in Ref. we have computed the spectrum of a generic area and as a particular case we have considered the de Sitter geometry. The result is a quantization process whose quanta can be identified with wormholes of Planckian size. In this paper, we would like to continue the investigation of such a model of spacetime foam, by looking at the problem of the cosmological constant, from the Casimir-like energy point of view, and as a consequence of the quantization process of Ref., we will give indications about the black hole mass quantization. The rest of the paper is structured as follows, in section II and III we briefly recall the results reported in Refs., by looking at the mass quantization process. In section IV we approach the cosmological constant problem. We summarize and conclude in section V. Units in which $`\mathrm{}=c=k=1`$ are used throughout the paper.
## II Spacetime foam: the model
In the one-wormhole approximation we have used an eternal black hole, to describe a complete manifold $``$, composed of two wedges $`_+`$ and $`_{}`$ located in the right and left sectors of a Kruskal diagram. The spatial slices $`\mathrm{\Sigma }`$ represent Einstein-Rosen bridges with wormhole topology $`S^2\times R^1`$. Also the hypersurface $`\mathrm{\Sigma }`$ is divided in two parts $`\mathrm{\Sigma }_+`$ and $`\mathrm{\Sigma }_{}`$ by a bifurcation two-surface $`S_0`$. The line element we shall consider is
$$ds^2=N^2\left(r\right)dt^2+\frac{dr^2}{1\frac{2MG}{r}}+r^2\left(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2\right)$$
(6)
or by defining the proper distance from the throat $`dy=\pm dr/\sqrt{1\frac{2MG}{r}}`$, metric $`\left(\text{6}\right)`$ becomes
$$ds^2=N^2dt^2+g_{yy}dy^2+r^2\left(y\right)d\mathrm{\Omega }^2.$$
(7)
$`N`$, $`g_{yy}`$, and $`r`$ are functions of the radial coordinate $`y`$ continuously defined on $``$. The throat of the bridge is at $`r=2MG`$ or $`y=0`$. We choose $`y`$ to be positive in $`\mathrm{\Sigma }_+`$ and negative in $`\mathrm{\Sigma }_{}`$. If we make the identification $`N^2=12MG/r`$, the line element $`\left(\text{6}\right)`$ reduces to the Schwarzschild metric written in another form. The boundaries $`{}_{}{}^{2}S_{+}^{}`$ and $`{}_{}{}^{2}S_{}^{}`$ are located at coordinate values $`y=y_+`$ and $`y=y_{}`$ respectively. The physical Hamiltonian defined on $`\mathrm{\Sigma }`$ is
$$H_P=HH_0=\frac{1}{16\pi l_p^2}_\mathrm{\Sigma }d^3x\left(N+N_i^i\right)+H_{S_+}H_S_{}$$
(8)
where $`l_p^2=G`$ and $`H_{S_+}H_S_{}`$ represents the boundary hamiltonian defined by
$$+\text{ }\frac{1}{8\pi l_p^2}_{S_+}d^2xN\sqrt{\sigma }\left(kk^0\right)\frac{1}{8\pi l_p^2}_S_{}d^2xN\sqrt{\sigma }\left(kk^0\right),$$
(9)
where $`\sigma `$ is the two-dimensional determinant coming from the induced metric $`\sigma _{ab}`$ on the boundaries $`S_\pm `$. The volume term contains two constraints
$$\{\begin{array}{c}=G_{ijkl}\pi ^{ij}\pi ^{kl}\left(\frac{l_p^2}{\sqrt{g}}\right)\left(\frac{\sqrt{g}}{l_p^2}\right)R^{\left(3\right)}=0\hfill \\ ^i=2\pi _{|j}^{ij}=0\hfill \end{array},$$
(10)
both satisfied by the Schwarzschild and Flat metric on-shell, respectively. The supermetric is $`G_{ijkl}=\frac{1}{2}\left(g_{ik}g_{jl}+g_{il}g_{jk}g_{ij}g_{kl}\right)`$ and $`R^{\left(3\right)}`$ denotes the scalar curvature of the surface $`\mathrm{\Sigma }`$. In particular on the boundary we shall use the quasilocal energy definition. Quasilocal energy $`E_{ql}`$ is defined as the value of the Hamiltonian that generates unit time translations orthogonal to the two-dimensional boundaries, i.e.
$`E_{ql}=E_+E_{},`$
$`E_+={\displaystyle \frac{1}{8\pi l_p^2}}{\displaystyle _{S_+}}d^2x\sqrt{\sigma }\left(kk^0\right)`$
$$E_{}=\frac{1}{8\pi l_p^2}_S_{}d^2x\sqrt{\sigma }\left(kk^0\right).$$
(11)
where $`\left|N\right|=1`$ at both $`S_+`$ and $`S_{}`$. $`E_{ql}`$ is the quasilocal energy of a spacelike hypersurface $`\mathrm{\Sigma }=\mathrm{\Sigma }_+\mathrm{\Sigma }_{}`$ bounded by two boundaries $`{}_{}{}^{3}S_{+}^{}`$ and $`{}_{}{}^{3}S_{}^{}`$ located in the two disconnected regions $`M_+`$ and $`M_{}`$ respectively. We have included the subtraction terms $`k^0`$ for the energy. $`k^0`$ represents the trace of the extrinsic curvature corresponding to embedding in the two-dimensional boundaries $`{}_{}{}^{2}S_{+}^{}`$ and $`{}_{}{}^{2}S_{}^{}`$ in three-dimensional Euclidean space. By using the expression of the trace
$$k=\frac{1}{\sqrt{h}}\left(\sqrt{h}n^\mu \right)_{,\mu },$$
(12)
with the normal to the boundaries defined continuously along $`\mathrm{\Sigma }`$ as $`n^\mu =\left(h^{yy}\right)^{\frac{1}{2}}\delta _y^\mu `$, the value of $`k`$ depends on the function $`r,_y`$, where we have assumed that the function $`r,_y`$ is positive for $`S_+`$ and negative for $`S_{}`$. We obtain at either boundary that
$$k=\frac{2r,_y}{r}.$$
(13)
The trace associated with the subtraction term is taken to be $`k^0=2/r`$ for $`B_+`$ and $`k^0=2/r`$ for $`B_{}`$. Thus the quasilocal energy with subtraction terms included is
$$E_{ql}=(E_+E_{})=l_p^2[\left(r[1|r,_y\left|\right]\right)_{y=y_+}\left(r[1|r,_y\left|\right]\right)_{y=y_{}}].$$
(14)
It is easy to see that $`E_+`$ and $`E_{}`$ tend individually to the $`𝒜𝒟`$ mass $`M`$ when the boundaries $`{}_{}{}^{3}B_{+}^{}`$ and $`{}_{}{}^{3}B_{}^{}`$ tend respectively to right and left spatial infinity. It should be noted that the total energy is zero for boundary conditions symmetric with respect to the bifurcation surface, i.e.,
$$E=E_+E_{}=M+\left(M\right)=0,$$
(15)
where the asymptotic contribution has been considered. We can recognize that the expression which defines quasilocal energy is formally of the Casimir type. Indeed, the subtraction procedure present in Eq.$`\left(\text{11}\right)`$ describes an energy difference between two distinct situations with the same boundary conditions. The same behaviour appears in the entropy calculation for the physical system under examination. Indeed
$$S_{tot}=S_+S_{}=\frac{A^+}{4l_p^2}\frac{A^{}}{4l_p^2}\frac{A_H}{4l_p^2}\frac{A_H}{4l_p^2}=0,$$
(16)
where $`A^+`$ and $`A^{}`$ have the same meaning as $`E_+`$ and $`E_{}`$. Note that for both entropy and energy this result is obtained at the tree level. In particular, the quasilocal energy can be interpreted as the tree level approximation of the Casimir energy. We can also see Eqs.$`\left(\text{15}\right)`$ and $`\left(\text{16}\right)`$ from a different point of view. In fact these equations say that flat space can be thought of as a composition of two pieces: the former, with positive energy, in the region $`\mathrm{\Sigma }_+`$ and the latter, with negative energy, in the region $`\mathrm{\Sigma }_{}`$, where the positive and negative concern the bifurcation surface (hole) which is formed due to a topology change of the manifold. The most appropriate mechanism to explain this splitting seems to be a neutral black hole pair creation. To this purpose we begin by considering perturbations at $`\mathrm{\Sigma }`$ of the type
$$g_{ij}^{S,F}=\overline{g}_{ij}^{S,F}+h_{ij},$$
(17)
where $`\overline{g}_{ij}^S`$ is the spatial part of the Schwarzschild background and $`\overline{g}_{ij}^F`$ is the spatial part of the Flat background. In this framework we compute the effective energy defined by
$$E=\underset{\left\{\mathrm{\Psi }\right\}}{\mathrm{min}}\frac{\mathrm{\Psi }\left|H\right|\mathrm{\Psi }}{\mathrm{\Psi }|\mathrm{\Psi }}$$
(18)
with the help of the following rules
$$\begin{array}{c}\left[𝒟h_{ij}\right]h_{ij}\left(\stackrel{}{x}\right)\mathrm{\Psi }\left\{h_{ij}+\overline{g}_{ij}\right\}^2=0,\\ \\ \frac{{\scriptscriptstyle \left[𝒟h_{ij}\right]d^3xh_{ij}\left(\stackrel{}{x}\right)h_{kl}\left(\stackrel{}{y}\right)\mathrm{\Psi }\left\{h_{ij}+\overline{g}_{ij}\right\}^2}}{{\scriptscriptstyle \left[𝒟h_{ij}\right]\mathrm{\Psi }\left\{h_{ij}+\overline{g}_{ij}\right\}^2}}=K_{ijkl}(\stackrel{}{x},\stackrel{}{y}).\end{array}$$
(19)
In particular
$$E\left(M\right)=\frac{\mathrm{\Psi }\left|H^{Schw.}\right|\mathrm{\Psi }}{\mathrm{\Psi }|\mathrm{\Psi }}+\frac{\mathrm{\Psi }\left|H_{boundary}^{Schw.}\right|\mathrm{\Psi }}{\mathrm{\Psi }|\mathrm{\Psi }}$$
(20)
and
$$E\left(0\right)=\frac{\mathrm{\Psi }\left|H^{Flat}\right|\mathrm{\Psi }}{\mathrm{\Psi }|\mathrm{\Psi }}+\frac{\mathrm{\Psi }\left|H_{boundary}^{Flat}\right|\mathrm{\Psi }}{\mathrm{\Psi }|\mathrm{\Psi }}$$
(21)
so that the physical Hamiltonian is given by the difference (Casimir energy) $`\mathrm{\Delta }E\left(M\right)`$
$$=E\left(M\right)E\left(0\right)=\frac{\mathrm{\Psi }\left|H^{Schw.}H^{Flat}\right|\mathrm{\Psi }}{\mathrm{\Psi }|\mathrm{\Psi }}+\frac{\mathrm{\Psi }\left|H_{ql}\right|\mathrm{\Psi }}{\mathrm{\Psi }|\mathrm{\Psi }}.$$
(22)
The quantity $`\mathrm{\Delta }E\left(M\right)`$ is computed by means of a variational approach, where the WKB functionals are substituted with trial wave functionals. By restricting our attention to the graviton sector of the Hamiltonian approximated to second order, hereafter referred as $`H_{|2}`$, we define
$`E_{|2}={\displaystyle \frac{\mathrm{\Psi }^{}\left|H_{|2}\right|\mathrm{\Psi }^{}}{\mathrm{\Psi }^{}|\mathrm{\Psi }^{}}},`$
where
$`\mathrm{\Psi }^{}=\mathrm{\Psi }\left[h_{ij}^{}\right]=𝒩\mathrm{exp}\left\{{\displaystyle \frac{1}{4l_p^2}}\left[\left(g\overline{g}\right)K^1\left(g\overline{g}\right)_{x,y}^{}\right]\right\}.`$
After having functionally integrated $`H_{|2}`$, we get
$$H_{|2}=\frac{1}{4l_p^2}_\mathrm{\Sigma }d^3x\sqrt{g}G^{ijkl}\left[K^1(x,x)_{ijkl}+\left(\mathrm{}_2\right)_j^aK^{}(x,x)_{iakl}\right]$$
(23)
The propagator $`K^{}(x,x)_{iakl}`$ comes from a functional integration and it can be represented as
$$K^{}(\stackrel{}{x},\stackrel{}{y})_{iakl}:=\underset{N}{}\frac{h_{ia}^{}\left(\stackrel{}{x}\right)h_{kl}^{}\left(\stackrel{}{y}\right)}{2\lambda _N\left(p\right)},$$
(24)
where $`h_{ia}^{}\left(\stackrel{}{x}\right)`$ are the eigenfunctions of
$$\left(\mathrm{}_2\right)_j^a:=\mathrm{}\delta _j^a+2R_j^a.$$
(25)
This is the Lichnerowicz operator projected on $`\mathrm{\Sigma }`$ acting on traceless transverse quantum fluctuations and $`\lambda _N\left(p\right)`$ are infinite variational parameters. $`\mathrm{}`$ is the curved Laplacian (Laplace-Beltrami operator) on a Schwarzschild background and $`R_{j\text{ }}^a`$ is the mixed Ricci tensor whose components are:
$$R_j^a=diag\{\frac{2MG}{r^3},\frac{MG}{r^3},\frac{MG}{r^3}\}.$$
(26)
The minimization with respect to $`\lambda `$ and the introduction of a high energy cutoff $`\mathrm{\Lambda }`$ give to the Eq. $`\left(\text{22}\right)`$ the following form
$$\mathrm{\Delta }E\left(M\right)\frac{V}{32\pi ^2}\left(\frac{3MG}{r_0^3}\right)^2\mathrm{ln}\left(\frac{r_0^3\mathrm{\Lambda }^2}{3MG}\right),$$
(27)
where $`V`$ is the volume of the system and $`r_0`$ is related to the minimum radius compatible with the wormhole throat. We know that classically, the minimum radius is achieved when $`r_0=2MG`$. However, it is likely that quantum processes come into play at short distances, where the wormhole throat is defined, introducing a quantum radius $`r_0>2MG`$. Nevertheless, since we are interested to probe the energy contribution near the Planck scale we can fix the value of $`r_0=l_p`$ from now on. We now compute the minimum of $`\mathrm{\Delta }E\left(M\right)`$, after having rescaled the variable $`M`$ to a scale variable $`x=3MG/\left(r_0^3\mathrm{\Lambda }^2\right)=3M/\left(l_p\mathrm{\Lambda }^2\right)`$. Thus
$`\mathrm{\Delta }E\left(M\right)\mathrm{\Delta }E(x,\mathrm{\Lambda })={\displaystyle \frac{V}{32\pi ^2}}\mathrm{\Lambda }^4x^2\mathrm{ln}x`$
We obtain two values for $`x`$: $`x_1=0`$, i.e. flat space and $`x_2=e^{\frac{1}{2}}`$. At the minimum
$$\mathrm{\Delta }E\left(x_2\right)=\frac{V}{64\pi ^2}\frac{\mathrm{\Lambda }^4}{e}.$$
(28)
Nevertheless, there exists another part of the spectrum which has to be considered: the discrete spectrum containing one mode. This gives the energy an imaginary contribution, namely we have discovered an unstable mode. Let us briefly recall, how this appears. The eigenvalue equation
$$\left(\mathrm{}_2\right)_i^ah_{aj}=\alpha h_{ij}$$
(29)
can be studied with the Regge-Wheeler method. The perturbations can be divided in odd and even components. The appearance of the unstable mode is governed by the gravitational field component $`h_{11}^{even}`$. Explicitly
$`E^2H\left(r\right)`$
$$=\left(1\frac{2MG}{r}\right)\frac{d^2H\left(r\right)}{dr^2}+\left(\frac{2r3MG}{r^2}\right)\frac{dH\left(r\right)}{dr}\frac{4MG}{r^3}H\left(r\right),$$
(30)
where
$$h_{11}^{even}(r,\vartheta ,\varphi )=\left[H\left(r\right)\left(1\frac{2m}{r}\right)^1\right]Y_{00}(\vartheta ,\varphi )$$
(31)
and $`E^2>0`$. Eq.$`\left(\text{30}\right)`$ can be transformed into
$$\mu =\frac{\underset{0}{\overset{\overline{y}}{}}𝑑y\left[\left(\frac{dh\left(y\right)}{dy}\right)^2\frac{3}{2\rho \left(y\right)^3}h^2\left(y\right)\right]}{\underset{0}{\overset{\overline{y}}{}}𝑑yh^2\left(y\right)},$$
(32)
where $`\mu `$ is the eigenvalue, $`y`$ is the proper distance from the throat in dimensionless form. If we choose $`h(\lambda ,y)=\mathrm{exp}\left(\lambda y\right)`$ as a trial function we numerically obtain $`\mu =.701626`$. In terms of the energy square we have
$$E^2=\mathrm{.\hspace{0.17em}175\hspace{0.17em}41}/\left(MG\right)^2$$
(33)
to be compared with the value $`E^2=\mathrm{.\hspace{0.17em}19}/\left(MG\right)^2`$ of Ref.. Nevertheless, when we compute the eigenvalue as a function of the distance $`y`$, we discover that in the limit $`\overline{y}0`$,
$$\mu \mu \left(\lambda \right)=\lambda ^2\frac{3}{2}+\frac{9}{8}\left[\overline{y}^2+\frac{\overline{y}}{2\lambda }\right].$$
(34)
Its minimum is at $`\stackrel{~}{\lambda }=\left(\frac{9}{32}\overline{y}\right)^{\frac{1}{3}}`$ and
$$\mu \left(\stackrel{~}{\lambda }\right)=\mathrm{1.\hspace{0.17em}287\hspace{0.17em}8}\overline{y}^{\frac{2}{3}}+\frac{9}{8}\overline{y}^2\frac{3}{2}.$$
(35)
It is evident that there exists a critical radius where $`\mu `$ turns from negative to positive. This critical value is located at $`\rho _c=\mathrm{1.\hspace{0.17em}113\hspace{0.17em}4}`$ to be compared with the value $`\rho _c=1.445`$ obtained by B. Allen in . What is the relation with the large number of wormholes? As mentioned in Ref., when the number of wormholes grows, the space available for every single wormhole has to be reduced to avoid overlapping of the wave functions. Consider the simple case of two wormholes covering the hypersurface $`\mathrm{\Sigma }`$, namely $`\mathrm{\Sigma }=\mathrm{\Sigma }_1\mathrm{\Sigma }_2`$, $`\mathrm{\Sigma }_1\mathrm{\Sigma }_2=\mathrm{}`$. The second property assures that the boundaries and the support of the wave functional do not overlap, in agreement with the WKB approximation. $`\mathrm{\Sigma }_1`$and $`\mathrm{\Sigma }_2`$ have topology $`S^2\times R^1`$ with boundaries $`\mathrm{\Sigma }_1^\pm `$ and $`\mathrm{\Sigma }_2^\pm `$ with respect to each bifurcation surface. The total Hamiltonian, in this case, is $`H_T=H_1+H_2`$, i.e. (here we are looking at boundary terms, because in this discussion they are the only relevant ones)
$`H_T=2H={\displaystyle \frac{1}{8\pi l_p^2}}\left[2{\displaystyle _{S_+}}d^2x\sqrt{\sigma }\left(kk^0\right)2{\displaystyle _S_{}}d^2x\sqrt{\sigma }\left(kk^0\right)\right]`$
$`={\displaystyle \frac{1}{8\pi l_p^2}}[2\left(r[1|r,_y\left|\right]\right)_{y=y_+}2\left(r[1|r,_y\left|\right]\right)_{y=y_{}}]`$
$`={\displaystyle \frac{1}{8\pi l_p^2}}\left[2r_+\left(1\sqrt{1{\displaystyle \frac{2Ml_p^2}{r_+}}}\right)2r_{}\left(1\sqrt{1{\displaystyle \frac{2Ml_p^2}{r_{}}}}\right)\right]`$
$`={\displaystyle \frac{1}{8\pi l_p^2}}\left[2r_+\left(1\sqrt{1{\displaystyle \frac{2Ml_{2_w}^2}{2r_+}}}\right)2r_{}\left(1\sqrt{1{\displaystyle \frac{2Ml_{2_w}^2}{2r_{}}}}\right)\right]`$
$$=\frac{1}{8\pi l_p^2}\left[R_+\left(1\sqrt{1\frac{2Ml_{2_w}^2}{R_+}}\right)R_{}\left(1\sqrt{1\frac{2Ml_{2_w}^2}{R_{}}}\right)\right].$$
(36)
This means that the total quasilocal energy is the same of a single wormhole with boundaries satisfying the relation $`R_\pm =2r_\pm `$, or in other words the value of $`R_\pm `$ in presence of two wormholes is divided by two. Note that $`R_\pm `$ are the boundary values corresponding to the single wormhole case. This implies that if we put more and more wormholes, say $`N_w`$, the initial boundary located at $`R_\pm `$ will be reduced and $`GN_wG`$. This boundary reduction is important, because it is related to the disappearing of the unstable mode. Let us see how. If we fix the initial boundary at $`R_\pm `$, then in presence of $`N_w`$ wormholes, it will be reduced to $`R_\pm /N_w`$. This means that boundary conditions are not fixed at infinity, but at a certain finite radius and the $`ADM`$ mass term is substituted by the quasilocal energy expression under the condition of having symmetry with respect to each bifurcation surface. The effect on the unstable mode is clear: as $`N_w`$ grows, the boundary radius reduces more and more until it will reach the critical value $`\rho _c`$ below which no negative mode will appear corresponding to a critical wormholes number $`N_{w_c}`$. To this purpose, suppose to consider $`N_w`$ wormholes and assume that there exists a covering of $`\mathrm{\Sigma }`$ such that $`\mathrm{\Sigma }=\underset{i=1}{\overset{N_w}{}}\mathrm{\Sigma }_i`$, with $`\mathrm{\Sigma }_i\mathrm{\Sigma }_j=\mathrm{}`$ when $`ij`$. Each $`\mathrm{\Sigma }_i`$ has the topology $`S^2\times R^1`$ with boundaries $`\mathrm{\Sigma }_i^\pm `$ with respect to each bifurcation surface. On each surface $`\mathrm{\Sigma }_i`$, quasilocal energy gives
$$E_{i\text{ }\mathrm{ql}}=\frac{1}{8\pi l_p^2}_{S_{i+}}d^2x\sqrt{\sigma }\left(kk^0\right)\frac{1}{8\pi l_p^2}_{S_i}d^2x\sqrt{\sigma }\left(kk^0\right).$$
(37)
Thus if we apply the same procedure of the single case on each wormhole, we obtain
$$E_{i\text{ }\mathrm{ql}}=(E_{i+}E_i)=l_p^2\left(r[1|r,_y\left|\right]\right)_{y=y_{i+}}l_p^2\left(r[1|r,_y\left|\right]\right)_{y=y_i}.$$
(38)
Note that the total quasilocal energy is zero for boundary conditions symmetric with respect to each bifurcation surface $`S_{0,i}`$. If we assume this kind of symmetry for boundary conditions, it is immediate to recognize that the vanishing of the boundary term is guaranteed beyond the semiclassical approximation, because every term on one wedge of the hypersurface $`\mathrm{\Sigma }_i`$ will be compensated by the term on the other wedge of the same hypersurface $`\mathrm{\Sigma }_i`$, giving therefore zero energy contribution. We are interested in a large number of wormholes, each of them contributing with a term of the type $`E_{i\text{ }\mathrm{ql}}`$. If the wormholes number is $`N_w`$, we obtain (semiclassically, i.e., without self-interactions)
$$H_{tot}^{N_w}=H^1+H^2+\mathrm{}+H^{N_w}.$$
(39)
Thus the total energy for the collection is
$`E_{|2}^{tot}=N_wH_{|2}.`$
The same happens for the trial wave functional which is the product of $`N_w`$ t.w.f.. Thus
$$\mathrm{\Psi }_{tot}^{}=\mathrm{\Psi }_1^{}\mathrm{\Psi }_2^{}\mathrm{}\mathrm{}\mathrm{\Psi }_{N_w}^{}=𝒩\mathrm{exp}N_w\left\{\frac{1}{4l_p^2}\left[\left(g\overline{g}\right)K^1\left(g\overline{g}\right)_{x,y}^{}\right]\right\}.$$
(40)
Thus for the $`N_w`$ wormholes, one gets
$$\mathrm{\Delta }E_{N_w}(x,\mathrm{\Lambda })N_w\frac{V}{32\pi ^2}\mathrm{\Lambda }^4x^2\mathrm{ln}x,$$
(41)
where we have defined the usual scale variable $`x=3M/\left(l_p\mathrm{\Lambda }^2\right)`$. Then at one loop the cooperative effects of wormholes behave as one macroscopic single field multiplied by $`N_w^2`$, but without the unstable mode. At the minimum, $`\overline{x}=e^{\frac{1}{2}}`$
$$\mathrm{\Delta }E\left(\overline{x}\right)=N_w\frac{V}{64\pi ^2}\frac{\mathrm{\Lambda }^4}{e}.$$
(42)
This means that we have obtained a minimum of the effective energy away by the flat space, indicating that another configuration has to be considered for the ground state of quantum gravity. Let us examine the implications on the area quantization, entropy and the cosmological constant.
## III Area spectrum and Entropy
Bekenstein made the proposal that a black hole does have an entropy proportional to the area of its horizon
$$S_{bh}=const\times A_{hor}.$$
(43)
In particular, in natural units one finds that the proportionality constant is set to $`1/4G=1/4l_p^2`$, so that the entropy becomes
$$S=\frac{A}{4G}=\frac{A}{4l_p^2}.$$
(44)
Following Bekenstein’s proposal on the quantization of the area for nonextremal black holes we have
$$A_n=\alpha l_p^2\left(n+\eta \right)\text{ }\eta >1\text{ }n=1,2,\mathrm{}$$
(45)
Many attempts to recover the area spectrum have been done, see Refs. for a review. Note that the appearance of a discrete spectrum is not so trivial. Indeed there are other theories, based on spherically symmetric metrics in a mini-superspace approach, whose mass spectrum is continuous. The area is measured by the quantity
$$A\left(S_0\right)=_{S_0}d^2x\sqrt{\sigma }.$$
(46)
$`\sigma `$ is the two-dimensional determinant coming from the induced metric $`\sigma _{ab}`$ on the boundary $`S_0`$. We would like to evaluate the mean value of the area
$$A\left(S_0\right)=\frac{\mathrm{\Psi }_F\left|\widehat{A}\right|\mathrm{\Psi }_F}{\mathrm{\Psi }_F|\mathrm{\Psi }_F}=\frac{\mathrm{\Psi }_F\left|\widehat{_{S_0}d^2x\sqrt{\sigma }}\right|\mathrm{\Psi }_F}{\mathrm{\Psi }_F|\mathrm{\Psi }_F},$$
(47)
computed on
$$|\mathrm{\Psi }_F=\mathrm{\Psi }_1^{}\mathrm{\Psi }_2^{}\mathrm{}\mathrm{}\mathrm{\Psi }_{N_w}^{}.$$
(48)
Since we are working with spherical symmetric wormholes we consider $`\sigma _{ab}=\overline{\sigma }_{ab}+\delta \sigma _{ab}`$, where $`\overline{\sigma }_{ab}`$ is such that $`_{S_0}d^2x\sqrt{\overline{\sigma }}=4\pi \overline{r}^2`$ and $`\overline{r}`$ is the radius of $`S_0`$. To the lowest level in the expansion of $`\sigma _{ab}`$ we obtain that
$$A\left(S_0\right)=\frac{\mathrm{\Psi }_F\left|\widehat{A}\right|\mathrm{\Psi }_F}{\mathrm{\Psi }_F|\mathrm{\Psi }_F}=4\pi \overline{r}^2.$$
(49)
Suppose to consider the mean value of the area $`A`$ computed on a given macroscopic fixed radius $`R`$. On the basis of our foam model, we obtain $`A=\underset{i=1}{\overset{N}{}}A_i`$, with $`A_iA_j=\mathrm{}`$ when $`ij`$. Thus
$$A=4\pi R^2=\underset{i=1}{\overset{N}{}}A_i=\underset{i=1}{\overset{N}{}}4\pi \overline{r}_i^2.$$
(50)
In Refs. we have considered, as a first approximation, the limit $`\overline{r}_il_p`$ and we have obtained
$$A=NA_{l_p}=N4\pi l_p^2.$$
(51)
Nevertheless an improvement of Eq.$`\left(\text{51}\right)`$ is possible if we introduce a scale variable $`x_i=\overline{r}_i/l_p`$ which leads to
$$A=4\pi l_p^2\underset{i=1}{\overset{N}{}}x_i^2=4\pi l_p^2N\overline{x^2}=4\pi l_p^2N\alpha .$$
(52)
Thus the number $`\alpha `$ appearing in Eq.$`\left(\text{45}\right)`$, here comes from an averaging process. Note that the $`4\pi `$ factor is a consequence of the $`S^2`$ wormhole topology which is an intrinsic feature of our foam model. Comparison of Eq.$`\left(\text{52}\right)`$with the Bekenstein area spectrum gives
$$4\pi l_p^2N\alpha =4l_p^2N\mathrm{ln}2.$$
(53)
This fixes the coefficient $`\alpha `$ to
$$\frac{\mathrm{ln}2}{\pi }=\alpha $$
(54)
and the entropy is
$$S=\frac{A}{4l_p^2}=\frac{4l_p^2N\mathrm{ln}2}{4l_p^2}=N\mathrm{ln}2.$$
(55)
$`N`$ is such that $`NN_{w_c}`$ and $`N_{w_c}`$ is the critical wormholes number above which we have the stability of our foam model. On the other hand if we apply the same reasoning of Refs., applied to the quantity
$$\frac{A}{4\pi l_p^2}=\underset{i=1}{\overset{N}{}}x_i^2,$$
(56)
produces an extra-factor of the form $`\mathrm{ln}2/\pi `$, when compared with the Hawking’s coefficient $`1/4`$. This factor can be absorbed by choosing a suitable normalization constant when we apply the partition of the integer $`N`$. In any case we are led to the Bekenstein-Hawking relation between entropy and area
$$S=\frac{A}{4l_p^2}.$$
(57)
We can use Eq.$`\left(\text{52}\right)`$ to compute the entropy for some specific geometries, for example, the Schwarzschild geometry
$$S=\frac{4\pi \left(2MG\right)^2}{4G}=4\pi M^2G=4\pi M^2l_p^2=N\mathrm{ln}2.$$
(58)
Thus the Schwarzschild black hole mass is quantized in terms of $`l_p`$ giving therefore the relation
$$M=\frac{\sqrt{N}}{2l_p}\sqrt{\frac{\mathrm{ln}2}{\pi }},$$
(59)
which is in agreement with the results presented in Refs.. This implies also that the level spacing of the transition frequencies is
$$\omega _0=\mathrm{\Delta }M=\left(8\pi Ml_p^2\right)^1\mathrm{ln}2$$
(60)
and the Schwarzschild radius is quantized in terms of $`l_p.`$ Indeed
$$R_S=2MG=2Ml_p^2=\sqrt{N}l_p\sqrt{\frac{\mathrm{ln}2}{\pi }}.$$
(61)
## IV The cosmological constant
Einstein introduced his cosmological constant $`\mathrm{\Lambda }_c`$ in an attempt to generalize his original field equations. The modified field equations are
$$R_{\mu \nu }\frac{1}{2}g_{\mu \nu }R+\mathrm{\Lambda }_cg_{\mu \nu }=8\pi GT_{\mu \nu }.$$
(62)
By redefining
$$T_{tot}^{\mu \nu }T^{\mu \nu }\frac{\mathrm{\Lambda }_c}{8\pi G}g^{\mu \nu },$$
(63)
one can regain the original form of the field equations
$$R_{\mu \nu }\frac{1}{2}g_{\mu \nu }R=8\pi GT_{\mu \nu },$$
(64)
at the prize of introducing a vacuum energy density and vacuum stress-energy tensor
$$\rho _\mathrm{\Lambda }=\frac{\mathrm{\Lambda }_c}{8\pi G};T_\mathrm{\Lambda }^{\mu \nu }=\rho _\mathrm{\Lambda }g^{\mu \nu }.$$
(65)
If we look at the Hamiltonian in presence of a cosmological term, we have the expression
$$H=_\mathrm{\Sigma }d^3x(N\left(+\rho _\mathrm{\Lambda }\sqrt{g}\right)+N^i_i),$$
(66)
where $``$ is the usual Hamiltonian density defined without a cosmological term. We know that the effect of vacuum fluctuation is to inducing a cosmological term. Indeed by looking at Eq.$`\left(\text{42}\right)`$, we have that
$$\frac{\mathrm{\Delta }H}{V}=N_w\frac{\mathrm{\Lambda }^4}{64e\pi ^2}.$$
(67)
On the other hand, the WDW equation in presence of a cosmological constant is
$$\left[\frac{16\pi l_p^2}{\sqrt{g}}G_{ijkl}\pi ^{ij}\pi ^{kl}\frac{\sqrt{g}}{16\pi l_p^2}\left(R2\mathrm{\Lambda }_c\right)\right]\mathrm{\Psi }\left[g_{ij}\right]=0.$$
(68)
By integrating over the hypersurface $`\mathrm{\Sigma }`$ and looking at the expectation values computed on the state $`|\mathrm{\Psi }_F`$ of Eq.$`\left(\text{48}\right)`$ the WDW equation becomes
$`\mathrm{\Psi }_F\left|{\displaystyle _\mathrm{\Sigma }}d^3x\left[{\displaystyle \frac{16\pi l_p^2}{\sqrt{g}}}G_{ijkl}\pi ^{ij}\pi ^{kl}{\displaystyle \frac{\sqrt{g}}{16\pi l_p^2}}R\right]\right|\mathrm{\Psi }_F`$
$$=\mathrm{\Psi }_F\left|\frac{\mathrm{\Lambda }_c}{8\pi l_p^2}_\mathrm{\Sigma }d^3x\sqrt{g}\right|\mathrm{\Psi }_F=\frac{\mathrm{\Lambda }_c}{8\pi l_p^2}_\mathrm{\Sigma }d^3x\sqrt{g}=\frac{\mathrm{\Lambda }_c}{8\pi l_p^2}V_c.$$
(69)
$`V_c`$ is the cosmological volume. The first term of Eq.$`\left(\text{69}\right)`$ is formally the same that generates the vacuum fluctuation $`\left(\text{67}\right)`$. Thus, by comparing the second term of Eq.$`\left(\text{69}\right)`$ with Eq.$`\left(\text{67}\right)`$, we have
$$\frac{\mathrm{\Lambda }_c}{8\pi l_p^2}V_c=N_w\frac{\mathrm{\Lambda }^4}{64e\pi ^2}V_w.$$
(70)
Therefore
$$\mathrm{\Lambda }_c=N_w^2\frac{\mathrm{\Lambda }^4l_p^2}{V_c8e\pi }V_w.$$
(71)
The cosmological volume has to be rescaled in terms of the wormhole radius, in such a way to obtain that $`V_cN_w^3V_w`$ and we have rescaled the Planck length as in page II. This is the direct consequence of the boundary rescaling, namely $`R_\pm `$ $``$ $`R_\pm /N_w`$. Thus
$$\mathrm{\Lambda }_c=\frac{\mathrm{\Lambda }^4l_p^2}{N_w8e\pi }.$$
(72)
This is the value of the induced cosmological constant. On the other hand, if we apply the area quantization procedure of Eq.$`\left(\text{58}\right)`$ to the de Sitter geometry, one gets
$$S=\frac{3\pi }{l_p^2\mathrm{\Lambda }_c}=N\mathrm{ln}2,$$
(73)
that isA relation relating $`\mathrm{\Lambda }`$ and $`G`$, via an integer $`N`$ appeared also in Ref.. Nevertheless in Ref., $`N`$ represents the number of scalar fields and the bound from above and below
$`\left|2GN\mathrm{\Lambda }/32\right|\sqrt{3}`$
comes into play, instead of the equality (74).
$$\frac{3\pi }{\mathrm{ln}2l_p^2N}=\mathrm{\Lambda }_c.$$
(74)
Thus the cosmological constant $`\mathrm{\Lambda }`$ is “quantized” in terms of $`l_p`$. Note that when the wormholes number $`N`$ is quite “large”, $`\mathrm{\Lambda }0.`$ We could try to see what is the rate of change between an early universe value of the cosmological constant and the value that we observe. In inflationary models of the early universe is assumed to have undergone an early phase with a large effective $`\mathrm{\Lambda }\left(10^{10}10^{11}GeV\right)^2`$ for GUT era inflation, or $`\mathrm{\Lambda }\left(10^{16}10^{18}GeV\right)^2`$ for Planck era inflation. A subsequent phase transition would then produce a region of space-time with $`\mathrm{\Lambda }\left(10^{42}GeV\right)^2`$, i.e. the space in which we now live. For GUT era inflation, we have (we are looking only at the order of magnitude)
$$10^{20}10^{22}GeV^2=\frac{1}{N}10^{38}GeV^2N=10^{16}10^{18},$$
(75)
while for Planck era inflation we have
$$10^{32}10^{36}GeV^2=\frac{1}{N}10^{38}GeV^2N=10^610^2,$$
(76)
to be compared with the value of $`\left(10^{42}GeV\right)^2`$ which gives a wormholes number of the order of
$$10^{84}GeV^2=\frac{1}{N}10^{38}GeV^2N=10^{122}.$$
(77)
In our model this very huge number represents the maximum wormholes number of Planck size that can be stored into an area of radius equal to the cosmological radius. This is in agreement with observational data of Eq.$`\left(\text{5}\right)`$. If we compare the previous value of $`\mathrm{\Lambda }_c`$ with the value of Eq.$`\left(\text{74}\right)`$, one gets
$$\mathrm{\Lambda }_c=\frac{\mathrm{\Lambda }^4l_p^2}{N_w8e\pi }=\frac{3\pi }{\mathrm{ln}2l_p^2N_w},$$
(78)
namely we have a constraint on the U.V. cut-off
$$\mathrm{\Lambda }^4=\frac{24e\pi ^2}{\mathrm{ln}2l_p^4}.$$
(79)
The probability to realize a foamy spacetime is measured by
$$\mathrm{\Gamma }_{\mathrm{N}\mathrm{holes}}=\frac{P_{\mathrm{N}\mathrm{holes}}}{P_{\mathrm{flat}}}\frac{P_{\mathrm{foam}}}{P_{\mathrm{flat}}}.$$
(80)
In a Euclidean time this is
$$P\left|e^{\left(\mathrm{\Delta }E\right)\left(\mathrm{\Delta }t\right)}\right|^2\left|\mathrm{exp}\left(N_w\frac{\mathrm{\Lambda }^4}{e64\pi ^2}\right)\left(V\mathrm{\Delta }t\right)\right|^2.$$
(81)
From Eq.(70), we obtain
$$P\left|\mathrm{exp}\left(\frac{\mathrm{\Lambda }_c}{8\pi l_p^2}V_c\right)\left(\mathrm{\Delta }t\right)\right|^2.$$
(82)
To be concrete we can consider again the de Sitter case. Thus
$$\mathrm{\Delta }t=2\pi \sqrt{\frac{3}{\mathrm{\Lambda }_c}}$$
(83)
and the cosmological volume is given by
$$V_c=\frac{4\pi }{3}\left(\sqrt{\frac{3}{\mathrm{\Lambda }_c}}\right)^3,$$
(84)
namely
$$\mathrm{exp}\left(3\pi /l_p^2\mathrm{\Lambda }_c\right).$$
(85)
Thus we recover the Hawking result about the cosmological constant approaching zero. Note that the vanishing of $`\mathrm{\Lambda }_c`$ is related to the growing of the wormholes number.
## V Conclusions
In this paper we have continued the investigation of our spacetime foam model presented in Refs., where we have obtained a “quantization” process in the sense that we can fill spacetime with a given integer number of disjoint non-interacting wormholes. At first look, it seems that our foam model looks promising, since in this framework we have reproduced certain features that a quantum theory of gravity must possess. Nevertheless a lot of points must be clarified. First of all the rôle of the Planckian cutoff that here is computed by comparing a tree level quantity (the entropy of the de Sitter space) with a one-loop quantity (the induced cosmological constant or the Casimir energy). Secondly, the effect of quantum fluctuation has to be inserted in the entropy computation. This could cause a modification of Eqs.(59) and (74) and therefore of estimate (77). On the other hand, as a first consequence we have obtained that the area operator has a discrete spectrum, whose quanta are Planck size wormholes. This is in agreement with the quantized area proposed heuristically by Bekenstein and also with the loop quantum gravity predictions of Refs.. Note that in order to have stability, it is the energy configuration that forces spacetime to be filled with $`N`$ wormholes of the Planckian size. Since the area is related to the entropy via the Bekenstein-Hawking relation, as a direct application, a “mass quantization” of a Schwarzschild black hole whose mass is $`M`$ is obtained, in agreement with Refs.. The second consequence of our model is the generation of a positive cosmological constant induced by vacuum fluctuations. Due to the uncertainty relation
$$\mathrm{\Delta }E\frac{A}{L^4}N_w\frac{V}{64\pi ^2}\frac{\mathrm{\Lambda }^4}{e}A\mathrm{\Lambda }^4.$$
(86)
The negative fourth power of the cutoff (or the inverse of the fourth power of the region of dimension L) is a clear signal of a Casimir-like energy generated by vacuum fluctuations. As a consequence a positive cosmological constant is induced by such fluctuations.
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# 1 Introduction
## 1 Introduction
The extended (soliton-like) objects quantisation problem considered in this paper have more than twenty years old history, see the review papers . But absence of progress in this field evokes anxiety, noting a number of unsolved by this reason important physical problems.
One of possible solutions of this problem was offered in . The aim of this article is to show that the approach described in this paper leads to the strong coupling perturbation theory over inverse interaction constant $`1/g`$.
Our approach is based on the idea that the measure $`DM`$ of the functional integral representation for
$$\rho (E)=𝑑u_1𝑑u_2|<u_2;E|u_1;E>|^2$$
(1.1)
is $`\delta `$-like (Diracian):
$$DM(u)=\underset{x,t}{}du(x,t)\delta \left(\frac{\delta S(u)}{\delta u(x,t)}j(x,t)\right),$$
(1.2)
where $`j(x,t)`$ is the random force of quantum excitations. We will consider the symplest one particle quantum problem and the states in (1.1) are described by the boundary values of coordinate $`u_i`$ and energy $`E`$. In Sec.2.1 the physical basis of (1.2) will be discussed and full derivation will be given in Sec.2.2.
The $`\delta `$-function of (1.2) means that the strict space-time local equality:
$$\frac{\delta S(u)}{\delta u(x,t)}=j(x,t)$$
(1.3)
defines the complete set of necessary contributions. The general properties of theory defined on the $`\delta `$-like measure is listed in Sec.2.3. Note, (1.3) is not the consequence of Hamiltonian variational principle and (1.2) will be derived in Sec.2.2, proceeding from the conservation of total probability (unitarity condition) .
Eq.(1.3) shows that a transformation of kinetic part of Lagrangian without fail induce the tangent transformation of quantum source $`j(x,t)`$ . We will use this possibility to describe the quantum dynamics in useful terms. Namely, we will apply the canonical transformation of (1.2) to the collective coordinates. They will have a meaning of (action, angle)-type variables and will form the cotangent foliation $`T^{}\mathrm{\Omega }`$ to the incident phase space $`\mathrm{\Omega }`$. So, used in this paper transformation is the ordinary momentum mapping of classical mechanics :
$$J:(u,p)(\xi ,\eta ),$$
(1.4)
where $`p`$ is the conjugate to $`u`$ momentum, see Secs.4.1 and 4.2. The Hamiltonian description will be useful by this reason.
In this paper we restrict ourselves by one dimensional quantum mechanics assuming that
$$v(u;g)=\frac{1}{g}v(g^{1/2}u;1)\frac{1}{g}v(g^{1/2}u),$$
(1.5)
where, for simplicity, $`v(u)`$ is the potential hole with one minimum at $`u=0`$ and $`g`$ is the interaction constant. Then the nontrivial solution $`u_c`$ of eq.(1.3) would be singular at $`g=0`$:
$$u_c=O(g^{1/2}).$$
(1.6)
in the lowest order over $`j`$.
We will find that the transformed perturbation theory presents the expansion over $`1/g`$ and the expansion coefficients are simply calculable iff $`T^{}\mathrm{\Omega }`$ is the homogeneous and isotropic space
This result cardinally distinguished from the week coupling perturbation theory developed in . We will derive in Sec.3 this ordinary perturbation series over $`g`$ using the measure (1.3) to show exactly where we turn from habitual way to formulate new perturbation theory. Strictly speaking, there is not any connections among both perturbation theories and they are dual to each other, see Sec.5.
The paper is organised as follows. In Sec.2 we will find the integral representation for $`\rho `$ with measure (1.2). In Sec.3 we will describe the week-coupling perturbation theory. In Sec.4 the mapping (1.4) for the quantum system is described to show the decomposition over $`1/g`$ and the rule of calculation of corresponding coefficient will be given. In concluding Sec.5 we will offer the dynamical interpretation of new perturbation theory.
## 2 Unitary definition of the functional measure
Starting this section we will try to explain the role of unitarity in definition of functional measure, Sec.2.1. Then we will show as the d’Alembert’s variational principle may be $`derived`$ for quantum systems (Sec.2.2) and, at the end, the general properties of theory on the $`\delta `$-like measure will be offered.
### 2.1 Formulation of method
To calculate the bound state energies $`E_n`$ it is enough to consider the trace:
$$R(E)=\underset{n}{}𝑑u\frac{\psi _n(u)\psi _n^{}(u)}{EE_ni\epsilon }=\underset{n}{}\frac{1}{EE_ni\epsilon }=\mathrm{Sp}\frac{1}{E𝐇i\epsilon },$$
(2.1)
where $`𝐇`$ is the Hamiltonian operator and $`\epsilon +0`$ and the wave functions $`\psi _n`$ ortho-normalizability was used.
The semiclassical approximation leads to
$$R(E)\underset{k=0}{\overset{\mathrm{}}{}}e^{ik(S_1(u_c)\pi )}=\frac{1}{1+e^{iS_1(u_c)}},$$
(2.2)
where $`S_1(u_c)`$ is the action on the elementary (one period) closed path trajectory $`u_c=u_c(E)`$ . The position of poles in (2.2) defines the value of $`E_n`$.
Note now that
$$\frac{1}{EE_ni\epsilon }=\mathrm{P}\frac{1}{EE_n}+i\pi \delta (EE_n)\mathrm{at}\epsilon =0,$$
i.e. it is not necessary to calculate the real part since it did not contain the masurable value of energy $`E_n`$:
$$\mathrm{P}\frac{1}{EE_n}=0E=E_n.$$
By this reason, following to , we will calculate much more simple quantity
$`\epsilon \rho (E)=\epsilon {\displaystyle \underset{n_1,n_2}{}}{\displaystyle 𝑑x_1𝑑x_2\frac{\psi _{n_1}(x_1)\psi _{n_1}^{}(x_2)}{EE_{n_1}i\epsilon }\frac{\psi _{n_2}^{}(x_1)\psi _{n_2}(x_2)}{EE_{n_2}+i\epsilon }}=`$
$`=\epsilon {\displaystyle \underset{n}{}}\left|{\displaystyle \frac{1}{EE_ni\epsilon }}\right|^2={\displaystyle \frac{1}{2i}}{\displaystyle \underset{n}{}}\left\{{\displaystyle \frac{1}{EE_ni\epsilon }}{\displaystyle \frac{1}{EE_n+i\epsilon }}\right\}=`$
$`=\pi {\displaystyle \underset{n}{}}\delta (EE_n)=\pi \mathrm{Sp}\delta (E𝐇)=\mathrm{Im}R(E).`$ (2.3)
Therefore, we wish exclude from consideration the unnecessary contributions<sup>2</sup><sup>2</sup>2‘Unnecessary’ means for us the unmeasurable quantity in given experiment. with $`EE_n`$. It should be noted that we exclude continuum of contributions contained in $`\mathrm{Re}\{1/EE_ni\epsilon \}`$ and leave the set of point-like contributions $`\mathrm{Im}\{1/EE_ni\epsilon \}`$, but with infitite amplitudes.
We can find that
$$\rho (E)\underset{k=\mathrm{}}{\overset{+\mathrm{}}{}}e^{ik(S_1(u_c)\pi )}=2\pi \underset{n}{}\delta (S_1(u_c)(2n+1)\pi )$$
(2.4)
So, as in (2.2), the aim of quantum perturbation theory is to define the corrections to the phase $`S_1(u_c)`$
In terms of integrals the cancellation phenomena, shown in (2.3), looks as follows:
$$\rho (E)=\underset{n}{}\left|\frac{1}{EE_ni\epsilon }\right|^2=\underset{n}{}_0^{\mathrm{}}𝑑T_+𝑑T_{}e^{\epsilon (T_++T_{})+i(EE_n)(T_+T_{})}$$
(2.5)
To see the effect of cancellation let us introduce new time variables $`T`$ and $`\tau `$:
$$T_\pm =T\pm \tau .$$
(2.6)
The Jacobian of transformation gives: $`0T\mathrm{}`$ and $`T\tau T`$. But in the integral (2.5) $`T(1/\epsilon )\mathrm{}`$ are essential at $`\epsilon 0`$. By this reason we can put $`|\tau |\mathrm{}`$. In result,
$$\epsilon \rho (E)=2\pi \epsilon _0^{\mathrm{}}𝑑Te^{2\epsilon T}_{\mathrm{}}^+\mathrm{}\frac{d\tau }{\pi }e^{2i(EE_n)\tau }$$
(2.7)
In the last integral all contributions, except for the case $`E=E_n`$, are cancelled.
Described cancellation is not accidental, or approximate, being the consequence of optical theorem, i.e. is the consequence of unitarity condition. The $`\delta `$-likeness of measure (1.2) has the same nature as the $`\delta `$-function in the r.h.s. of (2.3), i.e. the $`\delta `$-like measure will arise when the absorption part of amplitudes is calculated.
Note, we start from claculation of modulo squire of amplitudes since we know the path integral reprisantation for them. Then, using the unitarity condition we find the correct measure for imaginary part of the amplitude.
### 2.2 Functional $`\delta `$-like measure
We will use following path-integral representation for amplitude
$$a(u_1,u_2;E)=i_0^{\mathrm{}}𝑑Te^{iET}_{u(0)=u_1}^{u(T)=u_2}Due^{iS_{C_+(T)}(u)},Du=\underset{tC_+(T)}{}\frac{du(t)}{\sqrt{2\pi }},$$
(2.8)
to calculate
$$\rho (E)=𝑑u_1𝑑u_2\left|a(u_1,u_2;E)\right|^2.$$
(2.9)
The action $`S_{C_+(T)}(u)`$ is defined on the Mills complex time contour :
$$C_\pm (T):tt\pm i\epsilon ,\epsilon +0,0tT$$
(2.10)
Inserting $`a(u_1,u_2;E)`$ into (2.9) we find:
$$\rho (E)=_0^{\mathrm{}}𝑑T_+𝑑T_{}e^{iE(T_+T_{})}_{u_+(0)=u_{}(0)}^{u_+(T_+)=u_{}(T_{})}Du_+Du_{}e^{iS_{C_+(T_+)}(u_+)iS_{C_{}(T_{})}(u_{})}$$
(2.11)
Note crucial for us the ‘closed-path’ boundary conditions:
$$u_+(0)=u_{}(0),u_+(T_+)=u_{}(T_{}).$$
(2.12)
We will introduce new variables $`T`$ and $`\tau `$, see (2.6). The integral over $`\tau `$ will be calculated perturbatively. In zero order over $`\tau `$ we would have from (2.12):
$$u_+(0)=u_{}(0),u_+(T)=u_{}(T).$$
(2.13)
It should be underlined that this is unique solution of the boundary condition (2.12) which did not contradict to the quantum uncertainty principle (other solutions of (2.12) would involve constraints for time derivatives of coordinate).
If we introduce now new coordinates $`u`$ and $`x`$:
$$u_\pm (t)=u(t)\pm x(t),$$
(2.14)
then (2.13) gives:
$$x(0)=x(T)=0$$
(2.15)
and $`u(0)`$ and $`u(T)`$ are arbitrary. We will see that this ‘minimal’ boundary condition is sufficient to define the integrals over $`\tau `$ and $`u`$.
Let us expand the closed path action
$$S_{cl}(u\pm x;T\pm \tau )(S_{C_+(T+\tau )}(u+x)S_{C_{}(T\tau )}(ux))$$
over $`\tau `$:
$$S_{cl}(u\pm x;T\pm \tau )=S_{cl}(u\pm x;T)2\tau H_T(u)2\stackrel{~}{H}_T(u;\tau ),$$
(2.16)
where the Hamiltonian at the time moment $`T`$
$$H_T(u)=\frac{}{T}S_{C_+(T)}(u).$$
(2.17)
is $`x`$ independent because of (2.15). The remainder term $`\stackrel{~}{H}_T(u;\tau )`$ contains higher powers over $`\tau `$:
$$\stackrel{~}{H}_T(u;\tau )=\underset{n=1}{\overset{\mathrm{}}{}}\frac{\tau ^{2n+1}}{(2n+1)!}\frac{d^{2n}}{dT^{2n}}H_T(u).$$
Therefore, the conditions (2.15) factorize $`\tau `$ and $`x(t)`$ dependence: the $`x`$ dependence is contained in the $`\tau `$ independent quantity $`S_{cl}(u\pm x;T)`$ only. So, we may construct the perturbation theory over $`\tau `$ and $`x`$ independently.
Let us consider now expansion over $`x`$:
$$S_{cl}(u\pm x;T)=S_{P(T)}(u)2\mathrm{R}\mathrm{e}_{C_+(T)}𝑑tx(t)\frac{\delta S_{C_+}(u)}{\delta u(t)}2\stackrel{~}{V}_T(u,x),$$
(2.18)
where the first term in this decomposition is:
$$S_{P(T)}(u)=(S_{C_+(T)}(u)S_{C_{}(T)}(u)).$$
(2.19)
If the motion is periodic then $`S_{P(T)}(u)`$ is not equal to zero even on the real time axis . In semiclassical approximation
$$S_P(T)(u_c)=kS_1(u_c),k=0,1,\mathrm{},$$
and is $`T`$ independent. The reason of this conclusion is explained in Sec.3.3. As usual,
$$2\mathrm{R}\mathrm{e}_{C_+}𝑑t=_{C_+}𝑑t+_C_{}𝑑t$$
(2.20)
since for arbitrary analytic function $`f(tC_+)=f^{}(tC_{})`$.
Following formal trick will be useful. We can write:
$$e^{2i\stackrel{~}{H}_T(u;\tau )}=\underset{n}{}\frac{\tau ^n}{n!}K_n(u;T),$$
where
$$K_n(u;T)=\frac{d^n}{d\tau _1^n}e^{2iH_T(u;\tau _1)}|_{\tau _1=0}\widehat{\tau }_1^ne^{2i\stackrel{~}{H}_T(u;\tau _1)}.$$
On other hand,
$$(2i\tau )^n=\frac{d^n}{d\epsilon ^n}e^{2i\epsilon \tau }|_{\epsilon =0}\widehat{\epsilon }^ne^{2i\epsilon \tau }.$$
Therefore,
$$e^{2i\stackrel{~}{H}_T(u;\tau )}=\underset{n}{}\frac{(\widehat{\tau }_1\widehat{\epsilon }/2i)^n}{n!}e^{2i\epsilon \tau }e^{2i\stackrel{~}{H}_T(u;\tau _1)}=e^{i\widehat{\tau }_1\widehat{\epsilon }/2}e^{2i\epsilon \tau }e^{2i\stackrel{~}{H}_T(u;\tau _1)}.$$
(2.21)
The expansion of the operator $`e^{i\widehat{\tau }_1\widehat{\epsilon }/2}`$ will generate corresponding perturbation series.
The same operator can be introduced for expansion over the local quantity $`x`$:
$$e^{2i\stackrel{~}{V}_T(u,x)}=e^{\frac{i}{2}\mathrm{Re}_{C_+}𝑑t\widehat{j}(t)\widehat{x}_1(t)}e^{2i\mathrm{Re}_{C_+}𝑑tj(t)x(t)}e^{2i\stackrel{~}{V}_T(u,x_1)}.$$
(2.22)
Note, the eqs.(2.21), (2.22) linearise the arguments of corresponding exponents. Then, using (2.16), (2.18) and (2.21), (2.22) we find that
$$\rho (E)=2\pi _0^{\mathrm{}}𝑑Te^{i𝐊(\epsilon \tau ,\mathrm{𝐣𝐱})}DM(u)\delta (E+\epsilon H_T(u))e^{iS_{P(T)}(u)}e^{2i\stackrel{~}{H}_T(u;\tau )2i\stackrel{~}{V}_T(u,x)},$$
(2.23)
where expansion over the operator
$$𝐊=\frac{1}{2}\left(\widehat{\tau }\widehat{\epsilon }+\mathrm{Re}_{C_+}𝑑t\widehat{j}(t)\widehat{x}(t)\right)$$
(2.24)
gives the perturbation series. At the very end of calculations all auxiliary variables $`\tau ,\epsilon ,j`$ and $`x`$ should be taken equal to zero.
The measure in (2.23) is defined as follows:
$$DM(u)=\underset{t}{}du\delta \left(\frac{\delta S(u)}{\delta u}j\right)=\underset{t}{}du\delta (\ddot{u}+v^{}(u)j)$$
(2.25)
and the $`\delta `$-function is defined by the equality:
$$\underset{t}{}\delta (\ddot{u}+v^{}(u)j)=_{x(0)=0}^{x(T)=0}\underset{t}{}\frac{dx}{\pi }e^{2i\mathrm{Re}{\scriptscriptstyle 𝑑tx(\ddot{u}+v^{}(u)j)}}.$$
(2.26)
Argument of this $`\delta `$-function did not contain the boundary values $`u(0)`$ and $`u(T)`$. But this is not important since to solve the second order equation (2.28) two constant of integration is necessary.
The exponent in (2.26) is equal to the sum: $`\mathrm{Re}x\mathrm{Re}(\ddot{u}+v^{}(u)j)+\mathrm{Im}x\mathrm{Im}(\ddot{u}+v^{}(u)j)`$, being defined on the complex time contour. This means that
$$\underset{t}{}\delta (\ddot{u}+v^{}(u)j)=\underset{tC}{}\delta (\mathrm{Re}\{\ddot{u}+v^{}(u)j\})\delta (i\mathrm{Im}\{\ddot{u}+v^{}(u)j\}),$$
(2.27)
where $`C=C_++C_{}`$. So, the measure (2.25) defines both the real and imaginary part of contributions.
By definition, $`(\ddot{u}+v^{}(u)j)`$ is the total force, then the product $`(\ddot{u}+v^{}(u)j)x`$ is the virtual work. In classical mechanics this work should be equal to zero, since the classical motion is time reversible (d’Alembert). Then, noting that virtual deviation is arbitrary, one finds the local condition:
$$\ddot{u}+v^{}(u)j=0.$$
(2.28)
when the motion is time reversible.
In quantum case the virtual work is not equal to zero (quantum corrections shift the energy levels), but the integral over $`x(t)`$ gives the same result (2.28). We can conclude that the unitarity condition of quantum mechanics allows to $`derive`$ the d‘Alembert’s variational principle of classics mechanics , see also .
### 2.3 Properties of theory with $`\delta `$-like measure
The eq.(2.28) should be solved expanding over $`j(t)`$:
$$u_j(t)=u_c(t)+𝑑t^{}G(t,t^{};u_c)j(t^{})+(\mathrm{higher}\mathrm{powers}\mathrm{of}j)$$
(2.29)
where $`u_c`$ is the solution of homogeneous equation:
$$\ddot{u}+v^{}(u)=0$$
(2.30)
and $`G(t,t^{};u_c)`$ is the Green function:
$$(_t^2+v^{\prime \prime }(u_c))G(t,t^{};u_c)=\delta (tt^{}).$$
(2.31)
The eq.(2.30) have in our case the trivial constant solution
$$u_0:\dot{u}_0=0,v^{}(u_0)=0$$
(2.32)
and nontrivial one
$$u_c=u_c(t):\dot{u}_c(t)0,\ddot{u}_c+v^{}(u_c)0.$$
(2.33)
Because of definition of the $`\delta `$-function and since there is not any special restriction on the contributions both one should be taken into account:
$$\rho (E)=\rho _0(E)+\rho _c(E).$$
(2.34)
This means that one should sum over all possible topological classes of trajectory, if the single class is unable to cover all phase space. Each class of trajectories belongs to restricted domain of phase space: $`\mathrm{\Omega }=W^0\times W^c`$ in our case. Each sub-domain $`W^i`$ is restricted by the bifurcation lines . This means that $`\rho _0`$ can not be achieved by analytical continuation of $`\rho _c`$ (for instance, taking $`E=0`$ in $`\rho _c`$ for the semiclassical approximation).
It is evident, one should leave in the sum (2.34) the term with higher volume $`V_{W^i}`$, $`i=0,c`$. This is just the domain of $`u_cW^c`$ trajectory, and one can put out $`\rho _0`$ since the sub-domain of $`u_c^0W^0`$ is the point . Indeed, it will be shown that $`\rho _cV_{tr}=\mathrm{}`$, where $`V_{tr}`$ is the volume of time translations mode (zero frequency mode). At the same time, $`\rho _0O(1)`$. Therefore, iff the time translation invariance is unbroken, one can say that $`\rho _0`$ is realised on the measure $`O(1)/V_{tr}=0`$.
The ability to classify contributions by the trajectory topology classes becomes possible since there is not in (2.34) the $`u_c^0`$ and $`u_c`$ interference term. This is evident consequence of the orthogonality of corresponding Hilbert spaces. Therefore, the choice of solution means choice of corresponding vacuum.
## 3 WKB perturbation theory
It can be shown that (2.23) restores ordinary WKB expansion. The first step of this calculations is to find the solution of inhomogeneous equation (2.28), Sec.3.1. Then we may find that this perturbation theory counts positive powers of $`g`$, Sec.3.2. At the end the zero frequency modes problem will be discussed.
### 3.1 Tree decomposition
Let us consider the tree decomposition (2.29) more carefully for the potential
$$v(u;g)=\frac{1}{2}w_0^2u^2+\frac{1}{4}gu^4.$$
(3.1)
It is evident:
$$v(u;g)=\frac{1}{g}v(g^{1/2}u)$$
(3.2)
The decomposition (2.29) can be written in the form:
$$u_j(t)=u_c(t)+\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{n!}\underset{k=1}{\overset{n}{}}\left\{dt_ij(t_i)\right\}G_n(t,t_1,\mathrm{}.,t_n;u_c)$$
(3.3)
It easy to show that the $`n`$-point Green function
$$G_n=O(g^{(n1)/2}).$$
(3.4)
Indeed, inserting (3.3) into the equation:
$$\ddot{u}+\omega _0^2u+gu^3=j$$
(3.5)
we find:
$$(_t^2+\omega _0^2+3gu_c^2)G_1(t,t_1;u_c)=\delta (tt_1)$$
(3.6)
The operator $`(_t^2+\omega ^2+3gu_c^2)`$ is $`g`$ independent since $`u_c=O(1/g^{1/2})`$ and, therefore, $`G_1=O(g^0)`$.
Note also, the operator ($`_t^2+\omega ^2+3gu_c^2`$) is translationally noninvariant. By this reason considered perturbation theory is sufficiently complicated so that only first corrections has been computed till now.
The equation for $`G_2`$ have the form:
$$(_t^2+\omega ^2+3gu_c^2)G_2(t,t_1,t_2;u_c)+6gu_cG_1(t,t_1;u_c)G_1(t,t_2;u_c)=0.$$
(3.7)
Therefore, in accordance with (3.4), $`G_2=O(g^{1/2})`$. In result, the analysis of higher orders over $`j`$ would justify (3.4).
The interactions generating functional $`V_T`$ computed for the case (3.1) has the form:
$$\stackrel{~}{V}_T(u,x)=2g\mathrm{Re}_{C_+}𝑑tx^3(t)u(t)+O(\epsilon ),$$
(3.8)
where the $`O(\epsilon )`$ term is proportional to the imaginary part of $`S_{cl}`$.
The operator $`𝐊`$ is linear over $`\widehat{x}=\delta /\delta x`$. Therefore, action of $`\mathrm{exp}\{i𝐊\}`$ will give:
$$\rho _c:e^{2i\stackrel{~}{𝐕}_T(u,\widehat{j}/2i)}:e^{iS_{P(T)}(u)}e^{2i\stackrel{~}{H}_T(u;\tau )}\delta (E+\epsilon H_T(u_j)),$$
(3.9)
where the colons prescribe normal product, when the operator should stay to the left of all functions on which it may act, and the unimportant for present analyses integrations were not mentioned.. The expansion of the operator exponent gives the perturbation series:
$$\rho _c\underset{n}{}\frac{(2i)^n}{n!}:\stackrel{~}{𝐕}_T^n(u_j,\widehat{j}/2i):e^{iS_{P(T)}(u_j)}e^{2i\stackrel{~}{H}_T(u_j;\tau )}\delta (E+\epsilon H_T(u_j)).$$
(3.10)
Let us consider the self-interaction part for the beginning. This means that the shifting energy levels renormalisation of $`S_P`$ and $`\stackrel{~}{H}_T`$ is not considered. In other words, we omit the action of operators $`\widehat{j}(t_i)`$ on $`\mathrm{exp}\{iS_P(u_j)2i\stackrel{~}{H}_T(u_j;\tau )\}`$:
$$\rho _c\delta (E+\epsilon H_T(u_j))e^{iS_{P(T)}(u_c)}e^{2i\stackrel{~}{H}_T(u_c;\tau )}\underset{n}{}\frac{(21)^n}{n!}:\stackrel{~}{𝐕}_T^n(u_j,\widehat{j}/2i):.$$
Then the lowest order contribution is $`\stackrel{~}{𝐕}`$. So, in the first order we find:
$$\widehat{j}^3u_jG_3=O(g^2),$$
(3.11)
where the estimation (3.4) was used and the prescription that the auxiliary variable $`j`$ should be taken equal to zero was taken into account. In the second order $`\stackrel{~}{𝐕}_T^2=O(\widehat{j}^6)`$ contribution have following order over $`g`$:
$$g^2\widehat{j}^6u_j^2=O(g^4),$$
(3.12)
and so on.
In result, one can find that the $`n`$-th order in expansion of $`\mathrm{exp}\{i𝐊\}`$ gives $`O(g^{2n})`$ expansion if the self-interactions only are included in $`\rho (E)`$.
As follows from decomposition (3.3) and estimation (3.4) the action of operator $`\widehat{j}`$ on $`u_j`$ gives coefficient $`g^{1/2}`$. The renormalisation of $`S_P`$ and $`\stackrel{~}{H}_T`$ start from $`1/g`$ terms, but higher orders would contain the positive powers of $`g`$ since they are produced by the actions of $`\widehat{j}(t_i)`$.
### 3.2 Connection with WKB expansion
One can show another argument that considered above perturbation theory is nothing new but is the ordinary expansion around $`u_c`$ developed in early publications . Let us use for this purpose the substitution:
$$u(t)u_c(t)+u(t)$$
(3.13)
Then
$`\rho _c(E)=2\pi {\displaystyle _0^{\mathrm{}}}dTe^{i𝐊(\epsilon \tau ,\mathrm{𝐣𝐱})}{\displaystyle }DM(u_c,u)\delta (E+\epsilon H_T(u_c+u))\times `$
$`\times e^{iS_{P(T)}(u_c+u)}e^{2i\stackrel{~}{H}_T(u_c+u;\tau )2i\stackrel{~}{V}_T(u_c+u,x)},`$ (3.14)
where
$$DM(u_c,u)=\underset{t}{}du\delta \left(\frac{\delta S(u_c+u)}{\delta u}+j\right)$$
(3.15)
We should take into account that $`u_c`$ depends on the integration constants $`\xi `$ and $`\eta `$. Therefore, if $`(\xi ,\eta )`$ form the manifold $`W^c`$, as was mentioned in Sec.2.3, one should sum over all solutions $`u_cW^c`$, see Sec.3.3.
We want to show now that (3.14) may be reduced to the product of two path integrals. Indeed, using (2.26) and (2.16), (2.18) we find from (3.14) that
$`\rho _c(E)=2{\displaystyle _0^{\mathrm{}}}dT{\displaystyle _{\mathrm{}}^+\mathrm{}}d\tau ^{}e^{i𝐊}{\displaystyle }DuDx^{}e^{2i(E+\epsilon H_T(u_c+u))\tau ^{}}e^{2iH_T(u_c+u)\tau }\times `$
$`\times e^{S_{cl}(u_c+u\pm x^{};T\pm \tau ^{})}e^{2i\mathrm{Re}{\scriptscriptstyle 𝑑tx(S(u_c+u)/u)}}e^{2i\mathrm{Re}{\scriptscriptstyle 𝑑tx\{(S(u_c+u)/u)+j\}}}.`$ (3.16)
The action of operator $`\mathrm{exp}\{i𝐊\}`$ leads to substitutions: $`xx^{},\tau \tau ^{}`$ and $`\epsilon 0,j0.`$ Taking this into account we find:
$$\rho _c(E)=\left|_0^{\mathrm{}}𝑑Te^{iET}Due^{iS_{C_+}(u_c+u)}\right|^2,$$
(3.17)
where the functional integral should be calculated perturbatively over $`u`$. Note, calculation of amplitudes is useful since eliminates the doubling of degrees of freedom.
### 3.3 Zero modes
The defined by eq.(3.17) $`\rho _c(E)`$ stay undefined till the procedure of summation over all $`u_cW`$ is not formulated. Following to the equality:
$$\underset{\{u_c\}}{}=_W𝑑\xi 𝑑\eta \sigma (u;\xi ,\eta )$$
we should define the density $`\sigma (u;\xi ,\eta )`$ of states in the domain $`(\xi ,\xi +d\xi ;\eta ,\eta +d\eta )`$. The Faddeev-Popov $`ansatz`$ is used for this purpose .
By definition, $`(\xi ,\eta )`$ are the constants of integration and they may be chosen arbitrarily. For example, wee may take $`(\xi ,\eta )`$ as the initial coordinate and momentum of particle on the trajectory $`u_c`$.
But the ‘field-theoretical’ definitions would be much more useful for us, see Sec.4. One may note that the dependence on $`(\xi ,\eta )`$ indicates the symmetry breaking. Then $`\eta `$ may be taken as the generator $`J`$ of broken symmetry and $`\xi `$ as the canonically conjugate to it coordinate $`\mathrm{\Theta }`$. It will be important for us that $`(\xi ,\eta )`$ define the solution $`u_c`$ unambiguously. In other words, we will use the ordinary mechanical statement that $`(\xi ,\eta )`$ form a manifold $`W^c`$ and $`u_c`$ belongs to it $`completely`$. So, we would assume that the equations:
$$\xi =\mathrm{\Theta }(u_c,\dot{u_c}),\eta =J(u_c,\dot{u_c})$$
(3.18)
define the integration constants of $`u_c`$ unambiguously.
Then, to define the density $`\sigma `$, we may insert into the initial representation (2.23) the unite (Faddeev-Popov $`ansatz`$):
$$1=_W\underset{t}{}d\xi d\eta \delta (\xi \mathrm{\Theta }(u,\dot{u}))\delta (\eta J(u,\dot{u})$$
(3.19)
Note, by definition $`\eta `$ should coincide with conserved generator. But nevertheless we consider $`\eta =\eta (t)`$ and the same for $`\xi `$. This assumption is necessary since the quantum case is considered.
We can change order of integration and integrate firstly over $`u`$ using the $`\delta `$-function of the measure $`DM`$. Lagrange equation (2.30) should be solved taking into account the constraints (3.18):
$$\rho _c(E)=2\pi _W𝑑\xi (0)_0^{\mathrm{}}𝑑Te^{i𝐊(\epsilon \tau ,\mathrm{𝐣𝐱})}DM_c(u)\delta (E+\epsilon H_T(u))e^{iS_{P(T)}(u)}e^{2i\stackrel{~}{H}_T(u;\tau )2i\stackrel{~}{V}_T(u,x)},$$
(3.20)
where $`\xi (0)`$ is the initial phase and the constraint measure
$$DM_c(u)=\underset{t}{}du\delta \left(\frac{\delta S(u)}{\delta u}+j\right)\delta \left(\xi _0\mathrm{\Theta }(u,\dot{u})\right)$$
(3.21)
was introduced. In our problem the value of $`J`$ is restricted by $`\delta (E+\epsilon H_T(u))`$. It was used in (3.20) that $`DM_c`$ is $`\xi `$ independent since Lagrange equation is invariant against $`\xi `$ variations and $`j=j(t)`$ is the auxiliary variable. This means that $`\rho _c(E)`$ defined in (3.20) is proportional to the volume
$$V_{tr}=𝑑\xi (0)$$
of the time translation mode.
It is important here to trace on the following question. One can note that (3.20) gives $`\rho _cV_{tr}^1`$. On other hand, as follows from (3.17), one may expect $`\rho _cV_{tr}^2`$. It is evident that this discrepancy is the consequence of loaded into formalism condition of the orthogonality of Hilbert spaces, see Sec.2.3.
Remembering definition of $`\rho `$ as the squire of amplitudes, we may insert the Faddeev-Popov’s unite defined on the whole time contour $`C=C_++C_{}`$, see (2.27), to take into account the input condition that the trajectories $`u_+(tC_+)`$ and $`u_{}(tC_{})`$ are absolutely independent. This means that, generally speaking, the boundary conditions for this trajectories should not coincide and, therefore, if we introduce $`\xi (tC_\pm )|_{t=0}\xi _\pm `$, one should have in mind that, generally speaking, $`\xi _+\xi _{}`$. Then integration over $`\xi _+`$ and $`\xi _{}`$ should be performed independently.
But we have considered the closed-path contributions, see (2.13). This gives restriction for the $`u_\pm `$ boundary properties. Taking into account (2.15), we can find, considering the periodic orbits, that
$$\xi _+=\xi _{}\pm kP_1(E),k=0,1,2,\mathrm{},$$
(3.22)
where $`P_1(E)`$ is the elementary period. Just this solution leads to $`S_P0`$ and the necessary summation over $`k`$ gives the energy levels quantisation condition (2.4), see also .
## 4 Mapping on the cotangent manifold
The necessity to search a new form of the perturbation theory is caused by extremal complexity of the WKB perturbation theory described above.
The quantum nature of collective variables $`(\xi ,\eta )`$ was mentioned previously by many authors . We would like continue this idea considering them as a new quantum variables. For this purpose we will use the $`\delta `$-like definition of measure, the definition of the interactions generating functional $`\stackrel{~}{V}_T`$ and the perturbations generating functional $`\mathrm{exp}\{𝐊\}`$, to count the possible excitations of the field $`u_cW^c`$, see Sec.4.1. In Sec.4.2 we will show the structure of new perturbation theory.
### 4.1 Procedure of mapping
Let us return to the Faddeev-Popov unite
$$1=D\xi D\eta \underset{t}{}\delta (\xi \mathrm{\Phi }(u,\dot{u}))\delta (\eta J(u,\dot{u}))$$
(4.1)
It is assumed, as was offered in Sec.3.3,
$$\underset{t}{}\underset{tC=C_++C_{}}{}.$$
The first order formalism will be useful for us . Corresponding measure
$$DM(u,p)=\underset{t}{}dudp\delta \left(\dot{u}\frac{H_j(u,p)}{p}\right)\delta \left(\dot{p}+\frac{H_j(u,p)}{u}\right),$$
(4.2)
where the total Hamiltonian
$$H_j(u,p)=\frac{1}{2}p^2+v(u)ju$$
(4.3)
includes the energy of quantum excitations $`ju`$. It is evident that the integration over $`p`$ gives incident measure (2.25).
Inserting (4.1) into the functional integral with measure (4.2) we find that we have four equations for $`u`$ and $`p`$:
$$\dot{u}=\frac{H_j(u,p)}{p},\dot{p}=\frac{H_j(u,p)}{u}$$
(4.4)
and
$$\xi (t)=\mathrm{\Phi }(u,\dot{u}),\eta (t)=J(u,\dot{u}).$$
(4.5)
In previous section the first pare of equations (4.4) was used to calculate the functional integral.
But now we would like use second one (4.5). It is possible iff $`u_c`$ belongs to the space $`W^c`$ completely and $`W^c`$ is a manifold. This condition means that the eqs.(4.5) have unique solution $`(u_c,p_c)`$ and this solution transform (4.4) into identity at least at $`j=0`$. Let $`u_c(\xi ,\eta )`$ and $`p_c(\xi ,\eta )`$ are the solutions of (4.5). One can recognise in our description the ordinary canonical transformation (1.4), i.e. it defines the cotangent foliation $`W^c=T^{}\mathrm{\Omega }`$. But eqs.(4.4) and (4.5) should be solved simultaneously. So, inserting $`u_c,p_c`$ into (4.4) we should use the ‘excited’ by $`j`$ solutions $`\xi _j(t)`$ and $`\eta _j(t)`$.
So, we wish to adopt the statement that the random (Gaussian) walk, induced by the same operator $`\mathrm{exp}\{i𝐊\}`$, covers both $`W^c=(\xi ,\eta )`$ and $`\mathrm{\Omega }=(u,p)`$ spaces densely. By this reason one may choose one of them arbitrarily.
The corresponding Jacobian of transformation $`\mathrm{\Delta }`$ is $`\delta `$-like:
$$\mathrm{\Delta }=\underset{t}{}\delta (\dot{u}_c\frac{H_j(u_c,p_c)}{p_c})\delta (\dot{p}_c\frac{H_j(u_c,p_c)}{u_c}),$$
(4.6)
and
$$det^1(u_c,p_c)=\underset{t}{}dudp\delta (\xi \mathrm{\Phi }(u,\dot{u}))\delta (\eta J(u,\dot{u}))=1$$
(4.7)
since the transformation is canonical. This allows to diagonalise $`\mathrm{\Delta }`$ and mapping into the $`W^c`$ space leads to following path integral representation:
$$\rho _c(E)=2\pi 𝑑Te^{i𝐊}DM(\xi ,\eta )\delta (E+\epsilon h(\xi ,\eta ;T))e^{iS_{P(T)}(u_c)}e^{2i\stackrel{~}{h}(u_c;\tau ,T)2i\stackrel{~}{V}_T(u_c,x)},$$
(4.8)
where the measure
$$DM(\xi ,\eta )=\underset{t}{}d\xi d\eta \delta \left(\dot{\xi }\frac{h_j(\xi ,\eta )}{\eta }\right)\delta \left(\dot{\eta }+\frac{h_j(\xi ,\eta )}{\xi }\right)$$
(4.9)
and $`h_j`$ is the transformed Hamiltonian:
$$h_j(\xi ,\eta )=h(\eta )ju_c(\xi ,\eta ).$$
(4.10)
In result of mapping the problem of calculation of functional integral was reduced to solution of equations:
$$\dot{\xi }=\frac{h_j(\xi ,\eta )}{\eta }=\omega (\eta )j\frac{u_c(\xi ,\eta )}{\eta },\dot{\eta }=\frac{h_j(\xi ,\eta )}{\xi }=j\frac{u_c(\xi ,\eta )}{\xi },$$
(4.11)
where one can choose, for example,
$$\omega (\eta )=\frac{h(\eta )}{\eta }=1.$$
(4.12)
This means that in this case
$$\eta =H(u,p),\xi =^u\frac{dy}{\sqrt{2(Hv(y))}}.$$
(4.13)
It is evident that the solution of this equations gives $`u_c(\xi ,\eta )`$ and $`p_c(\xi ,\eta )`$ unambiguously.
### 4.2 Structure of transformed perturbation theory
We want to show now that $`\rho (E)`$, defined in (4.8), has the strong coupling expansion. Let us start for this purpose from the ‘tree decomposition’ of the equations (4.11):
$$\dot{\xi }=\frac{h_j(\xi ,\eta )}{\eta }=1j\frac{u_c(\xi ,\eta )}{\eta },\dot{\eta }=\frac{h_j(\xi ,\eta )}{\xi }=j\frac{u_c(\xi ,\eta )}{\xi }.$$
(4.14)
We will consider following decomposition of the solutions $`\xi _j`$ and $`\eta _j`$:
$`\xi _j(t)=\xi _0(t)+{\displaystyle \underset{n}{}}{\displaystyle \frac{1}{n!}}{\displaystyle \{dt_ij(t_i)\}\xi _n(t;t_1,\mathrm{},t_n)},`$
$`\eta _j(t)=\eta _0(t)+{\displaystyle \underset{n}{}}{\displaystyle \frac{1}{n!}}{\displaystyle \{dt_ij(t_i)\}\eta _n(t;t_1,\mathrm{},t_n)}.`$ (4.15)
Inserting (4.15) into the (4.14) we find equation for the $`n`$-point Green functions $`\xi _n(t;t_1,\mathrm{},t_n)`$ and $`\eta _n(t;t_1,\mathrm{},t_n)`$. It can be shown:
$$\xi _n=O(g^{n/2}),\eta _n=O(g^{n/2}).$$
(4.16)
Indeed, in zero order over $`j`$ we have:
$$\xi _0=\xi (0)+t,\eta _0=\eta (0)$$
(4.17)
since $`W^c`$ is the homogeneous and isotropic manifold. Then, in the first order over $`j`$:
$$\dot{\xi }_1(t;t_1)=\delta (tt_1)\frac{u_c(\xi _0(t),\eta _0)}{\eta _0}=O(g^{1/2}),\dot{\eta }_1(t;t_1)=\delta (tt_1)\frac{u_c(\xi _0(t),\eta _0)}{\xi _0}=O(g^{1/2})$$
(4.18)
since the derivatives of $`u_c`$ are unable to change the $`g`$ dependence. In second order we have the equations:
$`\dot{\xi }_2(t;t_1,t_2)=\delta (tt_1)\left\{\xi _1(t;t_2){\displaystyle \frac{u_c(\xi _0(t),\eta _0)}{\eta _0\xi _0}}+\eta _1(t;t_2){\displaystyle \frac{u_c(\xi _0(t),\eta _0)}{\eta _0\eta _0}}\right\}=O(g^1),`$
$`\dot{\xi }_2(t;t_1,t_2)=\delta (tt_1)\left\{\xi _1(t;t_2){\displaystyle \frac{u_c(\xi _0(t),\eta _0)}{\xi _0\xi _0}}+\eta _1(t;t_2){\displaystyle \frac{u_c(\xi _0(t),\eta _0)}{\xi _0\eta _0}}\right\}=O(g^1)`$ (4.19)
And so on. So, each power of $`\widehat{j}`$ adds $`g^{1/2}`$. This proves the estimation (4.16). It is important to note that eqs.(4.19) are trivially integrable. Therefore, we can calculate $`(\xi _n,\eta _n)`$ for arbitrary $`n`$.
The operator $`𝐊`$ is linear over $`\widehat{x}`$. So, the result of its action gives the normal ordered structure:
$$:e^{2i\stackrel{~}{𝐕}_T(u_c,\widehat{j}/2i)}:e^{iS_{P(T)}(u_c)}e^{2i\stackrel{~}{H}_T(u_c;\tau )}\delta (E+\epsilon h(\xi _j,\eta _j;T)),$$
(4.20)
and at the very end of calculations one should take the auxiliary variables $`x`$ equal to zero.
Let us consider once more the $`gu^4`$ theory. Then,
$$\stackrel{~}{V}_T(u_c,\widehat{j}/2i)=O(\widehat{j}^3),$$
Therefore, leaving the self-interaction parts only, in the lowest order we would have the contribution
$$\stackrel{~}{V}_T(u_c,\widehat{j}/2i)g\widehat{j}^3u_c=O(1/g)$$
(4.21)
where (4.16) was used. So, the lowest order of new perturbation theory is $`1/g`$. In result, the $`n`$-th order is $`\stackrel{~}{V}_T(u_c,\widehat{j}/2i)^n1/g^n`$.
The action of $`\stackrel{~}{V}_T^n(u_c,\widehat{j}/2i)`$ on $`e^{iS_{P(T)}(u_c)}e^{2i\stackrel{~}{H}_T(u_c;\tau )}\delta (E+\epsilon h(\xi _j,\eta _j;T))`$ did not alter this conclusion since the derivative of $`u_c`$ can not change the $`g`$ dependence.
## 5 Conclusion
We conclude this paper by notation that it is impossible the transformed theory reduce to amplitude representation. Indeed, let us return to (4.8) and use the Fourier definition of the $`\delta `$-functions:
$`\rho _c(E)={\displaystyle _0^{\mathrm{}}}dT{\displaystyle _{\mathrm{}}^+\mathrm{}}d\tau ^{}e^{i𝐊}{\displaystyle }D\xi D\eta Dx_\xi Dx_\eta e^{2i(E+\epsilon h(\eta ;T))\tau ^{}}\times `$
$`\times e^{2i\mathrm{Re}_{C_+}𝑑tx_\xi \{\dot{\xi }h_j(\xi ,\eta )/\eta \}}e^{2i\mathrm{Re}_{C_+}𝑑tx_\eta \{\dot{\eta }+h_j(\xi ,\eta )/\xi \}}e^{iS_{P(T)}(u_c)}e^{2i\stackrel{~}{h}(u_c;\tau ,T)2i\stackrel{~}{V}_T(u_c,x)},`$ (5.1)
where, see (2.16),
$$2\stackrel{~}{h}(u_c;\tau ,T)=S_{cl}(u_c\pm x;T\pm \tau )S_{cl}(u_c\pm x;T)+2\tau h(\eta ).$$
(5.2)
Using this definition, and remembering that the action of operator $`\mathrm{exp}\{i\widehat{e}\widehat{\tau }/2\}`$ gives $`\tau =\tau ^{}`$ and $`\epsilon =0`$, we find:
$`\rho _c(E)={\displaystyle _0^{\mathrm{}}}dT{\displaystyle _{\mathrm{}}^+\mathrm{}}d\tau e^{2iE\tau }e^{i\mathrm{Re}_{C_+}𝑑t\widehat{j}\widehat{x}/2}{\displaystyle }D\xi D\eta Dx_\xi Dx_\eta e^{iS_{cl}(u_c\pm x;T\pm \tau )iS_{cl}(u_c\pm x;T)}\times `$
$`\times e^{2i\mathrm{Re}_{C_+}𝑑tx_\xi \delta S(u_c)/\delta \eta }e^{2i\mathrm{Re}_{C_+}𝑑tx_\eta \delta S(u_c)/\delta \xi }e^{iS_{P(T)}(u_c)}e^{2i\stackrel{~}{V}_T(u_c,x)},`$ (5.3)
if the transformed action
$$S(u_c)=𝑑t\{\eta \dot{\xi }h(\eta )\}.$$
Action of the perturbation generating operator gives:
$`\rho _c(E)={\displaystyle _0^{\mathrm{}}}dT{\displaystyle _{\mathrm{}}^+\mathrm{}}d\tau e^{2iE\tau }{\displaystyle }D\xi D\eta Dx_\xi Dx_\eta e^{iS_{cl}(u_c\pm x_c;T\pm \tau )}\times `$
$`\times e^{2i\mathrm{Re}_{C_+}𝑑t\{x_\xi (\delta S(u_c)/\delta \eta )x_\eta (\delta S(u_c)/\delta \xi )\}}e^{2i\mathrm{Re}_{C_+}𝑑tx_c(\delta S(u_c)/\delta u_c)},`$ (5.4)
if (2.18) is used and, using the local coordinates of the $`W`$ space,
$$x_c=x_\xi \frac{u_c}{\eta }x_\eta \frac{u_c}{\xi }=\delta u_c\delta p_c$$
(5.5)
Now, if
$$\frac{\delta S(u_c)}{\delta \xi }=\frac{u_c}{\xi }\frac{\delta S(u_c)}{\delta u_c},\frac{\delta S(u_c)}{\delta \eta }=\frac{u_c}{\eta }\frac{\delta S(u_c)}{\delta u_c},$$
(5.6)
then we can write:
$$\rho _c(E)=_0^{\mathrm{}}𝑑T_{\mathrm{}}^+\mathrm{}𝑑\tau e^{2iE\tau }D\xi D\eta Dx_\xi Dx_\eta e^{iS_{cl}(u_c\pm x_c;T\pm \tau )}.$$
(5.7)
The quantities $`(x_\xi ,x_\eta )`$ and $`(\xi ,\eta )`$ have different meaning. First ones are the virtual variation of the ‘field’ $`u`$ along the corresponding axis of $`W^c`$ space, and the integrals over them should be calculated perturbatively, but last ones are the axis of the $`W^c=T^{}\mathrm{\Omega }`$ phase space. The closed path action
$$S_{cl}(u_c\pm x_c;T\pm \tau )=S_{C_+(T+\tau )}(u_c(\xi ,\eta )+x_c(\xi ,\eta ;x_\xi ,x_\eta ))S_{C_{}(T\tau )}(u_c(\xi ,\eta )x_c(\xi ,\eta ;x_\xi ,x_\eta )).$$
(5.8)
It is evident from (5.7) the transformed representation can not be written in the factorized form of product of two amplitudes.
We interpret this conclusion as impossibility of the canonical transformations in the path integrals (2.8) since on the cotangent manifolds the quantum excitations induce the phase space flows in which all degrees of freedom are mixed.
In traditional terms this means the problem of time ordering of nonlinear operators. Our success is based on the observation that the unitarity condition unambiguously defines the perturbation theory in the (‘linear’) representation, where we may disentangle all time orderings. Fixing this procedure in the structure of $`DM`$, $`\beta K`$ and $`\stackrel{~}{V}_T`$ one can do arbitrary transformations. But, the payment for this success is necessity to work in terms of less habitual absorption part of amplitude and, by this reason, one should be careful interpreting our perturbation theory as a general, see .
But, in conclusion, quantising the nonlinear waves our strong coupling perturbation theory seems much more attractive since we can perform the calculation in this theory up to the end, choosing $`W^c`$ as the homogeneous and isotropic space.
Acknowledgement
We would like to thank V.Kadyshevsky for interest to described perturbation theory.
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# DESY 00–039 ISSN 0418–9833 MZ-TH/00–06 TTP00–02 hep–ph/0003082 March 2000 Photon Plus Jet Cross Sections in Deep Inelastic 𝑒𝑝 Collisions at Order 𝑂(𝛼²𝛼_𝑠)
## 1 Introduction
In the past, measurements of prompt photon production at both fixed-target facilities and hadron-hadron colliders, have extensively been used to constrain the gluon distribution of the proton . Only recently the first data on prompt photon production in high-energy $`ep`$ collisions have been reported . Due to the presently still limited statistics the measurements are confined to prompt photons produced in photoproduction reactions, i.e. to $`ep`$ collisions with almost real exchanged photons ($`Q^20`$). As is well-known, photoproduction processes at high energies proceed by two distinct mechanisms. The incoming photon can couple either in a point-like manner to the incoming quark or antiquark (direct process) or hadronically as a source of quarks and gluons which in turn take part in the subsequent hard scattering process (resolved process). Therefore an important advantage of photoproduction measurements is to provide additional constraints on the quark and gluon content of the photon as suggested many years ago by Aurenche et al. .
By contrast, prompt photon production in deep inelastic scattering with large $`Q^2`$, $`Q^2\mathrm{\Gamma }>\mathrm{\hspace{0.17em}10}`$ GeV<sup>2</sup>, is fully determined by the direct process and does not need any non-perturbative input for the parton content of the photon<sup>1</sup><sup>1</sup>1At intermediate $`Q^2<10`$ GeV<sup>2</sup> there might be a significant resolved contribution which depends on the quark and gluon distribution of the virtual photon. as in photoproduction. Therefore this process is sensitive only to the parton distribution functions (PDF’s) of the proton. The possible information on the proton PDF’s would be complementary to the $`F_2`$ measurement from inclusive deep inelastic scattering, since up\- and down-type quarks contribute with different weights. Of course, the cross sections for $`epe\gamma X`$ at large $`Q^2`$ are smaller than the corresponding cross sections for almost real photons; but with the larger luminosities planned at HERA rather accurate measurements of various differential cross sections for $`Q^2>10`$ GeV<sup>2</sup> seem feasible.
In addition to the perturbative direct production, photons are also produced through the “fragmentation” of a hadronic jet into a single photon carrying a large fraction of the jet energy . This long-distance process is described in terms of the quark-to-photon and gluon-to-photon fragmentation functions (FF’s).
In order to unambiguously identify the prompt photon signal from the hadronic background it is necessary to apply isolation cuts in the experiment. This has the effect of reducing the cross section due to a suppression of the photon fragmentation contributions. On the other hand, it has the advantage of eliminating to a large extent the dependence on the photon fragmentation function, which is a non-perturbative input and must come from other experiments designed to measure them.
To obtain reliable predictions it is necessary to calculate the $`epe\gamma X`$ cross section in next-to-leading order (NLO) of the strong coupling constant $`\alpha _s`$ as has been done for $`p\overline{p}`$ collisions , $`e^+e^{}`$ annihilation , as well as photoproduction . The corresponding NLO calculation of cross sections for prompt photon production in $`ep`$ scattering with large $`Q^2`$ has not yet been done previously, neither for the technically simpler case of inclusive cross sections, i.e. without any photon isolation cut, nor for the case with isolation cuts. Applicability of perturbative QCD requires that the scattering process is characterized by a large transverse momentum, provided either by the momentum transfer $`Q^2`$ or a large transverse momentum of the hadronic final state. We consider only the case where both $`Q^2`$ is large and the hadronic final state is characterized by a large $`p_T`$. One specific possibility is to consider the case where in addition to the photon also one or more jets are observed in the final state. The detection of an additional jet may also help to identify the prompt photon events in the actual experiment. In leading order (LO) the photon is produced by the Compton process $`\gamma ^{}+q\gamma +q`$, where $`\gamma ^{}`$ is a photon of high virtuality emitted by the incoming electron. This partonic photon production process contributes to the $`\gamma +(1+1)`$-jet final state in $`ep`$ scattering (the proton remnant jet being counted as “$`+1`$” jet as usual). In NLO the final states are $`\gamma +(1+1)`$\- and $`\gamma +(2+1)`$-jets.
The first NLO calculations for this case were done by two of us and D. Michelsen . This calculation was restricted to the case of not too large $`Q^2`$ where it is possible to neglect the exchange of a $`Z`$ boson. Moreover, in this previous work the fragmentation contribution was not taken into account. Therefore, photon-quark collinear singularities could not be absorbed into the fragmentation functions. Instead, these singularities had been removed by explicit parton-level cutoffs. As a consequence, the result depended strongly on these parton-photon cutoffs, in particular for subprocesses with an incoming gluon. These cutoffs are difficult to control experimentally, where hadrons combined into jets are observed and not the partons needed to define the cutoffs. In subsequent work we included the fragmentation contribution thereby avoiding the need to use parton-level cutoffs. The isolation criteria, which limit the hadronic energy in the jet containing the photon, are thus physical, i.e. correspond to selection criteria in the experimental analysis. In a later paper we studied the sum of the $`\gamma +(1+1)`$-jet and $`\gamma +(2+1)`$-jet cross sections as a function of the momentum fraction carried by a photon inside a jet. We observed that this special cross section is sensitive to the fragmentation contribution, in particular to the quark-to-photon FF.
In , only a few observables have been calculated, as for example distributions with respect to the transverse momentum, $`p_T`$, and the rapidity, $`\eta `$, of the photon or the most energetic jet for one particular choice of the photon isolation cut. In this paper we take up the topic of this earlier work. Besides several other observables which are of interest for analyzing upcoming experimental data from HERA we shall present results for $`p_T`$ and $`\eta `$ distributions already considered in for different isolation cuts. We also study possible scale dependencies to estimate the reliability of our predictions. As in we use the $`\gamma ^{}p`$ center-of-mass system to define the kinematic variables. For most of the cross section calculations a particular cone algorithm is applied to define the parton jets and to isolate the photon signal. For a few cases we shall also make use of the $`k_T`$ cluster algorithm.
The plan of the paper is as follows. In section 2 a brief outline of the theoretical background to the cross section calculations as well as the technique of the calculation are given. The results for the various observables are presented and discussed in section 3. Section 4 contains a summary and some concluding remarks.
## 2 Subprocesses Through Next-to-Leading Order
### 2.1 Leading-Order Subprocesses
In leading order, the production of photons in deep inelastic electron scattering is described by the quark (antiquark) subprocess (see Fig. 1)
$$e(p_1)+q(p_3)e(p_2)+q(p_4)+\gamma (p_5)$$
(1)
where the particle momenta are given in parentheses. The momentum of the incoming quark is a fraction $`\xi `$ of the proton momentum $`P`$, $`p_3=\xi P`$. The proton remnant $`r`$ has the momentum $`p_r=(1\xi )P`$. It hadronizes into the remnant jet so that the process (1) gives rise to $`\gamma +(1+1)`$-jet final states. In the virtual photon-proton center-of-mass system the hard photon recoils against the hard jet back-to-back. In a leading-logarithmic calculation, the effects of higher-order processes show up only via the use of the scale-violating parton distributions of the proton. The PDF’s are calculated using collinear kinematics so that the event structure is the same as given by the lowest-order subprocesses, which are of order $`O(\alpha ^2)`$ <sup>2</sup><sup>2</sup>2Here and in the following we do not count the extra factor $`\alpha `$ from the $`ee\gamma ^{}`$ vertex..
To remove photon production by incoming photons with small virtuality and to restrict to the case where the scattered electron $`e(p_2)`$ is observed, one applies cuts on the usual deep inelastic scattering variables $`x`$, $`y`$, $`Q^2`$ as measured from the momentum of the scattered lepton. In particular we restrict the calculation to values of $`Q^210`$ GeV<sup>2</sup>; however, since very large $`Q^2`$ values are not relevant at HERA we can neglect contributions from $`Z`$ boson exchange. In addition, to have photons $`\gamma (p_5)`$ of sufficiently large energy we require an explicit cut on their transverse momentum. Finally, a cut on the invariant mass of the hadronic final state is also applied.
Both leptons and quarks emit photons. The subset of diagrams where the photon is emitted from the initial or final state lepton (leptonic radiation) is explicitly gauge invariant and can be considered separately. Similarly, the contribution from diagrams with a photon emitted from quark lines is called quarkonic radiation. In addition, there are also contributions from the interference of these two mechanisms. The emission of photons from leptons is described by pure QED and can be predicted with high reliability. Therefore, the contributions from leptonic radiation will be suppressed by cuts on the photon emission angle . In our numerical evaluation we include the remaining background from leptonic radiation as well as the interference contribution.
At lowest order, each parton is identified with a jet and the photon is automatically isolated from the quark jet by requiring a non-zero transverse momentum of the photon or jet in the $`\gamma ^{}p`$ center-of-mass frame. Therefore the photon fragmentation is absent in this order.
### 2.2 Subprocesses to Next-to-Leading Order
At next-to-leading order, processes with an additional gluon, either emitted into the final state or as incoming parton, have to be taken into account:
$$e(p_1)+q(p_3)e(p_2)+q(p_4)+\gamma (p_5)+g(p_6),$$
(2)
$$e(p_1)+g(p_3)e(p_2)+q(p_4)+\gamma (p_5)+\overline{q}(p_6),$$
(3)
where the momenta of the particles are again given in parentheses. Examples of diagrams for $`\gamma ^{}qq\gamma g`$ and $`\gamma ^{}gq\gamma \overline{q}`$ are shown in Fig. 2. In addition, virtual corrections (one-loop diagrams at $`O(\alpha _s)`$) to the LO processes (1) have to be included. The complete matrix elements for the processes (2) and (3) are given in . The processes (2) and (3) contribute both to the $`\gamma +(1+1)`$-jet cross section, as well as to the cross section for $`\gamma +(2+1)`$-jets in the final state. In the latter case each parton in the final state builds a jet on its own, whereas for $`\gamma +(1+1)`$-jets a pair of final state partons is experimentally unresolved. The criteria for combining two partons into one jet will be introduced later. The contributions (2) and (3) as well as the virtual corrections to (1) are of order $`O(\alpha ^2\alpha _s)`$.
### 2.3 Fragmentation Contributions
In addition to the direct production described in the last two subsections, photons can also be produced through the fragmentation of a hadronic jet into a single photon carrying a large fraction of the jet energy . This long-distance process is described in terms of quark-to-photon and gluon-to-photon fragmentation functions which absorb collinear singularities present in the NLO direct contributions of section 2.2. The corresponding fragmentation processes (see Fig. 3) are
$$e(p_1)+q(p_3)e(p_2)+q(p_4)+g(p_6),$$
(4)
$$e(p_1)+g(p_3)e(p_2)+q(p_4)+\overline{q}(p_6).$$
(5)
These processes are of order $`\alpha \alpha _s`$ whereas the photon FF is formally of order $`\alpha `$, so that the LO fragmentation contribution is formally of order $`O(\alpha ^2\alpha _s)`$, i.e. of the same order as the NLO direct contribution. The fragmentation photons, sometimes also called bremsstrahlung photons, are emitted predominantly along the direction of motion of the parent quark or gluon. Because of the pointlike nature of the photon-quark interaction, it is possible to calculate the leading-logarithmic behaviour of the photon FF, including the corrections due to additional gluon emissions. The resulting FF’s are in fact of order $`O(\alpha /\alpha _s)`$ since they possess a logarithmic growth coming from the integration over the momenta of unobserved partons. Therefore in the leading-logarithmic approximation the fragmentation contribution is obtained from $`O(\alpha /\alpha _s)`$ FF’s convoluted with corresponding $`O(\alpha \alpha _s)`$ cross sections for the two-body subprocesses (4) and (5). The resulting contribution is thus of order $`O(\alpha ^2)`$, i.e. of the same order as the LO non-fragmentation process (1). For this reason it is sometimes argued that the fragmentation contribution should be combined with the LO direct process to provide the full LO physical cross section. Consequently, the calculation of the full cross section up to NLO would then also require the computation of the NLO corrections to the fragmentation contributions. On the other hand, it is well-known, that the fragmentation contribution in LO depends strongly on the factorization scale $`\mu _F`$ which, however, is cancelled to a large extent by the $`\mu _F`$-dependence of the NLO contribution to the non-fragmentation part. For this reason and also since for an isolated photon the fragmentation contribution is small we shall take it into account only in LO in the same way as in our previous work .
The signature of the fragmentation contribution in LO is a photon balanced by a jet on the opposite side of the event and accompanied by nearly collinear hadrons on the same side of the event. This means that this contribution has a similar event structure as the LO direct contribution.
### 2.4 Calculational Details
The calculation of the NLO corrections was performed with the help of the phase space slicing method using a slicing parameter defined in terms of invariant masses. With this method it is straightforward to introduce the photon isolation requirement as well as to implement a jet definition which separates $`\gamma +(1+1)`$-jet from $`\gamma +(2+1)`$-jet final states. Phase space slicing based on invariant masses is also used to separate the photon-quark collinearly singular regions, however, using another independent cutoff parameter. The technical steps to apply the phase space slicing method in the present case are described in the following.
In the calculation of the contribution to the $`\gamma +(1+1)`$-jet cross section we encounter the well-known infrared singularities. They appear in those phase space regions where two partons are degenerate to one parton, i.e. when the gluon becomes soft or two partons become collinear to each other. The singularities are assigned either to the initial state (ISR) or to the final state (FSR). Contributions involving the product of an ISR and a FSR factor are separated by partial fractioning. The FSR singularities cancel against singularities from virtual corrections to the LO process (1). For the ISR singularities, the cancellation is incomplete and the remaining singular contributions have to be factorized and absorbed into the renormalized PDF’s of the proton.
To carry out these steps, the singularities are isolated in an analytic calculation using dimensional regularization. Since the corresponding calculations are too difficult for the exact cross sections of the processes (2) and (3) an approximate solution is required. To achieve this, the phase space slicing is used first to separate the singular regions in the 4-particle phase space. Then, in these regions the matrix elements are approximated by their most singular contributions. Only for these approximate expressions and only in the singular regions the calculation is performed analytically. The separation of singular regions is obtained by applying a slicing cut $`y_0^J`$ to the scaled invariant masses $`y_{ij}`$, where $`y_{ij}=(p_i+p_j)^2/W_{\mathrm{had}}^2`$ and the mass of the hadronic final state $`W_{\mathrm{had}}`$ is defined by $`W_{\mathrm{had}}^2=(P+qp_5)^2`$. The cut $`y_0^J`$ must be chosen small enough so that terms of order $`O(y_0^J)`$ which are neglected in the singular approximation are so small that an accuracy of a few per cent is achievable for the final result. Outside the singular regions the integrations are done numerically without any approximation and with 4 space-time dimensions. Physical cross sections, as defined in the next section, are obtained by adding the contributions from singular and non-singular regions as well as the virtual contributions and subtracting the remaining ISR collinear singularities. In the final results, the dependence on the slicing parameter $`y_0^J`$ cancels. This means the cut-off $`y_0^J`$ is purely technical. The independence on the slicing cut $`y_0^J`$ has been checked by explicit calculation for some special photon plus jet cross sections in . Further details and the derivation of the two-body matrix elements in the singular region together with the cancellation of the soft and collinear poles can be found in .
In addition, the squared matrix elements for the processes (2) and (3) have photonic infrared and collinear singularities, i.e. singularities due to soft or collinear photons. Since we require the photon to be observed in the detector the infrared singularity can not occur. In the numerical analysis we will introduce this condition by requiring a minimum on the transverse momentum of the photon. This cut removes also all collinear singularities due to initial state radiation.
Final state collinear singularities due to photons are present and are treated again with the help of the phase space slicing method in a similar way as the quark-gluon collinear contributions. The phase space slicing parameter used to treat the photonic singularities can be chosen independently and is denoted by $`y_0^\gamma `$. As before, it has to be chosen very small so that the matrix element can be approximated by its singular part. For the subprocess (2), the phase space slicing is described by the squared invariant masses $`y_{45}=(p_4+p_5)^2/W_{\mathrm{had}}^2`$ . In the gluon-initiated process (3) one has two singular regions which are controlled by the variables $`y_{45}`$ and $`y_{56}`$, respectively. In the regions $`y_{45}`$, $`y_{56}>y_0^\gamma `$ the cross section is evaluated numerically in the same way as in where these cuts were introduced as physical isolation cuts on the photon with sufficiently large isolation parameter $`y_0^\gamma `$. In this work the cuts on $`y_{45}`$ and $`y_{56}`$ are only technical since we include the contribution to the matrix element also in the regions $`y_{45}<y_0^\gamma `$ and $`y_{56}<y_0^\gamma `$. In these regions the matrix elements are collinearly divergent. The singularities are regulated by dimensional regularization, allowing us to absorb their divergent parts into the bare photon FF to yield the renormalized photon FF denoted by $`D_{q\gamma }`$. For the process (2) this procedure results in a contribution of the following form:
$$|M|_{\gamma ^{}q\gamma qg}^2=|M|_{\gamma ^{}qqg}^2D_{q\gamma }(z).$$
(6)
The matrix element $`|M|^2`$ on the right-hand side of (6) is the matrix element for the process $`\gamma ^{}qqg`$ whose Feynman diagram is shown in Fig. 3a. There exists a similar expression for the subprocess $`\gamma ^{}g\gamma q\overline{q}`$ (Fig. 3b). The photon FF $`D_{q\gamma }(z)`$ in (6) is given by
$$D_{q\gamma }(z)=D_{q\gamma }(z,\mu _F^2)+\frac{\alpha e_q^2}{2\pi }\left(P_{q\gamma }^{(0)}(z)\mathrm{ln}\frac{z(1z)y_0^\gamma W^2}{\mu _F^2}+z\right).$$
(7)
$`D_{q\gamma }(z,\mu _F^2)`$ in (7) stands for the non-perturbative FF describing the transition $`q\gamma `$ at the factorization scale $`\mu _F`$. This function will be specified in the next section. The second term in (7), if substituted in (6), is the finite part due to the collinear photon-quark (-antiquark) contribution to the matrix element $`|M|_{\gamma ^{}q\gamma qg}^2`$ integrated in the region $`y_{45}<y_0^\gamma `$ after absorption of the divergent part into the non-perturbative FF.
Again the parameter $`y_0^\gamma `$ is only a parameter used in intermediate steps of the calculation, introduced to separate divergent from finite contributions; the $`y_0^\gamma `$-dependence in (7) is canceled by the dependence of the numerically computed $`\gamma +(1+1)`$-jet cross section restricted to the region $`y_{45}>y_0^\gamma `$. Since the corresponding contributions to the matrix element in (6) are calculated in the collinear approximation, the result is valid only up to terms of order $`O(y_0^\gamma )`$. This requires to choose a very small value for $`y_0^\gamma `$. In Ref. it has been explicitly shown that the sum of all terms for the photon plus jet cross section becomes independent of $`y_0^\gamma `$ when $`y_0^\gamma `$ is chosen small enough.
In (7), $`P_{q\gamma }^{(0)}`$ is the LO quark-to-photon splitting function
$$P_{q\gamma }^{(0)}(z)=\frac{1+(1z)^2}{z}$$
(8)
and $`e_q`$ is the electric charge of quark $`q`$. The variable $`z`$ denotes the fraction of the quark momentum carried away by the photon. If the photon is emitted from the final state quark with 4-momentum $`p_4^{}=p_4+p_5`$, then $`z`$ can be related to the invariants $`y_{35}`$ and $`y_{34}`$:
$$z=\frac{y_{35}}{y_{34^{}}}=\frac{y_{35}}{y_{34}+y_{35}}.$$
(9)
The fragmentation contribution is proportional to the cross section for $`\gamma ^{}qqg`$ which is $`O(\alpha \alpha _s)`$ and is well-known. It must be convoluted with the function in (7) to obtain the contribution to the cross section for $`\gamma ^{}qq\gamma g`$ at $`O(\alpha ^2\alpha _s)`$. Equivalent formulas are used to calculate the fragmentation contributions to the channel (3) and in the case where the quarks in the initial and final state are replaced by an antiquark in (2).
## 3 Results
### 3.1 Kinematical Selection Cuts and Other Input
The results for the cross sections which will be presented in the following subsections are obtained for energies and kinematical cuts appropriate for HERA experiments. The energies of the incoming electron (positron) and proton are $`E_e=27.5`$ GeV and $`E_P=820`$ GeV, respectively. The cuts on the DIS variables are chosen as in our previous works :
$$\begin{array}{cc}Q^210\mathrm{GeV}^2,\hfill & W_{\mathrm{had}}10\mathrm{GeV},\hfill \\ 10^4x0.5,\hfill & 0.05y0.95\hfill \end{array}$$
where $`Q^2=q^2`$ and $`q`$ is the electron momentum transfer, $`q=p_1p_2`$ as usual.
To reduce the background from leptonic radiation we require
$$90^{}\theta _\gamma 173^{},\theta _{\gamma e}10^{}$$
(10)
where $`\theta _\gamma `$ is the emission angle of the photon measured with respect to the momentum of the incoming electron in the HERA laboratory frame. The cut on $`\theta _{\gamma e}(=\theta _{25})`$, the angle between the momenta of the photon and the scattered electron, suppresses radiation from the final state electron. In the case of photon plus jet cross sections, the photon and the hadron jets $`J`$ are required to have minimal transverse momenta in the $`\gamma ^{}p`$ center-of-mass system (i.e. in the rest system of $`p_1p_2+P`$ where the remnant has $`p_{T,r}=0`$),
$$p_{T,\gamma }5\mathrm{GeV},p_{T,J}6\mathrm{GeV}.$$
(11)
Note that an event is rejected if we do not find at least one jet with $`p_{T,J}`$ above the cut in (11). If there are two partons with $`p_{T,J}6`$ GeV, the event has a chance to be treated as a $`\gamma +(2+1)`$-jet event. If, after trying to recombine partons into jets (see below), only one jet has $`p_{T,J}6`$ GeV, the event contributes to the $`\gamma +(1+1)`$-jet class. In the latter case, it may happen that there is an additional parton not combined into a jet with $`p_{T,i}<6`$ GeV (in a JADE-like jet algorithm applied in the HERA laboratory frame, such a low-$`p_T`$ parton would be recombined with the remnant jet in most cases).
Different values of minimal $`p_T`$’s for the photon and the jet have to be chosen in order to avoid the otherwise present infrared sensitivity of the NLO predictions . This point will be studied in detail later. For inclusive photon cross sections, i.e. in the case where we do not perform a jet analysis of the hadronic final state, we replace the second of the conditions in (11) by a cut on the sum of transverse energies of all final state partons
$$E_T=\underset{i=q,\overline{q},g}{}\left|p_{T,i}\right|6\mathrm{GeV}.$$
(12)
The PDF’s of the proton are taken from (MRST). Recent updates of available PDF parametrizations (MRS99, CTEQ5) do not lead to markedly different results (roughly a 2 % (4 %) increase of the total cross section for MRS99 (CTEQ5) with respect to MRST; a version with enhanced $`d/u`$ ratio of MRS99 does not lead to observable differences in the range of $`x`$ and $`Q^2`$ considered here). $`\alpha _s`$ is calculated from the two-loop formula with the same $`\mathrm{\Lambda }`$-value ($`\mathrm{\Lambda }_{\overline{MS}}(n_f=4)=300`$ MeV) as used in the MRST parametrization of the proton PDF. The scale in $`\alpha _s`$ and the factorization scale are set equal to each other and fixed to the largest $`p_{T,J}`$, except when we present results with other scale choices. For completeness, we also mention that the slicing cuts have been fixed at $`y_0^J=10^4`$ and $`y_0^\gamma =10^5`$.
The dependence of the $`\gamma +`$jet cross sections on the choice of the photon FF has been studied earlier . In the present work we choose the FF of Bourhis et al. (BFG). This FF has been compared to the ALEPH $`e^+e^{}\gamma +1`$-jet cross section and also to the inclusive photon cross section measured by OPAL . Both data sets agreed well with predictions based on the BFG parametrization . In we studied the cross section differential with respect to the fraction of momentum $`z_\gamma `$ carried by a photon inside a jet for several other photon FF’s besides the BFG parametrization.
Photon isolation is implemented with the help of the cone isolation method similar to the one used in the ZEUS experiment for photon production with almost real photons . This method restricts the hadronic energy allowed in a cone around the jet axis of the jet containing the photon. The same method is used also to define jets emerging in the event sample of $`\gamma +2`$-parton-level jets when two partons are combined. In the $`\gamma ^{}p`$ center-of-mass frame, two partons $`i`$ and $`j`$ are combined into a jet $`J`$, when they obey the cone constraints $`R_{i,J}<R`$ and $`R_{j,J}<R`$, where
$$R_{i,J}=\sqrt{(\eta _i\eta _J)^2+(\varphi _i\varphi _J)^2}.$$
(13)
$`\eta _J`$ ($`=\mathrm{ln}\mathrm{tan}(\theta _J/2)`$) and $`\varphi _J`$ are the rapidity (polar angle $`\theta _J`$) and azimuthal angle of the recombined jet which are obtained from the formulae
$$\begin{array}{c}p_{T,J}=p_{T,i}+p_{T,j},\hfill \\ \eta _J=\frac{\eta _ip_{T,i}+\eta _jp_{T,j}}{p_{T,i}+p_{T,j}},\hfill \\ \varphi _J=\frac{\varphi _ip_{T,i}+\varphi _jp_{T,j}}{p_{T,i}+p_{T,j}}.\hfill \end{array}$$
(14)
For most of the results we choose $`R=1`$. If not, we shall state the value of $`R`$ explicitly. The azimuthal angle is defined with respect to the scattering plane given by the momentum of the beam and the momentum of the scattered electron. The rapidity is always defined positive in the direction of the proton remnant. The photon is treated like any other parton in the recombination process, i.e. when in (13) $`i`$ or $`j`$ is the photon, then $`J`$ is called the photon-jet. To qualify a jet as a photon-jet, we restrict the hadronic energy in the jet by
$$z_\gamma =\frac{p_{T,\gamma }}{p_{T,\gamma }+p_{T,\mathrm{had}}}=1ϵ_{\mathrm{had}}1ϵ_{\mathrm{had}}^0=z_{\mathrm{cut}}.$$
(15)
$`p_{T,\gamma }`$ and $`p_{T,\mathrm{had}}`$ denote the transverse momenta of the photon and the parton producing hadrons in this jet. Defining $`p_{T,\gamma \mathrm{jet}}=p_{T,\gamma }+p_{T,\mathrm{had}}`$, (15) is equivalent to the requirement
$$p_{T,\mathrm{had}}ϵ_{\mathrm{had}}^0p_{T,\gamma \mathrm{jet}},$$
(16)
i.e. the ratio of the total $`p_T`$ due to other particles (partons) than the photon is required not to exceed $`ϵ_{\mathrm{had}}^0`$. For our predictions we shall choose different values for $`ϵ_{\mathrm{had}}^0=(1z_{\mathrm{cut}})`$. In this parameter was set equal to $`ϵ_{\mathrm{had}}^0=0.11`$.
It is known that the cone algorithm is ambiguous for final states with more than three particles or partons . Since we have maximally three partons in the final state this problem is not relevant in our case. However, in some cases it may happen that two partons $`i`$ and $`j`$ qualify both as two individual jets $`i`$ and $`j`$, or as a combined jet $`J`$. In these exceptional cases we count only the combined jet $`J`$ to avoid double counting.
The cone algorithm is problematic in experimental analyses due to its seed-finding mechanism and due to overlapping cones. These problems are avoided with the $`k_T`$ jet finding algorithm . In our theoretical calculations, the $`k_T`$ algorithm can be incorporated quite easily: partons $`i`$ and $`j`$ (where one of them may be the photon) are combined if they fulfill the condition
$$R_{ij}<R\mathrm{with}R_{ij}=\sqrt{(\eta _i\eta _j)^2+(\varphi _i\varphi _j)^2}.$$
(17)
The resulting kinematic variables of the combined jet are calculated with the same formulae (14) as in the cone algorithm. The recombination condition (13) is equivalent to
$$R_{ij}<Rp_{ij}\mathrm{with}p_{ij}=\frac{p_{T,i}+p_{T,j}}{\mathrm{max}(p_{T,i},p_{T,j})}.$$
(18)
Therefore, choosing the same value for $`R`$, jets obtained with the cone algorithm (18) are slightly wider than those constructed with the $`k_T`$ algorithm (17). In order to demonstrate that our numerical routines work also with the $`k_T`$ algorithm we have calculated some representative cross sections with this jet definition as well.
### 3.2 Photon plus Jet Cross Sections
Now we shall present our numerical results. We start with various cross sections for the $`\gamma +(n+1)`$-jet final state since we think that these cross sections, although smaller than the fully inclusive photon cross section, will be measured first due to reduced background problems. It is clear that in NLO the final state may consist of two or three jets where one jet is always a photon jet. The remnant jet is not counted since it is produced with zero transverse momentum. The three-jet sample, equivalent to $`\gamma +(2+1)`$-jets in the notation of the previous sections, consists of all $`\gamma +(2+1)`$-parton-level jets which do not fulfill the cone constraint $`R_{i,J}<R`$ with $`R_{i,J}`$ given in (13). The $`\gamma +(1+1)`$-jet sample contains events where two partons (possibly a photon) are recombined into one jet or one parton does not obey the cut on transverse momenta (11). In the following we shall sum over the two samples with $`n=1`$ and $`n=2`$ jets. If there are two jets in an event, we order them according to their transverse momenta and call the one with larger $`p_T`$ “jet 1” and the one with smaller $`p_T`$ accordingly “jet 2”. Also, from now on we will use a simplified notation and denote kinematic variables of the jet containing the photon simply by $`p_{T,\gamma }`$, $`\eta _\gamma `$ and $`\varphi _\gamma `$.
In Figs. 4 and 5 we show results for the $`p_T`$ and $`\eta `$ dependence of the cross sections $`d\sigma /dp_T`$ and $`d\sigma /d\eta `$ for the photon jet and the jet with the largest $`p_T`$. In each of the four figures we have plotted three curves for three choices of $`z_{\mathrm{cut}}`$ defined in (15), $`z_{\mathrm{cut}}=0.5`$, 0.7 and 0.9. Together with the cone radius $`R`$, $`z_{\mathrm{cut}}`$ controls the amount of photon isolation. As to be expected the cross section decreases with the degree of photon isolation, i.e. with increasing $`z_{\mathrm{cut}}`$. Specifically in Fig. 4a we present $`d\sigma /dp_{T,\gamma }`$ as a function of $`p_{T,\gamma }`$ for the three $`z_{\mathrm{cut}}`$ values and for $`p_{T,\gamma }5`$ GeV. All other variables, in particular $`\eta _{\mathrm{jet}}`$, $`\eta _\gamma `$ and $`p_{T,\mathrm{jet}}`$ are integrated over the kinematically allowed ranges. We see that all three cross sections have a similar shape. In Fig. 5a $`d\sigma /dp_{T,\mathrm{jet1}}`$ as a function of $`p_{T,\mathrm{jet1}}`$ is shown. The qualitative behaviour of the cross sections for the different $`z_{\mathrm{cut}}`$’s is similar as in Fig. 4a.
For the $`\eta `$ distributions, shown in Figs. 4b, 5b for the photon jet and the most energetic jet, we have integrated over $`p_{T,\gamma }5`$ GeV and $`p_{T,\mathrm{jet1}}6`$ GeV. The shapes of the three curves for the different $`z_{\mathrm{cut}}`$ values are similar for both the cross section $`d\sigma /d\eta _\gamma `$ in Fig. 4b and the cross section $`d\sigma /d\eta _{\mathrm{jet1}}`$ in Fig. 5b. We note that the $`\eta _{\mathrm{jet1}}`$ distribution peaks at somewhat smaller rapidities than $`d\sigma /d\eta _\gamma `$.
Distributions with respect to $`p_T`$ and $`\eta `$ of the second jet in $`\gamma +(2+1)`$-jet events do not depend on the isolation cut since in this case each parton (photon) constitutes a jet by its own. Therefore we show corresponding figures in the next subsection when we will discuss the influence of the cone size $`R`$ on the cross sections.
In addition to predicting distributions in the transverse momenta and the rapidities of the photon and the hadron jets we also have calculated distributions for variables which characterize the correlation of two jets. One of these variables is
$$z_{\gamma 1}=\frac{\stackrel{}{p}_{T,\gamma }\stackrel{}{p}_{T,\mathrm{jet1}}}{p_{T,\mathrm{jet1}}^2}.$$
(19)
Note that $`z_{\gamma 1}`$ is defined with the help of the transverse momentum of the photon jet, i.e. $`p_{T,\gamma }`$ may include a contribution from accompanying hadronic energy. The dependence of the cross section on $`z_{\gamma 1}`$ characterizes the imbalance in transverse momentum of the photon and the most energetic jet. Similar variables have been used before for studies in the case of photon plus charm jet final states in $`p\overline{p}`$ collisions and of two-jet production in $`ep`$ scattering . The result for $`d\sigma /dz_{\gamma 1}`$ is shown in Fig. 6a for three photon isolation cuts $`z_{\mathrm{cut}}=0.5`$, 0.7 and 0.9. The cuts on transverse momenta and the cone parameters are as defined before.
For two-body processes such as the LO Compton subprocess $`\gamma ^{}q\gamma q`$, the final photon and the jet have balancing transverse momenta and the distribution is a $`\delta `$-function in $`(1z_{\gamma 1})`$. Also the fragmentation process contributes only to $`z_{\gamma 1}=1`$ since in this case the transverse momentum of accompanying hadronic energy is collinear with the photon, resulting in $`p_{T,\gamma }=p_{T,\mathrm{bare}\gamma }+p_{T,\mathrm{had}}=p_{T,\mathrm{jet1}}`$. Contributions with $`z_{\gamma 1}1`$ are due to the higher-order three-body contributions. Events with $`z_{\gamma 1}<1`$ typically result from configurations where a single photon is opposite in transverse momentum to a jet consisting of two partons. Since according to the recombination prescription (14) the scalar sum of transverse momenta is ascribed to the jet, not the vectorial sum, one finds always $`p_{T,\gamma }<p_{T,\mathrm{jet1}}`$ and thus $`z_{\gamma 1}<1`$ in this case. Moreover, the photon is never accompanied by a hadronic parton in this case and events with $`z_{\gamma 1}<1`$ are consequently not affected by the isolation cut. On the other hand, events with $`z_{\gamma 1}>1`$ are predominantly due to configurations with a photon-jet consisting of a photon and a quark or gluon opposite to a jet consisting of a single parton. In this case, a variation of the photon isolation cut $`z_{\mathrm{cut}}`$ has a strong effect on the differential cross section.
These features are clearly visible in Fig. 6a. The cross section increases when lowering $`z_{\mathrm{cut}}`$, i.e. when larger fragmentation contributions are included. We notice that the $`z_{\gamma 1}`$ distribution is not symmetric around $`z_{\gamma 1}=1`$. The cross section for $`z_{\gamma 1}<1`$ is larger than for $`z_{\gamma 1}>1`$, becoming more and more symmetric for less restrictive isolation cuts on $`z_\gamma `$. The residual asymmetric behaviour of this distribution for vanishing isolation cut is a dynamical property of the underlying cross section.
In the region of $`z_{\gamma 1}`$ near unity, one of the two final state partons of three-body contributions (not the photon) becomes soft and thus this region is sensitive to soft-gluon effects. In our calculation with an invariant mass cut slicing parameter $`y_0^J`$ these soft-gluon corrections are considered as two-body contributions. Their contribution depends on the slicing parameter $`y_0^J`$. To remove this dependence, i.e. to remove the infrared sensitivity, we must include a sufficiently large fraction of the three-body contribution from outside $`z_{\gamma 1}=1`$. Therefore we averaged the $`z_{\gamma 1}`$ distribution over sufficiently large bins and studied the sum of the $`\gamma +(1+1)`$\- and $`\gamma +(2+1)`$-jet cross sections $`d\sigma /dz_{\gamma 1}`$. We have chosen a bin width of $`\mathrm{\Delta }z_{\gamma 1}=0.2`$ around $`z_{\gamma 1}=1`$ and $`\mathrm{\Delta }z_{\gamma 1}=0.1`$ elsewhere.
It is clear that the cross section outside the bin at $`z_{\gamma 1}=1`$ has a stronger scale dependence than inside this bin since only three-parton terms contribute. The cross section inside the bin at $`z_{\gamma 1}=1`$ is a genuine NLO prediction with expected reduced scale dependence. The scale dependence will be studied later for some other distributions. The $`\delta `$-function behaviour at LO is of course in reality modified not only by NLO corrections, but also by non-perturbative effects originating from hadronization and a possible intrinsic transverse momentum of the initial parton. Our calculation does not include these latter effects.
In Fig. 6b we show the cross section $`d\sigma /dz_{12}`$ where the variable
$$z_{12}=\frac{\stackrel{}{p}_{T,\mathrm{jet1}}\stackrel{}{p}_{T,\mathrm{jet2}}}{p_{T,\mathrm{jet1}}^2}$$
(20)
with $`p_{T,\mathrm{jet1}}>p_{T,\mathrm{jet2}}`$ measures the correlation between the two jets in $`\gamma +(2+1)`$-jet events. This cross section peaks at $`z_{12}=0`$ as to be expected and decreases away from $`z_{12}=0`$. The point $`z_{12}=0`$ is the point with no second jet, i.e. the pure $`\gamma +(1+1)`$-jet region. This region is again infrared sensitive. Therefore we integrated here over the bin $`0.05<z_{12}<0.05`$. Outside this region we chose a bin size of $`\mathrm{\Delta }z_{12}=0.05`$. Note that the distribution shown in Fig. 6b includes the contribution from low-$`p_T`$ partons with $`p_T`$ below the cut in (11). We have also calculated the $`z_{12}`$-distribution restricting to $`\gamma +(2+1)`$-jet events where both jets have $`p_T6`$ GeV. In this case, the distribution would extend to larger values of $`z_{12}`$ with a maximum at $`z_{12}0.6`$. The asymmetric behaviour of the curve visible in Fig. 6b, i.e. the tail at large $`z_{12}`$, is due to the contribution from these $`\gamma +(2+1)`$-jet events.
Another interesting variable might be the azimuthal angle $`\varphi _\gamma `$ of the emitted photon. We define $`\varphi _\gamma `$ with respect to the plane spanned by the momenta of the beam and of the scattered lepton. In Fig. 7 we show the dependence of the cross section on $`\varphi _\gamma `$, again for the three $`z_{\mathrm{cut}}`$ values 0.5, 0.7 and 0.9. As is seen in this figure, the photon is emitted dominantly at $`\varphi _\gamma =0`$. We note that the distribution becomes flatter with decreasing $`z_{\mathrm{cut}}`$. It is symmetric within the statistical accuracy of the calculation. We do not present here a similar plot for $`d\sigma /d\varphi _{\mathrm{jet1}}`$, the cross section with respect of the azimuthal angle of the most energetic jet. It would show a distribution which peaks at $`\varphi _{\mathrm{jet1}}=\pi `$, since the dominant contribution to the cross section comes from configurations in which the photon and the jet with the largest $`p_T`$ are emitted back-to-back. For jet 2 there is no such correlation. The cross section $`d\sigma /d\varphi _{\mathrm{jet2}}`$ is independent of $`\varphi _{\mathrm{jet2}}`$ and, in NLO, does not change with $`z_{\mathrm{cut}}`$. In a similar way one can discuss the cross section as a function of $`\varphi _{\gamma 1}=\varphi _\gamma \varphi _{\mathrm{jet1}}`$ or $`\varphi _{\gamma 2}=\varphi _\gamma \varphi _{\mathrm{jet2}}`$. $`d\sigma /d\varphi _{\gamma 1}`$ is strongly peaking at $`\varphi _{\gamma 1}=\pi `$. Here it is again necessary to calculate $`d\sigma /d\varphi _{\gamma 1}`$ with a sufficiently large bin size around $`\varphi _{\gamma 1}=\pi `$ in order to avoid any infrared sensitivity. On the other hand, the distribution $`d\sigma /d\varphi _{\gamma 2}`$ is flat around $`\varphi _{\gamma 2}=0`$ and decreases towards $`\varphi _{\gamma 2}=\pi `$.
### 3.3 Cone Size Dependence of Jet Cross Sections
So far we presented results only for the cone jet algorithm with the cone radius fixed to $`R=1`$ for both the photon jet and purely hadronic jets. Sometimes it is advantageous to use smaller cone radii to suppress background processes. On the other hand the dependence of the cross sections on the cone radius is a genuine NLO effect since LO cross sections do not depend on the jet definition. To present an overview of the cone radius dependence we have recalculated some of the cross sections shown so far for $`R=1`$ for two smaller radii $`R=0.5`$ and 0.7. For the photon isolation parameter we fix now $`z_{\mathrm{cut}}=0.9`$.
In Figs. 8a, b we show the results for $`d\sigma /dp_{T,\gamma }`$ and $`d\sigma /d\eta _\gamma `$ with $`R=0.5`$, 0.7, and 1.0. All other cuts are chosen as before. The case with $`R=1`$ and $`z_{\mathrm{cut}}=0.9`$ was shown in Figs. 4a, b already. The distributions for the jet with largest $`p_T`$ are exhibited in Figs. 9a and b. $`d\sigma /dp_{T,\gamma }`$, $`d\sigma /d\eta _\gamma `$, $`d\sigma /dp_{T,\mathrm{jet1}}`$ and $`d\sigma /d\eta _{\mathrm{jet1}}`$ show very little dependence on the cone size $`R`$. For a less restrictive isolation cut of $`z_\gamma 0.5`$ these distributions would decrease with decreasing cone radius $`R`$ as it is known from other jet calculations.
In our previous work we studied the equivalent cross sections also for the two event classes with $`\gamma +(1+1)`$-jets and $`\gamma +(2+1)`$-jets separately. It turned out that the contribution for $`\gamma +(1+1)`$-jets is the dominant one. This is expected, since the contribution with $`\gamma +(2+1)`$-jets is an $`O(\alpha _S)`$ effect. Here we present now also results for the $`p_T`$ and $`\eta `$ distributions of jet 2 in this latter process. The results are shown in Fig. 10a ($`d\sigma /dp_{T,\mathrm{jet2}}`$) and Fig. 10b ($`d\sigma /d\eta _{\mathrm{jet2}}`$). For $`p_{T,\mathrm{jet2}}6`$ GeV the cross section is very much reduced as compared to the cross section in Fig. 9a. In fact, the dominating event configuration is with a photon and one jet balancing each other in $`p_T`$; a third jet with comparable $`p_T`$ is found in only a small portion of the events. The rapidity distribution $`d\sigma /d\eta _{\mathrm{jet2}}`$ plotted in Fig. 10b has a larger tail extending to larger rapidities as compared to the cross sections $`d\sigma /d\eta _\gamma `$ and $`d\sigma /d\eta _{\mathrm{jet1}}`$. The second energetic jet originates dominantly from the incoming quark and therefore is in many cases closer to the proton remnant, i.e. at positive rapidities, than the harder jet.
The cross sections $`d\sigma /dp_{T,\mathrm{jet2}}`$ and $`d\sigma /d\eta _{\mathrm{jet2}}`$ increase with decreasing $`R`$. This is the behaviour expected for cross sections which are of leading order in $`\alpha _s`$.
In our previous work we studied $`d\sigma /dz_\gamma `$ as a function of the fraction $`z_\gamma `$ of the momentum of the photon inside the photon jet. This cross section is expected to contain information on the photon fragmentation functions. The results presented in were obtained for the case $`R=1`$ only. To see how results change with $`R`$, we show in Fig. 11 $`d\sigma /z_\gamma `$ for the particular photon fragmentation function of Bourhis et al. for $`R=0.5`$, 0.7 and 1.0. As is seen, $`d\sigma /dz_\gamma `$ decreases with decreasing $`R`$ since this cross section is a superposition of leading and next-to-leading order contributions.
The $`R`$ dependence of the other cross sections considered above follows the same pattern. In cases where we have a superposition of LO and NLO contributions we encounter a decreasing cross section with decreasing $`R`$. In regions where the cross section receives contribution from $`O(\alpha _s)`$ only with no additional NLO corrections included, the cross section increases with decreasing $`R`$. Thus for example $`d\sigma /dz_{\gamma 1}`$ (cf. Fig. 6a) is decreasing with $`R`$ in the bin at $`z_{\gamma 1}=1`$ and increasing with $`R`$ outside this bin, i.e. for $`z_{\gamma 1}<0.9`$ and $`z_{\gamma 1}>1.1`$. The $`R`$ dependence of $`d\sigma /dz_{12}`$ (cf. Fig. 6b) is similar: inside the bin around $`z_{12}=0`$ the cross section decreases with decreasing $`R`$ and outside $`z_{12}=0`$ it increases with decreasing $`R`$.
### 3.4 Jet Algorithm Dependence
Measurements of cross sections for the production of a photon plus jets, as for example in $`\gamma p`$ collisions , have been performed with the help of the cone algorithm used to define jets and the photon isolation. However, the $`k_T`$ algorithm has definite advantages, in particular in the experimental analysis . Therefore we present a few distributions based on the $`k_T`$ algorithm as well. This algorithm is used for the recombination of two partons into a jet as explained in sect. 3.1 as well as for the definition of the photon jet. In Figs. 12a, 13a we show $`d\sigma /dp_{T,\gamma }`$ and $`d\sigma /dp_{T,\mathrm{jet1}}`$ calculated with the $`k_T`$ algorithm for $`R=0.5`$, 0.7 and 1. These two groups of curves have to be compared to the results with the cone algorithm in Figs. 8a, 9a. The qualitative behaviour of the two cross sections is similar; but there are quantitative differences. We note that the dependence on the parameter $`R`$, which controls the size of the jets, is now even more reduced as compared to the corresponding cross sections with the cone algorithm. The cross sections $`d\sigma /d\eta _\gamma `$ and $`d\sigma /d\eta _{\mathrm{jet1}}`$ for the $`k_T`$ algorithm are displayed in Figs. 12b, 13b. They can be compared with the corresponding cross section for the cone algorithm in Figs. 8b, 9b. Similar to the $`p_T`$ distributions the qualitative behaviour did not change. Again the $`R`$ dependence seems to be reduced for the $`k_T`$ algorithm. This is even more the case for the cross sections $`d\sigma /dp_{T,\mathrm{jet2}}`$ and $`d\sigma /d\eta _{\mathrm{jet2}}`$, shown in Figs. 14a, b which should be compared to the corresponding results in Figs. 10a, b. It is clear that the cross sections in Figs. 14a, b increase with decreasing $`R`$ in the same way as the cross sections for the cone algorithm in Figs. 10a, b.
A direct comparison of cross sections calculated with either the cone algorithm or the $`k_T`$ algorithm is shown in Fig. 15 (dashed and dotted curves). Here we have chosen $`R=1`$ and $`z_{\mathrm{cut}}=0.9`$. In Figs. 15a, b the cross sections $`d\sigma /dp_{T,\gamma }`$ and $`d\sigma /d\eta _\gamma `$ are plotted, respectively. As we can see, these cross sections hardly change when the cone algorithm is replaced by the $`k_T`$ algorithm. Only where $`d\sigma /d\eta _\gamma `$ is maximal, i.e. near $`\eta _\gamma =1.5`$, the cross section with the $`k_T`$ algorithm is approximately 5 % larger than with the cone algorithm.
It is expected that the requirement to observe additional jets reduces the cross section for the production of a high-$`p_T`$ photon. To study this reduction we have calculated also $`d\sigma /dp_{T,\gamma }`$ and $`d\sigma /d\eta _\gamma `$ for the inclusive case, i.e. without additional jets required in the final state and the same photon isolation constraint in terms of an isolation cone around the photon and the cut $`z_{\mathrm{cut}}=0.9`$. The results for these inclusive cross sections are shown in Figs. 15a, b as full curves. The total cross section without jet algorithm applied is increased by about 35 %. This is reflected in the inclusive cross section $`d\sigma /d\eta _\gamma `$ which is always larger than the corresponding distribution for final states which are required to contain at least one jet, over the full range of $`\eta _\gamma `$. In $`d\sigma /dp_{T,\gamma }`$ only the first two $`p_T`$-bins contain an appreciably larger cross section due to the removal of the jet requirement. This can be traced back to the different prescriptions used to remove events with small transverse momentum: for the analyses based on the cone or the $`k_T`$ jet algorithms, we used the $`p_T`$ cuts (11) which are applied to individual jets, whereas in the case without jet algorithm we applied the cut (12) to the sum of transverse momenta of all hadronic particles in the final state. As a consequence, an event with two low-$`p_T`$ partons, each with $`p_T<6`$ GeV will be rejected in the first case if these two partons are not recombined into one jet, while the event will be accepted in the second case if the sum of the transverse momenta of the two partons is larger than the cut of 6 GeV. This clearly affects only the bins with lowest $`p_T`$.
### 3.5 Scale Dependence of Jet Cross Sections
All results presented so far have been obtained for a renormalization and factorization scale $`\mu `$ fixed at $`\mu =p_{T,\mathrm{jet1}}`$ ($`\mu =E_T`$ of (12) for inclusive cross sections). In LO cross sections, the scale dependence is exclusively due to variations of $`\mu `$ in the parton distribution functions. At NLO we expect that additional terms containing an explicit $`\mu `$-dependence will reduce the scale dependence. Instead of studying the scale dependence of all the differential cross sections discussed so far separately, we have investigated the scale dependence of some components of the total cross section, i.e. integrated over the phase space allowed by the transverse momentum cuts (11). We define the scale in the form $`\mu ^2=f^2p_{T,\mathrm{jet1}}^2`$ and vary $`f`$ between $`f=1/4`$ and 4.
In Fig. 16a we have plotted the $`f`$ dependence of the $`\gamma +(1+1)`$-jet cross section in LO and NLO (denoted $`O(\alpha _s)`$) and of the $`\gamma +(2+1)`$-jet cross section. The LO cross section (denoted “Born” in Fig. 16) increases with $`f`$ by approximately 10 % in the range $`0.25<f<4`$. The cross section including corrections of $`O(\alpha _s)`$ is almost independent of $`f`$, i.e. no scale dependence inside the considered range of scales is visible. Here the decrease of the cross section due to the decrease of $`\alpha _s`$ with increasing $`f`$ is compensated by the increase originating from the scale dependence of the proton PDF’s. The cross section for the $`\gamma +(2+1)`$-jet final state, being of order $`O(\alpha _s)`$, decreases with increasing $`f`$ by approximately 25 %. This is a combined effect of the dependence of $`\alpha _s`$ and of the $`f`$ dependence of the parton distribution functions. The $`f`$ dependence of separate components (“real”, “sing” and “frag”) of the $`\gamma +(1+1)`$-jet cross section including $`O(\alpha _s)`$ terms is plotted in Fig. 16b and again compared with the LO cross section. “real” stands for the tree graph contributions in $`O(\alpha _s)`$, calculated with the slicing parameter as described in section 2.4. Since it is a tree-graph term it decreases with increasing $`f`$ due to the decrease of $`\alpha _s`$. The singular term “sing”, which includes virtual corrections and singular contributions below the slicing cut, is negative and decreases in absolute value by the same amount as “real”. The fragmentation contribution “frag” which includes both terms of the right-hand side of (7) is also negative and almost independent of $`f`$ as expected, since the factorization scale dependence ($`\mu _F`$ in (7)) cancels in first approximation.
### 3.6 Dependence on Low-$`p_T`$ Cuts
The choice of two different cuts for the transverse momentum of the photon and the jet is needed to avoid the otherwise present infrared sensitivity of the NLO predictions. This sensitivity is known from similar calculations of dijet cross sections in $`ep`$ collisions and must be avoided. The same problem was encountered in the calculation of inclusive two-jet cross sections in $`\gamma p`$ collisions , for the production of a prompt photon plus a charm quark in $`p\overline{p}`$ collisions and much earlier in NLO calculations of the inclusive cross section for photon-hadron and for two-photon production .
Above we have chosen the difference between the two $`p_{T,\mathrm{min}}`$ cuts for the photon and the jet, $`\mathrm{\Delta }=p_{T,J}^{\mathrm{min}}p_{T,\gamma }^{\mathrm{min}}`$, equal to 1 GeV (see (11)). $`\mathrm{\Delta }`$ should not be too small since then we would encounter the infrared sensitive region where the prediction of the cross section becomes unreliable. In order to obtain some information about possible choices for $`\mathrm{\Delta }`$ we have studied the $`\gamma +(n+1)`$-jet cross section $`d\sigma /dp_{T,J}`$ integrated over $`p_{T,J}p_{T,J}^{\mathrm{min}}`$ as a function of $`p_{T,J}^{\mathrm{min}}`$. The transverse momentum of the photon was always integrated over the range $`p_{T,\gamma }p_{T,\gamma }^{\mathrm{min}}=5`$ GeV. The results for $`\sigma (p_{T,J}^{\mathrm{min}})`$ are plotted in Fig. 17a. Starting at $`p_{T,J}^{\mathrm{min}}=6`$ GeV this cross section increases with decreasing $`p_{T,J}^{\mathrm{min}}`$ with almost constant slope. At about $`p_{T,J}^{\mathrm{min}}=5.5`$ GeV the slope decreases and approaches zero and even changes sign so that $`\sigma (p_{T,J}^{\mathrm{min}})`$ develops a maximum below $`p_{T,J}^{\mathrm{min}}=5.5`$ GeV. This change of slope is due to the infrared sensitivity in the point $`p_{T,J}^{\mathrm{min}}=p_{T,\gamma }^{\mathrm{min}}`$. To avoid this region one must choose $`\mathrm{\Delta }0`$. From the plot we observe that $`\mathrm{\Delta }0.5`$ GeV would be sufficient. Thus, in principle, we could have used a smaller value for this difference than was chosen in (11).
The cross section in the vicinity of $`p_{T,J}^{\mathrm{min}}=p_{T,\gamma }^{\mathrm{min}}`$ cannot be predicted reliably. It depends on the technical cut $`y_0^J`$ as soon as one approaches the limit $`\mathrm{\Delta }=0`$. In fact, the infrared sensitive region is very much influenced by non-perturbative effects which are not included in our calculation. In any case it would be interesting to measure $`\sigma (p_T^{\mathrm{min}})`$ in order to investigate this non-perturbative region.
In Fig. 17a we present $`\sigma (p_{T,J}^{\mathrm{min}})`$ also for $`p_{T,J}^{\mathrm{min}}`$ below 5 GeV, i.e. for $`\mathrm{\Delta }<0`$. If $`p_{T,J}^{\mathrm{min}}`$ is increased starting from $`p_{T,J}^{\mathrm{min}}=4`$ GeV, the slope of $`\sigma (p_{T,J}^{\mathrm{min}})`$ changes at about $`p_{T,J}^{\mathrm{min}}=4.5`$ GeV and $`\sigma (p_{T,J}^{\mathrm{min}})`$ becomes smaller towards $`p_{T,J}^{\mathrm{min}}=5`$ GeV compared to the behaviour with constant slope. This stronger decrease of $`\sigma (p_{T,J}^{\mathrm{min}})`$ above $`p_{T,J}^{\mathrm{min}}=4.5`$ has its origin again in the infrared sensitivity of this region. For $`p_{T,J}^{\mathrm{min}}5`$ GeV, and $`\mathrm{\Delta }<0`$ the cross section approaches the same value as we have obtained for $`\mathrm{\Delta }0`$ in the region $`\mathrm{\Delta }>0`$. This means, $`\sigma (p_{T,J}^{\mathrm{min}})`$ is not singular at $`p_{T,J}^{\mathrm{min}}=5`$ GeV, nor has it a discontinuity.
The observed behaviour is of course not only visible at the specific value of $`p_T^{\mathrm{min}}=5`$ GeV which we have chosen. A similar dip is seen for larger and smaller values as well. It becomes more pronounced for smaller values and gets washed out for larger $`p_T^{\mathrm{min}}`$.
For completeness we show in Fig. 17b also the cross section with the roles of $`p_{T,J}^{\mathrm{min}}`$ and $`p_{T,\gamma }^{\mathrm{min}}`$ interchanged, i.e. we fix $`p_{T,J}^{\mathrm{min}}=5`$ GeV and study the dependence on the cut for the photon transverse momentum $`p_{T,\gamma }^{\mathrm{min}}`$. The behaviour of the cross section in Fig. 17b is very similar to the first case and shows the same infrared sensitive region. Inside the region $`0.5`$ GeV $`<\mathrm{\Delta }<0.5`$ GeV the cross sections agree inside numerical errors. Only for larger $`|\mathrm{\Delta }|`$ the dependences on the minimal transverse momenta of the jet and the photon become different. The cross sections shown in Fig. 17 have been calculated with the cone algorithm using $`R=1`$ and $`z_{\mathrm{cut}}=0.9`$.
## 4 Summary and Concluding Remarks
We have presented results of a next-to-leading order calculation of isolated photon production in large-$`Q^2`$ $`ep`$ scattering. Contributions from quark-to-photon fragmentation are explicitly taken into account. We have discussed numerical results for $`\gamma +(1+1)`$-jet and $`\gamma +(2+1)`$-jet cross sections as functions of transverse momenta, rapidity and other observables derived from photon and/or jet kinematic variables. Infrared sensitive regions, as for example the region of equal photon and jet $`p_T`$, are studied in detail.
We investigated several of these cross sections with respect to their photon isolation, jet cone size, scale and jet algorithm dependences. It was found that these dependences are rather weak. In particular, the scale dependence of the integrated cross section is very small, giving quite some confidence in the reliability of our predictions. Also the results depended very little on the choice of modern parton distribution functions for the proton.
We expect that the measurement of photon plus jet cross sections at HERA will contribute to testing perturbative QCD in the process $`\gamma ^{}p\gamma X`$, in an area which has not been studied yet. The calculation covers the range of large $`Q^2`$ where the results do not depend on parton distribution functions of the virtual photon. At even larger $`Q^2`$ additional contributions from $`Z`$ boson exchange become important, which have been neglected in the present calculation. It is no major problem to incorporate these missing parts. They are non-negligible and must be considered when experimental data become available at very large $`Q^2`$.
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# Vacuum Polarization of Massive Scalar Fields on the Black Hole Horizon
## I Introduction
Since Hawking’s discovery of black holes radiation , it has been an important subject to investigate quantum effects in black hole spacetimes and a great deal of effort has been done. The expectation values of the regularized vacuum polarization $`\varphi ^2`$ and the stress energy tensor $`T_{\mu \nu }`$ are important quantities to understand quantum effects in black hole spacetimes. They represent the effect of quantum fluctuations and back reaction on black hole spacetimes and play an important role in the context of semi-classical theory of quantum gravity.
For a scalar field in the large mass limit, of which Compton length is much smaller than the characteristic curvature scale of background geometry, we can evaluate regularized $`\varphi ^2`$ and $`T_{\mu \nu }`$ by using the DeWitt-Schwinger approximation. This is an adiabatic expansion in terms of inverse power series of the mass of scalar field. One can evaluate divergent and finite parts of $`\varphi ^2`$ and $`T_{\mu \nu }`$ up to arbitrary order. Frolov and Zelnikov suggest that the finite parts of $`\varphi ^2`$ and $`T_{\mu \nu }`$ calculated by the DeWitt-Schwinger expansion coincides with the value calculated by assuming the Hartle-Hawking state for the scalar field, although the DeWitt-Schwinger expansion does not impose any condition on the state of the scalar field. The justification is confirmed numerically near the horizon.
The Hartle-Hawking state is defined by requiring that the Euclidean Green’s function is regular on the horizon. Its physical meaning is that a black hole is placed in a cavity and is equilibrium with thermal black-body radiation. The thermal state is defined by the Euclidean Green’s function by imposing periodicity in the imaginary time coordinate with a period of the inverse of the temperature. On the order hand, in the DeWitt-Schwinger expansion, the condition which characterizes the thermal state is not imposed. So it is not obvious whether the DeWitt-Schwinger expansion yields the same finite part of the Hartle-Hawking state near the horizon.
In this paper, we consider scalar fields in static spherically symmetric background spacetimes. By assuming that the mass of scalar fields is large, we calculate $`\varphi ^2`$ in the thermal state with arbitrary temperature using point splitting method near the horizon. A similar calculation is already done by Anderson et al., but their expression of the vacuum polarization diverges on the horizon which contradicts with their numerical calculation. Our purpose is to derive the vacuum polarization which is regular on the horizon, and resolve why the finite part of $`\varphi ^2`$ in the Hartle-Hawking state near the horizon can be approximated by the DeWitt-Schwinger expansion which does not assume any condition on the thermal state.
The plan of this paper is as follows. In Sec.II, we review the method to calculate $`\varphi ^2`$ in a thermal state using point splitting regularization. In Sec.III, we present our method to calculate $`\varphi ^2`$ in the large mass limit near the black hole horizon. Sec.IV is devoted to conclusion and discussion. Throughout the paper, we use units such that $`\mathrm{}=c=G=k_B=1`$. Our sign conventions are those of Misner, Thorne and Wheeler .
## II derivation of vacuum polarization $`\varphi ^2`$ in a thermal state
We consider a minimally coupled scalar field with mass $`m`$ in a thermal state at arbitrary temperature $`T`$ in a Reissner-Nordström spacetime. The metric in Euclidean section $`\tau =it`$ is given by
$$ds^2=fd\tau ^2+\frac{1}{f}dr^2+r^2d\mathrm{\Omega }^2,$$
(1)
where $`f=(rr_+)(rr_{})/r^2`$, $`r_+`$ and $`r_{}`$ are the location of inner and outer horizon, respectively. The surface gravity of a black hole is given by $`\kappa =(r_+r_{})/(2r_+^2)`$. The vacuum polarization $`\varphi ^2`$ can be computed from the Euclidean Green’s function by noting the identity
$$\varphi ^2_{\text{unreg}}=\frac{1}{2}G^{(1)}(x,x)=iG_F(x,x)=G_E(x,x),$$
(2)
where $`G^{(1)}`$ is the Hadamard Green’s function, $`G_F`$ is the Feynman Green’s function, and $`G_E`$ is the Euclidean Green’s function, respectively. $`G_E`$ is divergent at coincident limit and must be regularized. In this paper, the covariant point-splitting method is employed for the regularization of ultraviolet divergences. We start from an expression for $`G_E(x,x^{})`$ when the points $`x`$ and $`x^{}`$ are separated. Next we prepare an appropriate counterterm to subtract the divergence, and then take coincident limit. This procedure is shown schematically as
$$\varphi ^2_{\text{reg}}=\underset{x^{}x}{lim}\left[G_E(x,x^{})G_{DS}(x,x^{})\right],$$
(3)
where $`G_{DS}(x,x^{})`$ is a point splitting counterterm needed to regularize $`G_E(x,x^{})`$. The Euclidean Green’s function obeys the equation
$$\left[\mathrm{}m^2\right]G_E(x,x^{})=\frac{\delta ^4(x,x^{})}{g^{1/2}(x)},$$
(4)
where $`m`$ is the mass of the scalar field and d’Alembertian $`\mathrm{}`$ evaluated using the Euclidean metric (1). The point splitting counter term needed to renormalize $`G_E(x,x^{})`$ in arbitrary spacetime is
$$\begin{array}{cc}\hfill G_{DS}(x,x^{})& =\frac{1}{8\pi ^2\sigma }+\frac{m^2+(\xi \frac{1}{6})R}{8\pi ^2}\left[\gamma +\frac{1}{2}\mathrm{ln}\left[\frac{m^2\sigma }{2}\right]\right]\hfill \\ & \frac{m^2}{16\pi ^2}+\frac{1}{96\pi ^2}R_{\alpha \beta }\frac{\sigma ^{,\alpha }\sigma ^{,\beta }}{\sigma },\hfill \end{array}$$
(5)
where $`\xi `$ is coupling constant to the scalar curvature $`R`$, $`\sigma `$ is the one half the square of geodesic distance between the points $`x`$ and $`x^{}`$ and $`\gamma `$ is Euler’s constant.
For simplicity of calculation, the quantity $`\varphi ^2`$ is evaluated by timelike point separation $`\tau ^{}\tau `$. By assuming that the Green’s function is periodic in $`ϵ=\tau \tau ^{}`$ with period $`1/T`$, the expression for the vacuum polarization in the thermal state at temperature $`T`$ is given by
$`\varphi ^2_T`$ $`=\underset{\tau ^{}\tau }{lim}\left[G_E(x,\tau ;x,\tau ^{})G_{DS}(x,\tau ;x,\tau ^{})\right]`$
$`=\underset{\tau ^{}\tau }{lim}[{\displaystyle \frac{T}{4\pi }}{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}[(2l+1)p_{0l}q_{0l}{\displaystyle \frac{1}{r\sqrt{f}}}]`$
$`+{\displaystyle \frac{T}{2\pi }}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\mathrm{cos}[n2\pi T(\tau \tau ^{})]{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}[(2l+1)p_{nl}q_{nl}{\displaystyle \frac{1}{r\sqrt{f}}}]G_{DS}(x,\tau ;x,\tau ^{})]`$
$`\varphi ^2_{n=0}+\varphi ^2_{n1}\varphi ^2_{DS},`$ (6)
where
$$\varphi ^2_{DS}=\frac{m^2}{8\pi ^2}\left(\mathrm{ln}\frac{m\sqrt{f}ϵ}{2}+\gamma \frac{1}{2}\right)\frac{1}{4\pi ^2fϵ^2}+\left(\frac{f^{}^2}{192\pi ^2f}\frac{f^{\prime \prime }}{96\pi ^2}\frac{f^{}}{48\pi ^2r}\right).$$
(7)
The term $`1/(r\sqrt{f})`$ is subtracted to cancel the superficial divergence which comes from the choice of timelike point-splitting. $`p_{nl}`$ and $`q_{nl}`$ are independent solutions of the mode equation
$`{\displaystyle \frac{d^2\chi _{nl}}{dr^2}}+\left[{\displaystyle \frac{2}{r}}+{\displaystyle \frac{1}{f}}{\displaystyle \frac{df}{dr}}\right]{\displaystyle \frac{d\chi _{nl}}{dr}}\left[{\displaystyle \frac{(2\pi nT)^2}{f^2}}+{\displaystyle \frac{l(l+1)}{r^2f}}+{\displaystyle \frac{m^2}{f}}\right]\chi _{nl}=0`$ (8)
and satisfy the Wronskian condition
$$p_{nl}\frac{dq_{nl}}{dr}q_{nl}\frac{dp_{nl}}{dr}=\frac{1}{r^2f}.$$
(9)
Anderson et al. calculated the second and the third term of (II) analytically using the second order WKB approximation . Their result is
$$\begin{array}{cc}\hfill \varphi ^2_{n1}\varphi ^2_{DS}& =\frac{m^2}{16\pi ^2}\mathrm{ln}\left(\frac{m^2f}{16\pi ^2T^2}\right)+\frac{m^2}{16\pi ^2}\hfill \\ & +\frac{T^2}{12f}\frac{f_{}^{}{}_{}{}^{2}}{192\pi ^2f}\frac{f^{\prime \prime }}{96\pi ^2}\frac{f^{}}{48\pi ^2r}.\hfill \end{array}$$
(10)
The expression (10) was obtained by taking into account only the contribution from $`n1`$ mode. In the case of massless scalar fields, excluding the lowest frequency mode $`n=0`$ corresponds to imposing an infrared cutoff. The expression (10) gives good approximation for massless fields both near and far from the horizon in Schwarzschild and Reissner-Nordström spacetimes. This was confirmed by their numerical calculation.
When the result (10) is applied to the massive scalar field, however, two unpleasant features appear. Firstly, the regularized expression always has logarithmic divergence on the horizon $`f=0`$ even if the scalar field is in the Hartle-Hawking state. It is not consistent with their numerical work which shows regular behavior of $`\varphi ^2_T`$ on the horizon. Secondly, the finite terms of $`O(m^2)`$ and $`O(m^0)`$ are left in the regularized expression. This disagrees with the result of the DeWitt-Schwinger approximation by Frolov which does not contain such finite terms. Therefore the expression (10) cannot yield good approximation for the massive scalar field. In the next section, we present the method to improve their approximations and resolve these problems.
## III vacuum polarization near the horizon in the large mass limit
In this section, we present the method to resolve the unwanted behavior arose in the expression of $`\varphi ^2_T`$ near the horizon $`f0`$. We assume that the mass $`m`$ of the scalar field is large:
$$m\frac{1}{r_+}.$$
(11)
We evaluate $`n1`$ and $`n=0`$ contribution to $`\varphi ^2_T`$ separately.
### A $`n1`$ contribution
We can use the WKB approximation for $`n1`$ mode. The zeroth order WKB solution is given by
$`p_n={\displaystyle \frac{1}{\sqrt{2r^2W}}}\mathrm{exp}\left[{\displaystyle 𝑑r\frac{W}{f}}\right],q_n={\displaystyle \frac{1}{\sqrt{2r^2W}}}\mathrm{exp}\left[{\displaystyle 𝑑r\frac{W}{f}}\right],`$
$`;W^2=\kappa ^2n^2+m^2f+{\displaystyle \frac{\left(l+\frac{1}{2}\right)^2}{r^2}}f=(\alpha _n^2+{\displaystyle \frac{\left(l+\frac{1}{2}\right)^2}{r^2}})f,`$ (12)
$`\alpha _n^2=m^2+{\displaystyle \frac{\kappa ^2n^2}{f}},\kappa =2\pi T.`$
The WKB approximation is correct near the horizon. Using this solution, we have
$$\varphi ^2_{n1}=\frac{T}{2\pi }\underset{n=1}{\overset{\mathrm{}}{}}\mathrm{cos}(n\kappa ϵ)I_n,$$
(13)
where
$`I_n`$ $`{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}\left[(2l+1)p_nq_n{\displaystyle \frac{1}{r\sqrt{f}}}\right]`$
$`={\displaystyle \frac{1}{r\sqrt{f}}}{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}\left\{{\displaystyle \frac{l+\frac{1}{2}}{\left[r^2\alpha _n^2+\left(l+\frac{1}{2}\right)^2\right]^{1/2}}}1\right\}.`$ (14)
$`I_n`$ can be evaluated using the Plana sum formula under the assumption $`f0,\alpha _n\kappa n/f`$:
$`{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{l+\frac{1}{2}}{\left[r^2\alpha _n^2+\left(l+\frac{1}{2}\right)\right]^{1/2}}}`$ $`={\displaystyle \frac{1}{2}}{\displaystyle \frac{\frac{1}{2}}{\left[r^2\alpha _n^2+\frac{1}{4}\right]^{1/2}}}+{\displaystyle _0^{\mathrm{}}}𝑑l{\displaystyle \frac{l+\frac{1}{2}}{\left[r^2\alpha _n^2+\left(l+\frac{1}{2}\right)^2\right]^{1/2}}}`$
$`+i{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dl}{e^{2\pi l}1}}\left\{{\displaystyle \frac{il+\frac{1}{2}}{\left[r^2\alpha _n^2+\left(il+\frac{1}{2}\right)^2\right]^{1/2}}}(ll)\right\}`$
$`=\left[\left(l+{\displaystyle \frac{1}{2}}\right)^2+r^2\alpha _n^2\right]^{1/2}|_{l=0}^{\mathrm{}}+O\left({\displaystyle \frac{f}{r\kappa n}}\right).`$ (15)
Therefore, for $`n1`$,
$`I_n`$ $`={\displaystyle \frac{1}{r\sqrt{f}}}\left\{\left[{\displaystyle \frac{1}{4}}+r^2\left(m^2+{\displaystyle \frac{r^2\kappa ^2n^2}{f^2}}\right)\right]^{1/2}+O\left({\displaystyle \frac{f}{r\kappa n}}\right)\right\}`$
$`={\displaystyle \frac{\kappa n}{f}}{\displaystyle \frac{m^2}{2\kappa n}}+O({\displaystyle \frac{f}{rn^3}},{\displaystyle \frac{\sqrt{f}}{r^2\kappa n}}).`$ (16)
By taking $`n`$ sum, $`O(f/(rn^3),\sqrt{f}/(r^2\kappa n))`$ terms give finite values which becomes zero on the horizon and we can omit these terms. Using the formula
$$\begin{array}{cc}\hfill \kappa \underset{n=1}{\overset{\mathrm{}}{}}\frac{\mathrm{cos}(n\kappa ϵ)}{n\kappa }& =\frac{1}{2}\mathrm{ln}(\kappa ^2ϵ^2)+O(ϵ^2),\hfill \\ \hfill \kappa \underset{n=1}{\overset{\mathrm{}}{}}n\kappa \mathrm{cos}(n\kappa ϵ)& =\frac{1}{ϵ^2}\frac{\kappa ^2}{12}+O(ϵ^2),\hfill \end{array}$$
(17)
we have
$`\varphi ^2_{n1}`$ $`{\displaystyle \frac{\kappa ^2}{4\pi ^2f}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}n\mathrm{cos}(n\kappa ϵ){\displaystyle \frac{m^2}{8\pi ^2}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\mathrm{cos}(n\kappa ϵ)}{n}}`$
$`={\displaystyle \frac{1}{4\pi ^2fϵ^2}}+{\displaystyle \frac{\kappa ^2}{48\pi ^2f}}+{\displaystyle \frac{m^2}{16\pi ^2}}\mathrm{ln}(kappa^2ϵ^2).`$ (18)
Using the DeWitt-Schwinger counter term (7), the regularized expression of the expectation value becomes
$$\varphi ^2_{n1}\varphi ^2_{DS}=\frac{m^2}{16\pi ^2}\mathrm{ln}\left(\frac{m^2f}{16\pi ^2T^2}\right)+\frac{m^2}{16\pi ^2}+\frac{T^2}{12f}\frac{f_{}^{}{}_{}{}^{2}}{192\pi ^2f}\frac{f^{\prime \prime }}{96\pi ^2}\frac{f^{}}{48\pi ^2r}.$$
(19)
This gives the same expression of Anderson et al., who used the second order WKB solution. Near the black hole horizon, it is sufficient to use the zeroth order WKB solution to reproduce the result of the second order WKB approximation. On the horizon $`f=0`$, this expression diverge and we must examine the contribution from $`n=0`$ mode.
### B $`n=0`$ contribution
We must recall that the expression (10) does not include $`n=0`$ mode but there is no reason to exclude this contribution in the case of the massive scalar field. So we must investigate the contribution of $`n=0`$ mode. Since the WKB method breaks down for $`n=0`$ mode near the horizon, we cannot apply the approximation used by Anderson et al.. We solve $`n=0`$ mode function by the following method.
As we assume that the mass of the scalar field is large, a dimensionless constant
$$ϵ\frac{r_+r_{}}{4r_+[m^2r_+^2+l(l+1)]}$$
(20)
becomes smaller than unity for all value of $`l`$:
$$ϵ1(\text{for all }l)$$
(21)
We rescale the radial coordinate $`r`$ as follows:
$$x\sqrt{\frac{rr_+}{ϵr_+}}.$$
(22)
Assuming that $`x`$ is $`O(1)`$ quantity, this rescaling of the radial coordinate means we are concentrate on the region $`(rr_+)O(ϵ)`$, which is the region near the black hole horizon. Using $`x`$ and $`ϵ`$, the radial equation (8) for $`n=0`$ mode can be written as
$`{\displaystyle \frac{d^2\chi _{0l}}{dx^2}}+\left[{\displaystyle \frac{1}{x}}+F(x,ϵ)\right]{\displaystyle \frac{d\chi _{0l}}{dx}}\left[1+G(x,ϵ)\right]\chi _{0l}=0,`$ (23)
where
$`F(x,ϵ)`$ $`={\displaystyle \frac{2ϵx}{2\kappa r_++ϵx^2}},`$ (24)
$`G(x,ϵ)`$ $`={\displaystyle \frac{16ϵ}{2\kappa r_++ϵx^2}}\left[\kappa r_++2m^2r_+^2(2ϵ^2x^2+3ϵ^3x^4)\right].`$ (25)
We can express the approximate solution of the mode equation (23) by the power series expansion with respect to a small parameter $`ϵ`$:
$$\chi _{0l}(x,\eta )=\chi _{0l}^{(0)}(x)+ϵ\chi _{0l}^{(1)}(x)+ϵ^2\chi _{0l}^{(2)}(x)+\mathrm{}.$$
(26)
$`\chi _{0l}^{(n)}`$ obeys
$`{\displaystyle \frac{d\chi _{0l}^{(0)}}{dx^2}}+{\displaystyle \frac{1}{x}}{\displaystyle \frac{d\chi _{0l}^{(0)}}{dx}}\chi _{0l}^{(0)}=0,`$ (27)
$`{\displaystyle \frac{d\chi _0^{(n)}}{dx^2}}+{\displaystyle \frac{1}{x}}{\displaystyle \frac{d\chi _{0l}^{(n)}}{dx}}\chi _{0l}^{(n)}=`$
$`{\displaystyle \underset{k=1}{\overset{n}{}}}{\displaystyle \frac{1}{k!}}\left[{\displaystyle \frac{^kF}{ϵ^k}}|_{ϵ=0}{\displaystyle \frac{d\chi _{0l}^{(nk)}}{dx}}+{\displaystyle \frac{^kG}{ϵ^k}}|_{ϵ=0}\chi _{0l}^{(nk)}\right](\text{for }n1)`$ (28)
We can obtain solutions of these equations by using modified Bessel functions. For example, the lowest order solution is given by
$$p_{l0}^{(0)}=I_0(x),q_{l0}^{(0)}=\frac{2}{r_+r_{}}K_0(x).$$
(29)
Substituting the expansion (26) to
$$\varphi ^2_{n=0}=\frac{T}{4\pi }\underset{l=0}{\overset{\mathrm{}}{}}\left[(2l+1)p_{l0}q_{l0}\frac{1}{r\sqrt{f}}\right]$$
(30)
and taking summation about $`l`$ using the Plana sum formula, $`\varphi ^2_{n=0}`$ is written as the power series expansion with respect to $`(mr_+)^1`$. We calculated $`\varphi ^2_{n=0}`$ up to $`O\left((mr_+)^4\right)`$, which is accomplished by the $`ϵ`$ expansion of the mode function up to $`O(ϵ^3)`$. By taking into the account of the contribution of $`n=0`$ mode, we finally get the form of the vacuum polarization near the horizon as follows:
$`\varphi ^2_T`$ $`=\left(T{\displaystyle \frac{\kappa }{2\pi }}\right)\left\{{\displaystyle \frac{m^2}{16\pi \kappa }}+{\displaystyle \frac{m^2}{8\pi \kappa }}\left[2\gamma +\mathrm{ln}{\displaystyle \frac{m^2(rr_+)}{2\pi T}}\right]{\displaystyle \frac{1}{12\pi r_+}}\right\}`$
$`+\left\{T^2\left({\displaystyle \frac{\kappa }{2\pi }}\right)^2\right\}\left[{\displaystyle \frac{1}{24\kappa (rr_+)}}{\displaystyle \frac{r_{}}{48r_+^3\kappa ^2}}\right]`$ (31)
$`+{\displaystyle \frac{m^2}{16\pi ^2\kappa }}(\kappa +2\pi T)\mathrm{ln}\left({\displaystyle \frac{2\pi T}{\kappa }}\right)`$
$`+{\displaystyle \frac{T}{4\pi }}{\displaystyle \frac{1}{45m^2}}{\displaystyle \frac{1}{(r_+r_{})r_+^2}}\left[3{\displaystyle \frac{6r_{}}{r_+}}+{\displaystyle \frac{4r_{}^2}{r_+^2}}\right]`$
$`+{\displaystyle \frac{T}{4\pi }}{\displaystyle \frac{2}{315m^4}}{\displaystyle \frac{1}{(r_+r_{})r_+^4}}\left[10+39{\displaystyle \frac{r_{}}{r_+}}52{\displaystyle \frac{r_{}^2}{r_+^2}}+23{\displaystyle \frac{r_{}^3}{r_+^3}}\right]+O\left((mr_+)^6\right).`$ (32)
The divergences in $`\varphi ^2_T`$ appears only in the terms of $`O(m^2)`$ and $`O(m^0)`$. Taking $`f0`$ limit, $`O(m^2)`$ term diverge as $`\mathrm{ln}(rr_+)`$ and $`O(m^0)`$ as $`(rr_+)^1`$. These divergence disappear completely if the temperature of the scalar field equals $`T=\kappa /2\pi T_H`$, where $`T_H`$ is the Hawking temperature. Specifying $`T=T_H`$ is equivalent to require the regurality of $`\varphi ^2`$ on the horizon. In this case, the black hole is equilibrium with the thermal scalar field with temperature $`T=T_H`$. This is nothing but a Hartle-Hawking state.
We can evaluate a finite part of $`\varphi ^2_T`$ in the Hartle-Hawking state on the horizon. The leading order of $`\varphi ^2_T`$ at $`T=T_H`$ is $`O(m^2)`$ because not only divergent but also finite terms contained in $`O(m^2)`$ and $`O(m^0)`$ terms vanish. The result is given by
$$\begin{array}{cc}\hfill \varphi ^2_{T=T_H}& =\frac{1}{720\pi ^2m^2}\frac{1}{r_+^4}\left[3\frac{6r_{}}{r_+}+\frac{4r_{}^2}{r_+^2}\right]\hfill \\ & +\frac{1}{2520\pi ^2m^4}\frac{1}{r_+^6}\left[10+39\frac{r_{}}{r_+}52\frac{r_{}^2}{r_+^2}+23\frac{r_{}^3}{r_+^3}\right]+O\left(m^6\right).\hfill \end{array}$$
(33)
This expression agrees with the result of the DeWitt-Schwinger expansion .
## IV Conclusion and discussion
We calculated $`\varphi ^2_T`$ of a massive scalar field near the horizon of a Reissner-Nordström black hole using point splitting method. Our results are as follows: (i)all divergent and finite terms of $`O(m^2)`$ and $`O(m^0)`$ in $`\varphi ^2_T`$ are canceled to be zero if the temperature of the scalar field equals to the Hawking temperature $`T_H`$ of a black hole. (ii)at this temperature, the leading order of the regularized $`\varphi ^2_T`$ becomes $`O(m^2)`$ and its value on the horizon agrees with that of the DeWitt-Schwinger expansion up to $`O(m^4)`$. (i) and (ii) indicate that regularized $`\varphi ^2_T`$ in the Hartle-Hawking state near the black hole horizon is well approximated by the DeWitt-Schwinger expansion. This confirms the results of .
It is not trivial that the finite part of $`\varphi ^2_T`$ in the thermal state at $`T=T_H`$ is reproduced by the DeWitt-Schwinger expansion which does not specify any thermal state of the scalar field. It was shown that the leading order of the DeWitt-Schwinger expansion reproduces the value $`\varphi ^2_{T=0}`$ of scalar field with large mass in general static spherical symmetric spacetime . This result is natural because the scalar field with $`T=0`$ is in a vacuum state and the DeWitt-Schwinger expansion does not incorporate information on a thermal state.
Discrepancy between $`\varphi ^2_T`$ in the Hartle-Hawking state and that in he DeWitt-Schwinger expansion will appear far from the horizon, because the thermal effect dominates the gravitational effect in the asymptotic flat region. If the thermal effect is dominant or comparable to the contribution from curvature effect near the horizon, there would also be discrepancy between $`\varphi ^2_T`$ in the Hartle-Hawking state and that of the DeWitt-Schwinger expansion. However this is not seen from our results. This implies that the finite part of $`\varphi ^2_T`$ in the Hartle-Hawking state near the black hole horizon is dominated by the contribution from curvature effect. Why the contribution from thermal effects to finite parts can be negligible? The characteristic scale for thermal state is the inverse of the temperature $`1/T_HM`$ which is a scale length of thermal fluctuation. On the other hand, the Compton length of the scalar field is $`1/m`$. The contribution of the Hawking radiation to $`\varphi ^2`$ is $`T_H^4/m^2`$ and that of thermal excitation is $`\mathrm{exp}(m/T_H)`$. For $`m1/M`$, the effect of thermal excitation is exponentially small. Furthermore, near the black hole horizon, almost all part of thermally excited particle is absorbed by the black hole. This is the reason of the DeWitt-Schwinger approximation gives the same result of Hartle-Hawking state on the horizon.
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# Causality, delocalization and positivity of energy.
## 1 Introduction
Are there deviations from Einstein’s causality? G. C. Hegerfeldt has written : “Positivity of the Hamiltonian alone is used to show that particles, if initially localized in a finite region, immediately develop infinite tails.” This seems to imply superluminality. One of his examples is the Fermi problem of two atoms coupled by a radiation field. Consider the initial condition when one of the atoms is in an excited state, the other in the ground state, and no photons are present. The probability to find the second atom in an excited state is non-vanishing immediately after the initial moment, independently of the distance between the atoms , \- . Hegerfeldt’s arguments are based on the analyticity of the expectation values of the operator $`N(V)`$, which gives the probability to find a particle inside a finite volume $`V`$. He showed that a state in a quantum system with positive energy localized in a finite volume $`V`$ at the instant $`t=0`$, will develop infinite tails immediately afterwards. Positivity of energy plays an essential role in his proof. In this paper we present an illustration of Hegerfeldt’s theorem, without any appeal to superluminality. We apply Hegerfeldt’s consideration to wave packets. Moreover, we show that Hegerfeldt’s effect appears even for classical fields, if wave packets are constructed from positive frequencies (corresponding positive energy quantum fields).
We first study the positive-frequency solutions of the classical wave equation (section 2). We consider wave packets $`\mathrm{\Phi }(x,t)`$ localized at $`t=0`$. The localization is due to interference of the two complex solutions, each propagating causally
$$\mathrm{\Phi }(x,t)=\psi (xt)+\psi ^{}(x+t)$$
(1)
where “” denotes complex conjugation and we take $`c=1`$. We show that both these wave packets are delocalized. They present long tails, extending to arbitrary distances and decaying according to a power law. As we shall show, the “nonlocal effect” can also be understood from the point of view of the initial conditions. Indeed, in our construction of the solution (1) we shall use two conditions; one is the initial condition of the local shape of the field $`\mathrm{\Phi }(x,t)`$ and the other is the condition of the positivity of frequencies. The frequency positivity replaces the usual initial condition on the time derivative of the field $`\mathrm{\Phi }/t`$. We shall show that our initial condition with positive frequencies leads to the nonlocality of $`\mathrm{\Phi }/t`$ at $`t=0`$.
In section 3 we show that similar conclusions are obtained for the wave packet of a free field in relativistic quantum field theory. We construct an operator reminding of the position operator of Newton-Wigner . The expectation value of this operator with the state corresponding to our wave packet is local at $`t=0`$. However, it has infinite tails which are “hidden” at time $`t=0`$, but emerge immediately afterwards. We may call this effect a “curtain” effect. No superluminal propagation is involved. We note that at the same time other quantities such as the energy density have a nonlocal expectation value in the same state even for $`t=0`$. It should be also pointed out that for the Dirac equation there are no positive energy solutions which can be localized in a finite region (see ). This demonstrates that localization in relativistic quantum field theory cannot be “complete”.
## 2 Classical wave packets
Consider classical wave packets constructed by the solutions of the wave equation with positive frequency and localized at time $`t=0`$. We show that these wave packets will spread immediately over the whole space. Curiously we have not found any reference to this effect in the literature. We start from the wave equation on the real line ($`c=1`$).
$$\left(\frac{^2}{t^2}\frac{^2}{x^2}\right)\mathrm{\Phi }(x,t)=0$$
(2)
The general complex solution of Eq. (2) is, by the Fourier transform, of the form
$$\mathrm{\Phi }(x,t)=\frac{1}{2\pi }_{\mathrm{}}^{\mathrm{}}𝑑k\left\{\varphi _+(k)e^{i\omega _kt}+\varphi _{}(k)e^{i\omega _kt}\right\}e^{ikx}$$
(3)
where $`\omega _k=|k|`$ and where $`\varphi _\pm (k)`$ are arbitrary functions. To determine $`\varphi _+`$ and $`\varphi _{}`$ one can use the two initial conditions $`\mathrm{\Phi }(x,0)`$ and $`\dot{\mathrm{\Phi }}(x,0)`$. However, one can also consider the special class of positive-frequency solutions to Eq. (2), i.e. $`\varphi _{}(k)0`$ and
$$\mathrm{\Phi }_+(x,t)=\frac{1}{2\pi }_{\mathrm{}}^{\mathrm{}}𝑑k\varphi _+(k)e^{i\omega _kt}e^{ikx}$$
(4)
These positive-frequency solutions are determined by the initial condition $`\mathrm{\Phi }(x,0)`$. Note that relation (4) leads to a complex field for $`t0`$, even if $`\mathrm{\Phi }_+(x,0)`$ or $`\varphi _+(k)`$ are real. Consider as an example a localized (rectangular) wave packet with center $`x_0`$, and width $`2b`$ at time $`t=0`$:
$$\mathrm{\Phi }_{x_0,b}(0)=\frac{1}{2b}\mathrm{\Theta }(b|xx_0|)$$
(5)
The normalization has been chosen so that the integral of this function over $`x`$ is equal to one. Then the function $`\varphi (k)`$ is:
$$\varphi _+(k)=\frac{1}{2b}_{\mathrm{}}^+\mathrm{}𝑑xe^{ikx}\mathrm{\Theta }(b|xx_0|)$$
(6)
where $`\mathrm{\Theta }(x)`$ is the step function, which is $`0`$ for $`x`$ negative, and $`1`$ for $`x`$ positive. Then the function $`\mathrm{\Phi }(x,t)`$ in (4) is given by
$$\mathrm{\Phi }(x,t)=\frac{1}{4\pi b}_{\mathrm{}}^+\mathrm{}𝑑k_{x_0b}^{x_0+b}𝑑x^{}e^{i|k|t+ik(xx^{})}$$
(7)
This is a sum of two functions corresponding to two wave packets moving in opposite directions,
$$\mathrm{\Phi }(x,t)=\psi (xt)+\psi ^{}(x+t)$$
(8)
where
$$\psi (x)=\frac{1}{4\pi b}_{x_0b}^{x_0+b}𝑑x^{}_0^+\mathrm{}𝑑ke^{ik(xx^{})}$$
(9)
To evaluate the integral over $`k`$ we introduce the usual regularisation by adding a positive infinitesimal to $`x`$, which leads to
$$\psi (x)=\frac{1}{4\pi bi}_{x_0b}^{x_0+b}\frac{dx^{}}{xx^{}+i0}$$
(10)
After integration over $`x^{}`$ we obtain:
$$\psi (x)=\frac{i}{4\pi b}[\mathrm{ln}(xx_0+b+i0)\mathrm{ln}(xx_0b+i0)]$$
(11)
The logarithm of a complex number is given by
$$\mathrm{ln}(z)=\mathrm{ln}|z|+i(\mathrm{arg}(z)+2\pi n)$$
(12)
where $`n`$ is an integer. In order to have both terms in (11) on the same branch of the logarithm we take $`n=0`$ for both of them (due to the difference of the two terms in (11) the result does not depend on the particular value of $`n`$). The argument of $`x+i0`$ can be expressed as
$$\mathrm{arg}(x+i0)=\frac{\pi }{2}(1\mathrm{sign}(x))$$
(13)
where $`\mathrm{sign}(x)=x/|x|`$ is the sign of $`x`$. Then, inserting (12) and (13) into (11) we obtain
$$\psi (x)=\frac{1}{8b}\left(\mathrm{sign}(xx_0+b)\mathrm{sign}(xx_0b)\right)+\frac{i}{4\pi b}\mathrm{ln}\left|\frac{xx_0+b}{xx_0b}\right|$$
(14)
We see that the function $`\psi (x)`$ in (14) consists of a local real part ($`\mathrm{sign}`$) and a nonlocal imaginary part ($`\mathrm{log}`$). For $`t0`$ it is sufficient to replace $`x`$ by $`xt`$ in (8). Similar result is obtained for $`\psi ^{}(x+t)`$. As a result, the function $`\mathrm{\Phi }(x,t)`$ is also nonlocal because it is the superposition of the two complex functions $`\psi (xt)`$ and $`\psi ^{}(x+t)`$ in (8), which describe nonlocal objects moving with the speed of light in opposite directions. However, at $`t=0`$ the imaginary parts cancel each other (see Fig. 1), and we recover our localized initial condition (5), because only the real parts of these functions, which are local, remain. In all our figures time $`t`$ is measured in seconds ($`s`$), the coordinate $`x`$ is measured in “light seconds” ($`ls`$) and wave packet amplitudes are dimensionless. Fig. 2 corresponds to $`t=0.25s`$. At this moment the real local parts of $`\psi (xt)`$ and $`\psi ^{}(x+t)`$ have moved in opposite directions. The nonlocal imaginary parts of $`\psi (xt)`$ and $`\psi ^{}(x+t)`$ have also shifted in opposite directions and no more cancel each other. At this time, the two waves overlap and we have
$$|\mathrm{\Phi }(x,t)||\psi (xt)+\psi ^{}(x+t)||\psi (xt)|+|\psi ^{}(x+t)|$$
(15)
At $`t=1s`$ in Fig. 3 the overlapping is small and we have
$$|\mathrm{\Phi }(x,t)||\psi (xt)|+|\psi ^{}(x+t)|$$
(16)
We see that the initial condition $`\mathrm{\Phi }(x,0)`$ is local (Fig. 1) only because at $`t=0`$ the nonlocal parts cancel each other completely by destructive interference. We may describe the appearance of nonlocality as a sort of “curtain effect”. The nonlocal nature of each wave packet $`\psi (xt)`$ and $`\psi ^{}(x+t)`$ is hidden behind a “curtain” at the initial time and emerges immediately afterwards. Each of the nonlocal wave packets is complex and propagates at the speed of light.
In conclusion, we have illustrated Hegerfeldt’s theorem for classical wave packets. We see that the localization of wave packets corresponding to positive frequency is unstable and involves complex space structures.
Note that the localization of a positive frequency wave packet is not “complete”, because the time derivative of the function $`\mathrm{\Phi }(x,t)`$ is nonlocal even at $`t=0`$:
$$\left[\frac{\mathrm{\Phi }(x,t)}{t}\right]_{t=0}=\frac{i}{2\pi b}\left(\frac{1}{xx_0b}\frac{1}{xx_0+b}\right)$$
(17)
The wave equation being of second order demands two initial conditions: for the function itself and for its time derivative. The additional requirement of positivity of frequency replaces the second condition. There are no wave packets containing only positive frequency modes, which are localized together with their time derivative .
## 3 Relativistic quantum field
We turn now to relativistic quantum field theory. We show that the previous discussion is applicable to relativistic quantum particles. In this case the condition $`\omega _k>0`$ appears naturally since the energy $`E=\mathrm{}\omega _k`$ must be positive (we take $`\mathrm{}=1`$). We consider massless particles with no spin (“photons”). To simplify our consideration we use again a 1+1-dimensional spacetime. In terms of second quantization we have the scalar field operator
$$\widehat{\psi }(x,t)=_{\mathrm{}}^+\mathrm{}𝑑\stackrel{~}{k}\left(a_k^{}e^{i(\omega _ktkx)}+a_ke^{i(\omega _ktkx)}\right)$$
(18)
where $`d\stackrel{~}{k}=dk/(4\pi \omega _k)`$ is a relativistic invariant measure and
$$\omega _k=|k|$$
(19)
with $`c=1`$. The creation and annihilation operators $`a_k^{}`$ and $`a_k`$ of the photon with wave vector $`k`$ obey the commutation relation
$$[a_k^{},a_k^{}]=4\pi \omega _k\delta (kk^{})$$
(20)
We construct a wave packet from a linear combination of normal modes.
$$|\mathrm{\Phi }_{x_0,b}(t)=_{\mathrm{}}^+\mathrm{}𝑑\stackrel{~}{k}\varphi _{x_0,b}(k)e^{i\omega _kt}a_k^{}|0$$
(21)
Here $`|0`$ is the vacuum state for the field (18). The fact that our wave packet is obtained by the action of creation operators on the vacuum state implies that the state consists of normal modes with positive energy. As before, we chose the function $`\varphi _{x_0,b}(k)`$ so that the wave packet is localized at the time $`t=0`$ in a domain with center $`x_0`$ and width $`2b`$,
$$\varphi _{x_0,b}(k)=(2\omega _k)^{1/2}_{\mathrm{}}^+\mathrm{}𝑑xe^{ikx}\mathrm{\Phi }_{x_0,b}(x,0)$$
(22)
where
$$\mathrm{\Phi }_{x_0,b}(x,0)=\frac{1}{(2b)^{1/2}}\mathrm{\Theta }(b|xx_0|)$$
(23)
The function $`\mathrm{\Phi }_{x_0,b}(x,0)`$ is normalized to ensure the normalization of the state $`|\mathrm{\Phi }_{x_0,b}(0)`$ in (21).
Let us introduce the operator $`\rho (x)`$:
$$\rho (x)=a^{}(x)a(x)$$
(24)
where $`a^{}(x)`$ and $`a(x)`$ are defined by
$`a^{}(x)`$ $`=`$ $`{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑\stackrel{~}{k}(2\omega _k)^{1/2}e^{ikx}a_k^{}`$
$`a(x)`$ $`=`$ $`{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑\stackrel{~}{k}(2\omega _k)^{1/2}e^{ikx}a_k`$ (25)
These operators satisfy the commutation relation
$$[a(x),a^{}(x^{})]=\delta (xx^{})$$
(26)
This construction follows the ideas of the construction of positions operators by Newton and Wigner . We shall express the localization of our state $`|\mathrm{\Phi }_{x_0,b}(t)`$ in terms of the expectation value of the operator $`\rho (x)`$. We call a state localized if the expectation value of $`\rho (x)`$ in this state vanishes when $`x`$ is outside a finite region. Our choice of $`\varphi _{x_0,b}(k)`$ in (22) and $`\mathrm{\Phi }_{x_0,b}(x)`$ in (23) guarantees that the state $`|\mathrm{\Phi }_{x_0,b}(t)`$ is localized at $`t=0`$ in the domain $`[x_0b,x_0+b]`$, i.e.,
$$\mathrm{\Phi }_{x_0,b}(0)|\rho (x)|\mathrm{\Phi }_{x_0,b}(0)=\frac{1}{2b}\mathrm{\Theta }(b|xx_0|)$$
(27)
Using (25) we obtain the time evolution of this quantity
$$\mathrm{\Phi }_{x_0,b}(t)|\rho (x)|\mathrm{\Phi }_{x_0,b}(t)=_{\mathrm{}}^+\mathrm{}𝑑\stackrel{~}{k}𝑑\stackrel{~}{k}^{}(4\omega _k\omega _k^{})^{1/2}e^{i(kk^{})x}\mathrm{\Phi }_{x_0,b}(t)|a_k^{}a_k^{}|\mathrm{\Phi }_{x_0,b}(t)$$
(28)
where $`\mathrm{\Phi }_{x_0,b}(t)|a_k^{}a_k^{}|\mathrm{\Phi }_{x_0,b}(t)`$ is expressed using our form of the wave packet (21) as follows:
$$\mathrm{\Phi }_{x_0,b}(t)|a_k^{}a_k^{}|\mathrm{\Phi }_{x_0,b}(t)=_{\mathrm{}}^+\mathrm{}𝑑\stackrel{~}{l}𝑑\stackrel{~}{l}^{}\varphi _{x_0,b}^{}(l)\varphi _{x_0,b}(l^{})e^{i(\omega _l\omega _l^{})t}0|a_la_k^{}a_k^{}a_l^{}^{}|0$$
(29)
Using the commutation relation (20) we integrate (29) over $`l`$ and $`l^{}`$ and then insert the result into (28). Taking into account the positivity of energy (19) and the form of the function $`\varphi _{x_0,b}(k)`$ in (22) we obtain
$$\mathrm{\Phi }_{x_0,b}(t)|\rho (x)|\mathrm{\Phi }_{x_0,b}(t)=\left|\frac{1}{2\pi (2b)^{1/2}}_{x_0+b}^{x_0b}𝑑x^{}_{\mathrm{}}^+\mathrm{}𝑑ke^{i|k|t+ik(xx^{})}\right|^2$$
(30)
By comparison with (7) we see that this quantity is equal to the absolute value squared of the classical function $`\mathrm{\Phi }(x,t)`$ up to the normalization constant. Our discussion of nonlocality remains, therefore, also valid in the quantum case, and the expression inside the absolute value in (30) is a superposition of two nonlocal wave packets that move in opposite directions at the speed of light. In the appendix give a second example using an analogy with Fermi’s problem .
Let us note that in quantum field theory localization depends on the observable. If a state is local from the point of view of one observable, it can be nonlocal from the point of view of another. In our example the expectation value of the operator $`\rho (x)`$ is local in the state $`|\mathrm{\Phi }_{x_0,b}(0)`$. At the same time, the energy density of the field in the same state is nonlocal. Indeed, the energy density $`T_{00}(x)`$ of the free massless field (18) is
$$T_{00}(x)=\frac{1}{2}\left(\left(\frac{\widehat{\psi }}{t}\right)^2+\left(\frac{\widehat{\psi }}{x}\right)^2\right)$$
(31)
It contains the derivatives of the field operator. (As we have seen in classical case, the time derivative of the function $`\mathrm{\Phi }(x,t)`$ is nonlocal even at $`t=0`$). To determine the expectation value of $`T_{00}(x)`$ in the state $`|\mathrm{\Phi }_{x_0,b}(0)`$ taking into account the positivity of energy (19) we first calculate this expectation value for a finite $`t`$ and then take the limit $`t0`$. Using (18)-(23) obtain:
$$\mathrm{\Phi }_{x_0,b}(0)|T_{00}(x)|\mathrm{\Phi }_{x_0,b}(0)=\frac{1}{4\pi b}\left(\frac{1}{|xb|}+\frac{1}{|x+b|}\frac{1\mathrm{sign}(xb)\mathrm{sign}(x+b)}{\sqrt{|xb|}\sqrt{|x+b|}}\right)$$
(32)
where we put $`x_0=0`$ to simplify the expression. This quantity is obviously nonlocal.
## 4 Conclusion
Positivity of energy for quantum fields (or frequency for classical fields) leads to a decomposition of localized wave packets in terms of nonlocal wave packets with long tails. The long tails, which cancel each other initially, appear immediately afterwards as the nonlocal wavepackets move in opposite directions. In our examples the long tails decay with the distance $`x`$ according to $`b/x`$ for $`b/x1`$, where $`b`$ is the size of the localized wave packet. They are precursors to the usual wave propagation. Although we may have instant interactions, these are not the result of superluminal propagation, but of “preformed” structures.
We shall study the interaction between nonlocal structures in a separate paper. We have then “contact interactions”, due to the overlapping of the long tails. We shall also show that the photon clouds around atoms and molecules are nonlocal, which leads to the precursor effect and eliminates the apparent deviation from causality in Fermi’s two-atom problem. However, it is true that the two atoms “feel” each other instantaneously. Even inside a relativistic theory (the wave equation is Lorenz invariant) there is place for instantaneous interactions due to nonlocality.
Einstein’s relativistic events are associated to four dimensional points. Here we see nonlocal but still relativistic events that are due to the instability of localization as shown in the examples presented in this paper.
ACKNOWLEDGEMENTS
The authors would like to thank Prof. B. Pavlov and Prof. I. Antoniou for comments and fruitful discussions. This work was carried out with financial support of the International Solvay Institutes, the European Commission ESPRIT Project 28890 NTGONGS. This work was partially supported by the Engineering Research Program of the Office of Basic Energy Sciences at the Department of Energy Grant No. DE-FG03-94ER14465, the Robert A Welch Foundation Grant No. F-0365.
APPENDIX
Here we show an analogy of our problem with Fermi’s two-atom problem. We prepare our wave packet $`|\mathrm{\Phi }_{x_0,b}(t)`$ “localized” at $`t=0`$. At time $`t`$ we project this state on an second wave packet $`|\mathrm{\Phi }_{x_1,b}(0)`$, which is localized in a domain with center $`x_1x_0`$ and width $`2b`$, and plays the role of a measurement device. We choose $`x_1`$ so that $`x_1x_0>2b`$. Then, at $`t=0`$ these two states do not overlap and the scalar product $`\mathrm{\Phi }_{x_1,b}(0)|\mathrm{\Phi }_{x_0,b}(t)`$ vanishes. We consider this scalar product as a function of $`t`$, i.e., at each moment $`t`$ we project the evolving wave packet $`|\mathrm{\Phi }_{x_0,b}(x,t)`$ on the localized state $`|\mathrm{\Phi }_{x_1,b}(0)`$. As we have shown, the initially localized packet, which evolves in time, is nonlocal immediately after $`t=0`$ and our scalar product has a non-vanishing value. This can be interpreted as the second (localized) wave packet “feeling” the existence of the first one even at $`t<x_1x_02b`$ when the causal component of the first wave packet still did not reach the domain of localization of the second wave packet. The scalar product $`\mathrm{\Phi }_{x_1,b}(0)|\mathrm{\Phi }_{x_0,b}(t)`$ grows as the overlapping of the two wave packets increases (see Fig. 4). We expect some essential change of this growth when the main part of the first wave packet corresponding to the position of its local (“causal”) component reaches the domain of localization of the second wave packet. Then, as the moving component goes away, the scalar product decreases. To show this, we write this scalar product using (21) in the following form
$$\mathrm{\Phi }_{x_1,b}(0)|\mathrm{\Phi }_{x_0,b}(t)=_{\mathrm{}}^+\mathrm{}𝑑k𝑑k^{}\varphi _{x_1,b}^{}(k)\varphi _{x_0,b}(k^{})e^{i\omega _kt}0|a_ka_k^{}^{}|0$$
$`(A1)`$
We perform the integration over $`k^{}`$ with the help of the commutation relation (20). Then, inserting (22) and (23) we obtain
$$\mathrm{\Phi }_{x_1,b}(x,0)|\mathrm{\Phi }_{x_0,b}(x,t)=\frac{1}{4\pi b}_{x_1b}^{x_1+b}𝑑x^{}_{x_0b}^{x_0+b}𝑑x^{\prime \prime }_{\mathrm{}}^+\mathrm{}𝑑ke^{i|k|t+ik(x^{}x^{\prime \prime })}$$
$`(A2)`$
Using a regularisation similar to (7), we integrate over $`k`$ and obtain
$$\mathrm{\Phi }_{x_1,b}(0)|\mathrm{\Phi }_{x_0,b}(t)=\frac{1}{4\pi ib}_{x_1b}^{x_1+b}𝑑x^{}_{x_0b}^{x_0+b}𝑑x^{\prime \prime }\left(\frac{1}{tx^{}+x^{\prime \prime }i0}\frac{1}{t+x^{}x^{\prime \prime }i0}\right)$$
$`(A3)`$
The integration over $`x^{}`$ and $`x^{\prime \prime }`$ and the rearrangement of terms give us
$$\mathrm{\Phi }_{x_1,b}(0)|\mathrm{\Phi }_{x_0,b}(t)=\psi _1(x_1x_0+ti0)\psi _1(x_1x_0t+i0)$$
$`(A4)`$
where
$$\psi _1(x)=\frac{1}{4\pi ib}\left((x2b)\mathrm{ln}(x2b)+(x+2b)\mathrm{ln}(x+2b)2x\mathrm{ln}(x)\right)$$
$`(A5)`$
Then, using (12), (13) we come to
$$\mathrm{\Phi }_{x_1,b}(0)|\mathrm{\Phi }_{x_0,b}(t)=\psi _2(x_1x_0t)+\psi _2^{}(x_1x_0+t)$$
$`(A6)`$
where
$`\psi _2(x)`$ $`=`$ $`{\displaystyle \frac{1}{8b}}\left(|x2b|+|x+2b|2|x|\right)`$
$`+`$ $`{\displaystyle \frac{i}{4\pi b}}\left(|x2b|\mathrm{ln}|x2b|+|x+2b|\mathrm{ln}|x+2b|2|x|\mathrm{ln}|x|\right)`$
Fig. 5 shows the real part, the imaginary part and the absolute value of the scalar product $`\mathrm{\Phi }_{x_1,b}(0)|\mathrm{\Phi }_{x_0,b}(t)`$ as functions of $`t`$ for two different values of the wave packet’s width $`b`$. The real component is non-vanishing only in the time interval corresponding to the overlapping of the localized components. In contrast, the imaginary part is non-vanishing immediately after $`t=0`$, because it reflects the overlapping of nonlocal components. Fig. 5 also shows that the “causal” part of the effect is much bigger than the contribution of the long tails, if the size of the wave packet is much less than the distance between the domains of localization.
| 1 | The real part (dashed lines) and the imaginary part (solid lines) | |
| --- | --- | --- |
| | of $`\psi (xt)`$ (a), $`\psi ^{}(x+t)`$ (c) and $`\mathrm{\Phi }(x,t)=\psi (xt)+\psi ^{}(x+t)`$ (e) | |
| | as functions of $`x`$ at $`t=0`$; the absolute values $`|\psi (xt)|`$ (b), $`|\psi ^{}(x+t)|`$ (d) | |
| | and $`|\mathrm{\Phi }(x,t)||\psi (xt)|+|\psi ^{}(x+t)|`$ (f) as functions of $`x`$ at $`t=0`$. | |
| 2 | The real part (dashed lines) and the imaginary part (solid lines) | |
| | of $`\psi (xt)`$ (a), $`\psi ^{}(x+t)`$ (c) and $`\mathrm{\Phi }(x,t)=\psi (xt)+\psi ^{}(x+t)`$ (e) | |
| | as functions of $`x`$ at $`t=0.25s`$; the absolute values $`|\psi (xt)|`$ (b), $`|\psi ^{}(x+t)|`$ (d) | |
| | and $`|\mathrm{\Phi }(x,t)||\psi (xt)|+|\psi ^{}(x+t)|`$ (f) as functions of $`x`$ at $`t=0.25s`$. | |
| 3 | The real part (dashed lines) and the imaginary part (solid lines) of $`\psi (xt)`$ (a), | |
| | $`\psi ^{}(x+t)`$ (c) and $`\mathrm{\Phi }(x,t)=\psi (xt)+\psi ^{}(x+t)`$ (e) as functions of $`x`$ | |
| | at $`t=1s`$; the absolute values $`|\psi (xt)|`$ (b), $`|\psi ^{}(x+t)|`$ (d) and | |
| | $`|\mathrm{\Phi }(x,t)||\psi (xt)|+|\psi ^{}(x+t)|`$ (f) as functions of $`x`$ at $`t=1s`$. | |
| 4 | $`\mathrm{\Phi }_{0,b}(t)|\rho (x)|\mathrm{\Phi }_{0,b}(t)`$ (evolving object) and $`\mathrm{\Phi }_{2,b}(0)|\rho (x)|\mathrm{\Phi }_{2,b}(0)`$ | |
| | (right rectangle) with no overlap at $`t=0`$ (a), overlapping only by the | |
| | nonlocal tail at $`t=0.75s`$ (b), overlapping also by the local (causal) | |
| | component at $`t=1.25s`$ (c). | |
| 5 | The real part (a) and (b), the imaginary part (c) and (d), the absolute value | |
| | (e) and (f) of $`\mathrm{\Phi }_{2,b}(0)|\mathrm{\Phi }_{0,b}(t)`$ as functions of $`t`$ for $`b=0.5ls`$ ((a), (c), (e)) and | |
| | $`b=0.01ls`$ ((b), (d), (f)). | |
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# Acknowledgments
## Acknowledgments
This project is supported by the Ministry of Science and Technology of the Republic of Slovenia and by the Rector’s Fund of the University of Maribor. VR acknowledges the support of the work by the grant of the Ministry of Science and Technology of the Republic of Slovenia and the Abdus Salam ICTP (Trieste) Joint Programme.
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# Generating moment equations in the Doi model of liquid–crystalline polymers
## I Introduction
Kinetic theory is a powerful analytical tool for describing the dynamics of dilute and semi–dilute solutions of polymers in terms of a diffusion equation for the particle distribution function or, equivalently, by the full system of moment equations. In general, the moment system has to be truncated at some level. The problem of the “closure approximation” is very well–known in the literature, and an enormous amount of suggestions have been analyzed in the case where each moment couples only to a few higher–order moments. However, for some kinetic equations, the time evolution of each moment couples to an infinite set of higher moments, so that further analytical work is often precluded since closure approximations are less studied for this case. In this paper we present a simple method that allows to derive moment equations with a finite coupling valid for a wide class of kinetic equations.
In order to be specific, we consider a particularly important example: the Doi theory of liquid–crystalline polymers (LCP), subject to the Onsager excluded–volume potential . As it is well–known, in this model each moment equation depends on an infinite set of higher–order moments. In the original work , this problem was treated in two steps: First, the Onsager potential was replaced by a different, phenomenological potential of the Maier–Saupe type , which gives raise to a coupling to the next higher moment only. In the second step, the “decoupling” approximation was used to solve the resulting closure problem for the second moment. Subsequent extensive studies were focused on improvements of the second step . At the same time, we are not aware of improvements on the first step and closure approximations are limited to the Maier–Saupe potential up to now. However, it would be desirable to deal with the true Onsager potential, not only because it becomes exact in the limit of low concentrations of perfectly rigid rod–like molecules, but also because it contains no phenomenological parameters and therefore gives more quantitative predictions. In addition, the Onsager potential is preferred in the study of the influence of flow on the isotropic–nematic transition, since it gives a clear–cut prediction of the range of coexistence of the equilibrium isotropic and nematic phase, both in stationary and non–stationary flows. The method, which we propose in this work, leads to an approximation of the Onsager potential, which, to the lowest order, is at the same time as simple as the Maier–Saupe potential but also closer to the true Onsager potential. Moreover, corrections to this approximation can be obtained in a systematic manner.
## II The Doi model
Let $`\psi (𝒖;t)`$ be the probability distribution function for a rigid rod–like polymer molecule to be oriented parallel to the unit vector $`𝒖`$. The time evolution of $`\psi `$ in the presence of flow and the Onsager excluded–volume potential was given by Doi and may be written as:
$$_t\psi =\left[𝒖\times \left(𝜿𝒖\psi \right)\right]+\widehat{D}_\mathrm{r}\psi \left(\frac{\delta A}{\delta \psi (𝒖)}\right).$$
(1)
Here $`=𝒖\times /𝒖`$ is the rotational operator, $`/𝒖`$ the gradient on the unit sphere, $`𝜿`$ the gradient of the velocity, $`\widehat{D}_\mathrm{r}`$ the rotational diffusivity, $`\delta /\delta \psi `$ the functional derivative and $`A=A_0+A_1`$ is the free energy functional per molecule divided by $`k_\mathrm{B}T`$,
$`A_0`$ $`=`$ $`\mathrm{ln}\nu 1+\mathrm{ln}\psi (𝒖)`$ (3)
$`A_1`$ $`=`$ $`{\displaystyle \frac{U}{2}}\sqrt{1(𝒖𝒘)^2}.`$ (4)
$`U=2bL^2\nu `$ is the reduced excluded–volume, $`2b`$ and $`L`$ are the diameter and the length of the rod–like polymeric molecules, respectively, and $`\nu `$ is the number of molecules per unit volume. Here and below we use the following notations for averages: $`f(𝒖)=f(𝒖)\psi (𝒖)𝑑𝒖`$, and $`f(𝒖,𝒘)=f(𝒖,𝒘)\psi (𝒖)\psi (𝒘)𝑑𝒖𝑑𝒘`$. $`A_0`$ describes the loss of entropy with molecular alignment, while $`A_1`$ expresses the Onsager free energy of steric interaction in the second virial approximation . Following Doi and Edwards , the rotational diffusivity is approximated by
$$\widehat{D}_\mathrm{r}\overline{D}_\mathrm{r}=D_\mathrm{r}\left[\frac{4}{\pi }\sqrt{1(𝒖𝒘)^2}\right]^2,$$
(5)
where $`D_\mathrm{r}`$, the rotational diffusion coefficient for a rod in an isotropic, semi–dilute solution of like rods, is related to the rotational diffusion constant for a dilute solution, $`D_{\mathrm{r0}}`$, by $`D_\mathrm{r}=cD_{\mathrm{r0}}(\nu L)^2`$ with an empirical coefficient $`c`$. Nonlinearity of Eq. (1) in $`\psi `$ brought about by the potential (4) reflects the mean–field nature of the Onsager theory of the excluded–volume effect. The self–consistent potential, identified by Doi , is related to the free energy of interaction, $`V(𝒖)=k_\mathrm{B}T\delta A_1/\delta \psi (𝒖)`$. Various phases of the LCP are conveniently described by the order parameter $`\mathrm{SS}=𝒖𝒖(1/3)\mathrm{𝟏}`$, where $`\mathrm{𝟏}`$ is the unit tensor. It is reasonable therefore to look for approximate formulations of the dynamics in terms of the order parameter alone. However, as mentioned above, the time evolution equation for $`\mathrm{SS}`$ couples to an infinite number of moments of $`\psi `$. In the derivation given by Doi, this difficulty was circumvented by replacing the Onsager potential (4) by a different, phenomenological expression of the Maier–Saupe type :
$$A_1^{\mathrm{MS}}=a_0\frac{a_1}{2}U\mathrm{SS}:\mathrm{SS},$$
(6)
where $`a_0`$ and $`a_1`$ are parameters independent of $`\psi `$. A further separate treatment of the diffusivity (5) is also necessary. A compact presentation of the entire development is given by Doi and Edwards . The Doi model with the Maier–Saupe potential (6) constitutes the basic kinetic model of LCP used by many authors for analytical studies to derive equations for the order parameter. As is well–known, the kinetic equation (1) with the potential (6) does not give a closed equation for the order parameter but contains also the higher–order moment $`𝒖𝒖𝒖𝒖`$, and therefore constitutes a further problem of closure. The original Doi approach was based on the decoupling approximation for the fourth–order moments of $`\psi `$ in terms of $`\mathrm{SS}`$. Improvements on the decoupling approximation are currently under active research .
## III Generating moment equations
In this communication, we demonstrate that a different self–consistent treatment of the kinetic equation (1) is possible. Modifications concern only the relaxational part of the Eq. (1), specifically, the excluded–volume potential (4) and the diffusivity (5), and therefore we consider the case $`𝜿=0`$ in the sequel to simplify notations. Specifically, we employ the cumulant expansion of the potential (4) and the diffusivity (5). The leading term of this expansion results in an effective potential that differs from the Maier–Saupe potential (6), and which contains a non-polynomial dependence on the order parameter $`\mathrm{SS}`$.
In the second virial approximation, the free energy of interaction, $`A_1`$ can be written as $`A_1=(\nu /2)\beta (𝒖,𝒘)`$. If only excluded–volume interactions are present, the second virial coefficient $`\beta `$ corresponding to the Onsager expression (4) is of the form $`\beta (𝒖,𝒘)=\beta ((𝒖𝒘)^2)`$, with $`\beta (x)=2bL^2\sqrt{1x}`$. Specifically, expanding $`\beta (x)`$ in a Taylor series and interchanging summation and averaging in this expansion, we get $`\beta (x)=_{n=0}^{\mathrm{}}a_nx^n`$, where $`a_n`$ are numerical coefficients. Each average $`x^n`$ can be represented in terms of cumulants $`x^k_c`$ of order $`kn`$. Resummation of the series leads to
$$\beta (x)=\beta (x)+\underset{m=1}{\overset{\mathrm{}}{}}\frac{1}{m!}\left(\frac{x^2_c}{2}\right)^m\beta ^{(2m)}(x)+\mathrm{},$$
(7)
where $`\beta ^{(2m)}`$ is the $`2m`$–th derivative of $`\beta `$ and ellipses denote terms including third or higher–order cumulants as factors. Therefore, the functional $`A_1`$ (4) can be split as $`A_1=A_1^{(1)}+B`$, where $`A_1^{(1)}`$ is the free energy, corresponding to the total neglect of second and higher–order cumulants in each term of the expansion,
$$A_1^{(1)}=\frac{U}{2}\sqrt{1𝒖𝒖:𝒖𝒖}.$$
(8)
In terms of the order parameter $`\mathrm{SS}`$, $`A_1^{(1)}`$ may be rewritten as $`A_1^{(1)}=(U/\sqrt{6})\sqrt{1(3/2)\mathrm{SS}:\mathrm{SS}}`$. By the mean value theorem, it is easy to see that $`A_1^{(1)}`$ gives an upper bound to $`A_1`$, $`A_1A_1^{(1)}`$, for the present case of excluded–volume interactions. The functional $`B`$ contains the higher–order cumulants. While all powers of the second cumulants are displayed in Eq. (7), in general it is not a priori clear whether it is more important to keep powers of the second cumulant or higher cumulants. However, we generally expect the linear term in the second order cumulant to be most important. The corresponding term $`m=1`$ in Eq. (7), $`A_1^{(2)}=(1/2)x^2_c\beta ^{\prime \prime }(x)`$, gives the first non-vanishing contribution to $`B`$,
$$A_1^{(2)}=\frac{U}{16}[(𝒖𝒘)^2𝒖𝒖:𝒖𝒖]^2(1𝒖𝒖:𝒖𝒖)^{3/2}.$$
(9)
Keeping only the first $`N`$ cumulants in the expansion (7), the functional (4) is approximated by non–linear functions of the first $`2N`$ moments of $`\psi `$. Inserting this approximation in the time evolution equation (1) amounts to an approximation of the self–consistent Onsager potential $`V`$ in terms of a polynomial of order $`2N`$ in $`𝒖`$ but with non–linear coefficients. In this approximation, the time evolution of the $`2n`$–th moment contains only the first $`(2n+2N)`$–th moments. With this, moment equations can be generated, that approach the original equations in a systematic way, thereby containing only a finite number of moments at each stage.
## IV Testing the approximation
Clearly, the above procedure is most valuable if the first terms, $`A_1^{(1)}`$, etc. , already provide a good approximation to the full expression $`A_1`$. While the general validity of $`A_1^{(1)}`$ as a good approximation to $`A_1`$ is a rather delicate problem, it should be mentioned that it is so at least in two limiting cases. Namely, for the isotropic state, the value of $`A_1^{(1)}`$ differs from $`A_1`$ for less than 5%, while in the fully ordered state the approximation (8) becomes exact. Moreover, on the submanifold of distribution functions of the form
$$\psi _\alpha (𝒖)=\frac{\alpha }{4\pi \mathrm{sinh}\alpha }\mathrm{cosh}(\alpha 𝒖𝒏),$$
(10)
where $`𝒏`$ is an arbitrary unit vector, and $`0\alpha \mathrm{}`$, the functional $`A_1^{(1)}`$ turns out to approximate $`A_1`$ very well for all values of the parameter $`\alpha `$ between the isotropic state, $`\alpha =0`$, and the fully ordered state, $`\alpha =\mathrm{}`$. To show this, we plot in Fig. 1 the functions $`A_1(\alpha )`$, $`A_1^{(1)}(\alpha )`$ and $`A_1^{(1)}(\alpha )+A_1^{(2)}(\alpha )`$, that result upon inserting the ansatz (10) into (4), (8), (9), respectively. Note, that $`A_1(\alpha )`$, $`A_1^{(1)}(\alpha )`$ and $`A_1^{(2)}(\alpha )`$ can be calculated analytically. For convenience, we plot the functions against the scalar order parameter, defined as $`S=\sqrt{(3/2)\mathrm{SS}:\mathrm{SS}}`$. Including $`A_1^{(2)}`$ does not only reduce the error of the approximate value of $`A_1`$ in the isotropic state to 1.5%, but improves the accuracy of the approximation over the whole range of $`S`$. For comparison, we included in Fig. 1 also the free energy $`A_1^{\mathrm{MS}}`$, corresponding to the Maier–Saupe expression (6), thereby choosing the undetermined constant so that the limit of the fully ordered state is matched correctly. Note, however, that in any case $`A_1^{\mathrm{MS}}`$ decays asymptotically like $`1/\alpha `$, for $`\alpha 1`$, whereas $`A_1`$ and $`A_1^{(1)}`$ behave like $`1/\sqrt{\alpha }`$ in this regime. We included in Fig. 1 also the derivative of the above functions, since they are related to the self–consistent potential $`V`$. Fig. 1 shows that also the derivative of $`A_1^{(1)}`$ provides a good approximation to the derivative of $`A_1`$, with correct limiting behavior near the isotropic and fully ordered state. Note that including the first correction $`A_1^{(2)}`$ yields excellent agreement to the true Onsager prediction. The Maier–Saupe potential captures the main features but, besides an undetermined constant, shows the wrong behavior near the fully ordered state. The ansatz (10), originally proposed by Onsager , is known to approximate the equilibrium distribution very well. Therefore we conclude that $`A_1^{(1)}`$ represents a good approximation to $`A_1`$, at least on a representative subset of distribution functions.
## V Thermodynamic consistency
It is worth mentioning that the presentation given so far can easily be cast into the recently developed GENERIC formalism of nonequilibrium thermodynamics . In the absence of potential forces, the example of rigid dumbbells, which are equivalent to the model of rigid rods, is formulated within the GENERIC formalism in Ref. . The mean field potentials considered above can be included in a straightforward manner, if we recognize that $`A_1=S_1`$, where $`S_1`$ is the entropic contribution per molecule to the free energy of interaction divided by $`k_\mathrm{B}`$. Formulating the original model as well as the approximations within the GENERIC formalism guarantees that our treatment is in accordance with the principles of nonequilibrium thermodynamics. This becomes especially important if the present model is considered in nonisothermal situations. For example, the structure of the GENERIC formalism requires the polymeric contribution to the elastic stress to be
$$\sigma _{\alpha \beta }^\mathrm{e}=3\nu k_\mathrm{B}TS_{\alpha \beta }\nu k_\mathrm{B}T(𝒖\times \frac{\delta S_1}{\delta \psi })_\alpha u_\beta .$$
(11)
Eq. (11) agrees with the result of Doi , obtained upon varying the free energy functional.
## VI The lowest order approximation
In the sequel, we will adopt the lowest order approximation $`A=A^{(1)}=A_0+A_1^{(1)}`$, where $`A_1^{(1)}`$ is given by Eq. (8), and $`A_0`$ is given by Eq. (3). This amounts to neglect of all higher order correlations in Eq. (7), or, equivalently, setting $`B=0`$. Substituting $`A^{(1)}`$ instead of $`A`$ into Eq. (1), we derive
$$_t\psi =\widehat{D}_\mathrm{r}\left[\psi \psi \left(\frac{U𝒖𝒖:𝒖𝒖}{2\sqrt{1𝒖𝒖:𝒖𝒖}}\right)\right].$$
(12)
It is now possible to identify the self–consistent potential as
$$V^{(1)}(𝒖)=\left(\frac{Uk_\mathrm{B}T}{2}\right)\frac{1𝒖𝒖:𝒖𝒖}{\sqrt{1𝒖𝒖:𝒖𝒖}},$$
(13)
which can be compared to the expression obtained from inserting the Maier–Saupe free energy (6) into Eq. (1)
$$V_{\mathrm{MS}}(𝒖)=a_2a_1Uk_\mathrm{B}T𝒖𝒖:𝒖𝒖,$$
(14)
where $`a_2`$ is an arbitrary constant. The normalized equilibrium solutions to the Eq. (12) are $`\psi _{\mathrm{eq}}^{(1)}=Z^1\mathrm{exp}[V^{(1)}/k_\mathrm{B}T]`$.
The rotational diffusivity (5) is related to the free energy of interaction, since $`\overline{D}_\mathrm{r}=D_\mathrm{r}\left[\frac{4}{\pi }A_1/(U/2)\right]^2`$. Substituting $`A_1=A_1^{(1)}`$ gives
$$\overline{D}_\mathrm{r}^{(1)}=(3\pi ^2/32)D_\mathrm{r}[1(3/2)\mathrm{SS}:\mathrm{SS}]^1.$$
(15)
The diffusion coefficient $`\overline{D}_\mathrm{r}^{(1)}`$ (15) is positive in the entire physically meaningful range of the order parameter $`\mathrm{SS}`$. Expression (15) should be compared with the Doi phenomenological result:
$$\overline{D}_{\mathrm{rD}}=D_\mathrm{r}[1(3/2)\mathrm{SS}:\mathrm{SS}]^2.$$
(16)
While we have not found an argument which of the two powers, $`1`$ or $`2`$, is more consistent, it should be stressed that our derivation of the diffusion coefficient does not need any further assumptions or adjustable parameters, while the derivation of Eq. (16) requires the matching of $`\overline{D}_\mathrm{r}`$, resp. $`A_1^{\mathrm{MS}}`$ in both, the isotropic and the fully ordered state. Due to its relation to $`A_1^{(1)}`$, the diffusion coefficient $`\overline{D}_\mathrm{r}^{(1)}`$ (15) has a correct limit in the fully ordered state ($`D_\mathrm{r}/\overline{D}_\mathrm{r}^{(1)}=0`$ as soon as $`\mathrm{SS}:\mathrm{SS}=2/3`$ in the ordered state), while the opposite limit of the isotropic state ($`\overline{D}_\mathrm{r}=D_\mathrm{r}`$) is matched within $`8\%`$. Again, the first correction (9) reduces the error in this limit to less than $`3\%`$.
If we adopt (5) and approximate $`\overline{D}_\mathrm{r}`$ by $`\overline{D}_\mathrm{r}^{(1)}`$ (15), the time evolution of the order parameter $`\mathrm{SS}`$ can be derived from Eq. (12) by the so–called Prager procedure
$$_t\mathrm{SS}=6\overline{D}_\mathrm{r}^{(1)}\mathrm{SS}+6\overline{D}_\mathrm{r}^{(1)}\frac{U^{}}{\sqrt{1(3/2)\mathrm{SS}:\mathrm{SS}}}(\mathrm{SS}𝒖𝒖\mathrm{SS}:𝒖𝒖𝒖𝒖),$$
(17)
with $`U^{}=U/\sqrt{6}`$. This expression differs from the result of Doi not only in the diffusion coefficient and in the reduced excluded–volume $`U`$ due to the undetermined constant in the Maier–Saupe potential (14), but contains a non–polynomial dependence on the order parameter $`\mathrm{SS}`$, which becomes important in the nematic state.
## VII Conclusion
We have presented a systematic procedure that allows to derive approximate moment equations for the Doi model of LCP, which contain only a finite number of higher order moments. The first approximation for the Onsager excluded–volume interaction results in an effective potential (13) proportional to $`𝒖`$$`𝒖`$, but different from the Maier–Saupe form (14) and without free parameters. We find indications, that (13) approximates the true Onsager potential better than the Maier–Saupe potential. For higher accuracy, the first correction seems to be the most important contribution. All these approximations are in accordance with nonequilibrium thermodynamics.
Note, that we have not addressed the problem of solving the resulting kinetic equations or “closing” the moment equations. This work is currently under preparation. Nevertheless, for comparing Eq. (17) to the corresponding equation with the Maier–Saupe potential, we follow Refs. and consider the decoupling approximation $`\mathrm{SS}:𝒖𝒖𝒖𝒖=\mathrm{SS}:𝒖𝒖𝒖𝒖`$. If the order parameter is assumed to be of the form $`S_{\alpha \beta }=S(t)[n_\alpha n_\beta \delta _{\alpha \beta }]`$, the relaxation equation for the scalar parameter $`S`$ is found to be as follows:
$`_tS`$ $`=`$ $`6\overline{D}_\mathrm{r}^{(1)}{\displaystyle \frac{A^{(1)}(S,U^{})}{S}},`$ (19)
$`A^{(1)}(S,U^{})`$ $`=`$ $`{\displaystyle \frac{S^2}{2}}{\displaystyle \frac{U^{}}{9}}\sqrt{1S^2}\left(1{\displaystyle \frac{3S}{2}}+2S^2\right){\displaystyle \frac{U^{}}{6}}\mathrm{arcsin}(S),`$ (20)
where $`U^{}=U/\sqrt{6}`$. Due to a non–polynomial character of the function $`A^{(1)}`$ (20), the relaxation equation (19) differs formally from the Landau–de Gennes counterpart derived by Doi for the Maier–Saupe potential. Expansion of the function (20) around $`S=0`$ reproduces the result of Doi for the Maier–Saupe potential, subject to a renormalization of the strength of the excluded–volume potential, and a difference in the coefficient in front of the $`S^4`$ term. Moreover, the relaxation implied by $`A^{(1)}`$ (20) is qualitatively similar to the one given by Doi result and distinguishes the same three regimes. For $`U<U_1`$, $`A^{(1)}`$ has only one minimum at $`S=0`$, so that the system finally becomes isotropic. For $`U_1<U<U_2`$, a second local minimum occurs. The system either becomes isotropic or nematic depending on the initial value of $`S`$. Finally, for $`U>U_2`$, the isotropic state becomes unstable and the system always approaches a nematic state.
Due to the undetermined constant $`a_1`$ in the Maier–Saupe potential (14), the Doi theory predicts the values $`U_1`$ and $`U_2`$ also in terms of $`a_1`$. On the contrary, the self–consistent potential (13) contains no free parameters, so that $`U_1=3^{1/4}\sqrt{8}/(\sqrt{3}1)6.22`$ and $`U_2=3\sqrt{6}7.34`$ may directly be compared to the values in the Onsager theory $`U_1=8.38`$ and $`U_2=10.67`$. However, the values of the order parameter $`S_1=1/4`$ and $`S_2=1/2`$ at $`U_1`$ resp. $`U_2`$ in the Doi theory do not depend on $`a_1`$ and can therefore be compared to the values predicted by (13): $`S_1=(\sqrt{3}1)/20.37`$ and $`S_2=1/\sqrt{2}0.71`$.
In Fig. 2, the equilibrium order parameter $`S_{\mathrm{eq}}`$ is shown as a function of $`\nu /\nu _2`$. The lower solid line shows the prediction of the Doi theory, whereas the upper solid line corresponds to the approximation (13). For the Maier–Saupe potential, $`S_{\mathrm{eq}}0`$ is given by
$$S_{\mathrm{eq}}=\frac{1}{4}+\frac{3}{4}\sqrt{1\frac{8\nu _2}{3\nu }}.$$
(21)
For the free energy (20), $`S_{\mathrm{eq}}`$ is given implicitly as the solution to the algebraic equation
$$\frac{1+S_{\mathrm{eq}}2S_{\mathrm{eq}}^2}{\sqrt{1S_{\mathrm{eq}}^2}}=\frac{\nu _2}{\nu }.$$
(22)
For large values of $`\nu /\nu _2`$, the solution (22) approaches the value $`S_{\mathrm{eq}}=1`$ and asymptotically behaves like the solution of the Doi theory $`S_{\mathrm{eq}}1\nu _2/\nu `$, for large $`\nu `$. Note, that the decoupling approximation corrects the asymptotic behavior of the Maier–Saupe potential near the fully ordered state.
As is well–known, the detailed form of the interaction potential can have significant effect on the behavior of the order parameter in the nematic phase . Specifically, the amount of order at the transition is known to be much smaller in the Maier–Saupe theory than in the Onsager model. For comparison, we included in Fig. 2 the values of the order parameter obtained from minimizing the true Onsager free energy numerically , where $`\nu _2`$ now corresponds to the true nematic transition. Although the analysis of the phase transitions via the dynamical approach is affected by the use of the decoupling approximation, the prediction of the self–consistent approach is much closer to the true Onsager values than is the Maier–Saupe potential.
Finally, it should be mentioned that approximations to the Onsager potential like Eqs. (8) and (9) can also be used in the case of potential flows, following the approach of Thirumalai without additional assumptions.
To summarize, we have developed a direct approach to the Doi model with the Onsager potential. We have demonstrated that the resulting kinetic equation has much in common with the Doi model with the phenomenological Maier–Saupe potential. Corrections to the approximation developed here can be found in a systematic way from Eqs. (7) to (9) by taking into account higher order correlations. The approach to derive self–consistent moment equations is applicable to other kinetic equations which can be cast into the form (1).
## Figure captions
Free energy of excluded–volume interaction for distribution functions (10) plotted against the scalar order parameter $`S(\alpha )`$. From top to bottom: approximation $`A_1^{(1)}`$ (8), with first correction $`A_1^{(1)}+A_1^{(2)}`$ (9), true Onsager expression $`A_1`$ (4) and the Maier–Saupe free energy $`A_1^{\mathrm{MS}}`$ (6), when the limits $`S=0`$ and $`S=1`$ are matched. In the inset, the derivative of the above functions is shown as a function of $`S(\alpha )`$. The order of the curves from top to bottom is the same.
The equilibrium order parameter $`S_{\mathrm{eq}}`$ as a function of $`\nu /\nu _2`$. The figure shows the behavior due to the Maier–Saupe potential (21), lower curve, and the solution of (22) corresponding to $`A^{(1)}`$, upper curve, in the decoupling approximation. Full circles indicate the order parameter for the true Onsager potential in the static case (from ).
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# Quantum phase transitions in a linear ion trap
## I Introduction
More than two decades ago, when quantum optics was young, the quantum dynamics of collective spin systems interacting with a single bosonic degree of freedom was a major research problem. The model arose as an attempt to describe the interaction between a collection of two level atoms and a single mode of the radiation field. Walls and co workers were among the first to realise that such models provided ideal examples of the role of quantum fluctuations in the nonlinear interaction between matter and light. Quantum fluctuations were shown to drastically change the predictions of semiclassical theory in such systems. This phenomenon has appeared more recently in the discovery of quantum phase transitions in quantum spin glasses and other many body quantum systems. While the collective spin models did not directly apply to achievable experiments at the time, they did provide insight that subsequently proved important for many other quantum optical experiments including anti-bunching, squeezing, and cavity QED. In this paper we show that the models of a collective spin interacting with one or more bosonic modes can now be experimentally realised in modern ion trap systems of the kind proposed for quantum computation. An enormous effort has gone into making such systems work at the quantum level, with little interference form classical sources of noise, and a number of such experiments exist today. It would thus appear worthwhile to reconsider the collective spin models, and the associated quantum many-body effects exhibited by such systems, with a view to direct experimental realisation.
In particular we consider the Tavis-Cummings (TC) model, which can be realised in a linear ion trap of $`N`$ ions with the bosonic degree of freedom appearing as the quantised collective centre-of-mass motion. If each ion is coupled to the vibrational motion using an identical external (classical) laser detuned to the first red-sideband transition, the symmetry is such that the electronic degree of freedom for the ions can be described as a collective spin ($`N`$) and the reversible dynamics is well described by the TC model. The TC model is known to exhibit important nonlinear quantum effects including a quantum phase transition in which the (zero temperature) ground state undergoes a morphological change as a parameter is varied and averages of intensive quantities undergo a bifurcation.
## II The Tavis-Cummings model
The interaction Hamiltonian for N ions interacting with the centre of mass vibrational mode can be controlled by using different kinds of Raman laser pulses. A considerable variety of interactions has already been achieved or proposed . In this paper we consider the first red-sideband transition. The ion is assumed to be in a three dimensional anisotropic harmonic potential. Two dimensions are very tightly bound and are neglected. In the remaining dimension, an external laser couples the electronic state to the vibrational motion. If the vibrational frequency is large enough and the Lamb-Dicke limit applies the motional sidebands of the absorption of the electronic transition can be resolved and a laser detuned below the electronic resonance by one unit of the trap frequency can excite the electronic transition by absorbing one vibrational phonon, the additional energy required being made up by the laser. We will assume that the laser ( or lasers if a Raman process is used) is sufficeintly strong that it can be treated classically. Under these assumptions the Hamiltonian, in the interaction picture, is
$$H_I=\mathrm{}\mathrm{\Omega }\underset{i=1}{\overset{N}{}}(a\sigma _+^{(i)}+a^{}\sigma _{}^{(i)})$$
(1)
where the coupling constant is $`\mathrm{\Omega }=\eta \mathrm{\Omega }_0`$ where $`\eta ^2=E_r/(\mathrm{}M\omega _0)`$ is the Lambe-Dicke parameter with $`E_r`$ the recoil kinetic energy of the atom, $`\omega _0`$ is the trap vibrational frequency, and $`M`$ is the effective mass for the centre-of-mass mode. The Lamb-Dicke limit assumes $`\eta <<1`$, which is easily achieved in practice. The frequency, $`\mathrm{\Omega }_0`$ is the effective Rabi frequency for the electronic transition involved. The raising and lowering operators for each ion are defined by $`\sigma _{}=|ge|`$ and $`\sigma _+=|eg|`$. This sideband transition can be used to efficiently cool the ions to the collective centre-of-mass ground state, thus preparing the system in the vibrational ground state.
If the external laser field on each ion is identical (in amplitude and phase) the interaction Hamiltonian is
$$H_I=\mathrm{}\mathrm{\Omega }(a\widehat{J}_++a^{}\widehat{J}_{})$$
(2)
where we have introduced the bosonic annihilation operator $`a`$ for the centre-of-mass vibrational mode and where we have used the definition of the collective spin operators,
$$\widehat{J}_\alpha =\underset{i=1}{\overset{N}{}}\sigma _\alpha ^{(i)}$$
(3)
where $`\alpha =x,y,z`$ . Identical laser fields could easily be obtained by splitting a single, stabilised laser into multiple beams. The interaction Hamiltonian in Eq (2) specifies the Tavis-Cummings model. This model first appeared in quantum optics where the bosonic mode is the quantised field in a cavity. However this realisation is difficult to achieve experimentally. In contrast the vibrational mode realisation should be readily achieved. The dynamics resulting from this Hamiltonian is quite rich. Collective spin models of this kind were considered many decades ago in quantum optics. In much of that work however the collective spin underwent an irreversible decay. In the case of an ion trap model however we can neglect such decays due to the long lifetimes of the excited states. On the other hand heating of the vibrational centre-of-mass mode can induce irreversible dynamics in the system in a manner that has not been previously considered, and that is reminiscent of thermal effects in condensed matter physics.
We are interested in the driven Tavis-Cummings model in which the vibrational mode is subject to a linear forcing term which can easily be achieved by a suitable combination of Raman laser pulses, or by appropriate AC voltages applied to the trap electrodes. In this case the Hamiltonian, in the interaction picture, is given by
$$H_I=\mathrm{}\mathrm{\Omega }(a\widehat{J}_++a^{}\widehat{J}_{})+\mathrm{}E(a+a^{})$$
(4)
This may be written in terms of the hermitian canonical oscillator variables $`\widehat{X}=(a+a^{})/\sqrt{2}`$, $`\widehat{Y}=i(aa^{})/\sqrt{2}`$, and the canonical angular momentum variables $`\widehat{J}_x=(\widehat{J}_++\widehat{J}_{})/2`$, $`\widehat{J}_y=i(\widehat{J}_+\widehat{J}_{})/2`$, $`\widehat{J}_z=[\widehat{J}_+,\widehat{J}_{}]/2`$. It takes the form
$$H=\widehat{X}\widehat{J}_x\widehat{Y}\widehat{J}_y+\chi \widehat{X}$$
(5)
with $`\chi =E/\mathrm{\Omega }`$ and we have scaled the Hamiltonian by $`HH/\sqrt{2}\mathrm{\Omega }`$. This indicates that time is measured in units of $`\frac{1}{\sqrt{2}\mathrm{\Omega }}`$.
Alsing has shown that the ground state of this system, for weak driving, is a product state in which the bosonic mode is squeezed and the electronic states are rotated in the angular momentum space. We provide a direct proof of this statement below. However it is first useful to consider the dynamics of the equivalent semiclassical model as many of the results in the quantum case can be interpreted in terms of the features of the semiclassical model.
### A Semiclassical Tavis-Cummings model
The Tavis-Cummings model represents an interaction between a simple harmonic oscillator and a linear top for which there is a classical model which we now define. We choose the classical model so that the equations of motion are of the same form as the Heisenberg equations of motion for the quantum model. The classical Hamiltonian is defined as
$$=X𝒥_xY𝒥_y+\chi EX$$
(6)
where $`X,Y`$ are respectively the canonical oscillator position and momentum variables with the canonical Poisson bracket $`\{X,Y\}=1`$, while $`𝒥_k`$ are the three components of angular momentum for a classical top with the canonical Poisson brackets $`\{𝒥_i,𝒥_k\}=_kϵ_{ijk}𝒥_k`$. The equations of motion for a canonical coordinate $`w`$ is given as usual by Poisson bracket with the Hamiltonian $`\dot{w}=\{w,H\}`$. The equations of motion are,
$`\dot{X}`$ $`=`$ $`𝒥_Y`$ (7)
$`\dot{Y}`$ $`=`$ $`𝒥_X\chi `$ (8)
$`\dot{𝒥_x}`$ $`=`$ $`Y𝒥_z`$ (9)
$`\dot{𝒥_y}`$ $`=`$ $`X𝒥_z`$ (10)
$`\dot{𝒥_z}`$ $`=`$ $`X𝒥_y+Y𝒥_x`$ (11)
Note that these equations have a conservation law $`𝒥_x^2+𝒥_y^2+𝒥_z^2=\text{constant}`$.
We now justify this choice of classical Hamiltonian by noting that the Heisenberg equations of motion for the Hamiltonian Eq(4) have the same form as the semiclassical equations of motion with all variables replaced by the corresponding operators. We thus see that the semiclassical equations result form taking moments of the Heisenberg equations and factorising all product moments. The factorisation assumptions ignores correlations which scale as $`1/N`$ for the scaled operators $`\widehat{J}_\sigma /N`$. The conservation law $`𝒥_x^2+𝒥_y^2+𝒥_z=constant`$ is a reflection of the operator relation
$$\widehat{J}^2=\frac{N}{2}\left(\frac{N}{2}+1\right)$$
(12)
which in the semiclassical limit indicates that $`𝒥_x^2+𝒥_y^2+𝒥_z=\frac{N^2}{4}`$.
The classical equations have one nontrivial fixed point at $`X^{}=Y^{}=𝒥_y^{}=0`$ and $`𝒥_x^{}=\chi `$, $`𝒥_z^{}=\sqrt{N^2/4\chi ^2}`$. However as the conservation law requires that $`|𝒥_x|N/2`$ we see that we must have
$$\frac{2E}{N\mathrm{\Omega }}1\text{(below threshold)}$$
(13)
which corresponds to an energy of $`=0`$. We will refer to this as the below threshold case. As $`E`$ is increased from zero, the fixed point for the angular momentum system rotates about the $`𝒥_y`$ direction eventually reaching the equatorial plane at $`𝒥_x=N/2`$ at the threshold condition. The oscillator system always has zero amplitude below threshold. If we linearise around this fixed point we discover that it is an unstable hyperbolic point with time constant proportional to $`\frac{1}{\sqrt{𝒥_z^{}}}`$. Note that this time constant goes to infinity as the fixed point is approached as is typical for a hyperbolic fixed point.
We now consider the above threshold case
$$\frac{2E}{N\mathrm{\Omega }}1\text{(above threshold)}$$
(14)
Clearly the value of $`|𝒥_x|`$ cannot increase above $`N/2`$. Indeed there is no fixed point above threshold. However there is a special solution curve that continuously joins to the below threshold case for phase curves with $`=0`$.
To see this we consider making a canonical transformation by a rotation in both the $`XY`$ plane and in the $`𝒥_x,𝒥_y`$ plane (see figure 1). The canonical transformations are
$`X`$ $`=`$ $`\overline{X}\mathrm{cos}\theta +\overline{Y}\mathrm{sin}\theta `$ (15)
$`Y`$ $`=`$ $`\overline{Y}\mathrm{cos}\theta \overline{X}\mathrm{sin}\theta `$ (16)
$`𝒥_x`$ $`=`$ $`\overline{𝒥}_x\mathrm{cos}\theta \overline{𝒥}_y\mathrm{sin}\theta `$ (17)
$`𝒥_y`$ $`=`$ $`\overline{𝒥}_x\mathrm{cos}\theta +\overline{𝒥}_y\mathrm{sin}\theta `$ (18)
The Hamiltonian then takes the form
$$=\overline{X}(\overline{𝒥}_x+\chi \mathrm{cos}\theta )\overline{Y}(\overline{𝒥}_y\chi \mathrm{sin}\theta )$$
(19)
The phase curves with $`=0`$ now correspond to either
$`\overline{X}`$ $`=`$ $`0;\overline{𝒥}_y=\chi \mathrm{sin}\theta `$ (20)
or
$`\overline{Y}`$ $`=`$ $`0;\overline{𝒥}_x=\chi \mathrm{cos}\theta `$ (21)
These phase curves smoothly join the fixed point at threshold if $`𝒥_x=N/2`$ which implies
$$\mathrm{cos}\theta =\frac{N\mathrm{\Omega }}{2E}$$
(22)
These solutions are illustrated in figure 1. Note that as $`E\mathrm{}`$ we have that $`\overline{𝒥}_x`$ eventually points in the direction of $`𝒥_y`$ while phase curve in the oscillator phase space points along the $`Y`$ axis, indicating that for large driving the system is essentially a particle in a linear potential which accelerates at constant rate. These results were first obtained by Alsing and Carmichael.
## III Quantum states
First note that the ground state when there is no driving is $`|j,j_z|0_v`$ with a zero eigenvalue. This ground state corresponds to the fixed point of the semiclassical model with zero oscillator amplitude and angular momentum pointing in the $`𝒥_z`$ direction. We postulate that as the driving is increased form zero the ground state of the Hamiltonian Eq(4) is given by
$$|_0=S(r)R(\theta )|j,j_z|0_v$$
(23)
where $`|j,j_z|0_v`$ corresponds to all ions in the ground state and the vibrational mode in the ground state. The operator $`S(r)`$ is a squeezing operator defined by
$$S^{}(r)aS(r)=\mu a+\nu a^{}$$
(24)
with $`\mu =\mathrm{cosh}r,\nu =\mathrm{sinh}r`$.
The rotation operator $`R(\theta )`$ is defined by
$$R(\theta )=e^{\theta (\widehat{J}_+\widehat{J}_{})}$$
(25)
and corresponds to a rotation of $`2\theta `$ around the $`\widehat{J}_y`$ axis. Consider now
$$H_I|_0=SR\left(R^{}S^{}HSR\right)|j,j_z|0_v$$
(26)
If we now transform the Hamiltonian and require that
$$R^{}S^{}HSR|j,j_z|0_v=0$$
(27)
we find the following conditions,
$`\nu (1+\mathrm{cos}2\theta )`$ $`=`$ $`\mu (1\mathrm{cos}2\theta )`$ (28)
$`\mathrm{\Omega }j\mathrm{sin}2\theta `$ $`=`$ $`E`$ (29)
which requires that
$$\mathrm{cos}2\theta =e^{2r}$$
(30)
and the ground state energy is taken to be $`_0=0`$. The ground state is thus a product of a squeezed state for the vibrational mode and a rotated angular momentum state, rotated about the $`\widehat{J}_y`$ axis.
The above results are consistent with the semiclassical approximation. The mean amplitude of a squeezed vacuum state is zero, corresponding to the semiclassical fixed point at $`\overline{X}=\overline{Y}=0`$ while the rotation around the $`\widehat{J}_y`$ axis corresponds to the semiclassical fixed point at $`\overline{𝒥}_x=\chi `$.
If we continue to increase $`E`$ above the threshold value the system adiabatically follows a zero energy state, although this is no longer a ground state. In fact the canonical transformation used in the semiclassical analysis can be applied to the quantum operator valued Hamiltonian. The result is the same as the semiclassical case, Eq (19) with all variables replaced with the corresponding operators. The zero energy state then corresponds to the zero energy eigenstate of $`\widehat{Y}\mathrm{cos}\theta \widehat{X}\mathrm{sin}\theta `$ with $`\mathrm{cos}\theta =N\mathrm{\Omega }/2E`$. This is of course just a rotated, infinitely squeezed state. The electronic state is likewise a angular momentum eigenstate rotated from $`|j,j`$ in the equatorial plane (orthogonal to $`\widehat{J}_z`$). Thus above threshold the zero energy eigenstate deforms continuously from the state at threshold.
Let us summarise these results. For no driving the ground state corresponds to the oscillator in the ground state and all ions in the ground state. As the driving is increased, but kept below threshold, this state deforms to a squeezed oscillator state while the collective spin system begins to rotate about the $`\widehat{J}_y`$ axis. Note that the mean oscillator amplitude $`a`$ remains zero as does the mean of the $`y`$-component of the collective spin. As the driving increases through the threshold value, this state changes its character so that a non zero value of $`\widehat{J}_y`$ is acquired and the oscillator is infinitely squeezed in a direction at an angle $`\mathrm{cos}\theta =N\mathrm{\Omega }/2E`$ to the below threshold squeezing. This morphological change of the state as the driving passes the semiclassical critical point is a quantum phase transition. The quantum phase transition can be seen in the mean value for $`\widehat{J}_y`$ and $`\widehat{J}_z`$ as shown in figure 2. Below threshold the scaled mean values are given by
$`{\displaystyle \frac{\widehat{J}_y}{N/2}}`$ $`=`$ $`0`$ (31)
$`{\displaystyle \frac{\widehat{J}_z}{N/2}}`$ $`=`$ $`\sqrt{1x^2}`$ (32)
and above threshold we have
$`{\displaystyle \frac{\widehat{J}_y}{N/2}}`$ $`=`$ $`\sqrt{1{\displaystyle \frac{1}{x^2}}}`$ (33)
$`{\displaystyle \frac{\widehat{J}_z}{N/2}}`$ $`=`$ $`0`$ (34)
where $`x=2E/N\mathrm{\Omega }`$.
What are the experimental manifestations of this transition ? Needless to say no one is ever going to observe an infinitely squeezed state in an experiment. So what does happens at $`2\theta =\pi /2`$ when the electronic state is the $`\widehat{J}_x`$ eigenstate $`|jj_x`$ and the vibrational mode appears to be infinitely squeezed ? Is such a state physically possible ? Suppose for example we begin in the ground state of the Hamiltonian with no driving ($`E=0`$) which is simply $`|j,j_z|0_v`$, and adiabatically increase the driving strength. It would appear that the system would then adiabatically evolve into the squeezed vibrational state described above. If we were ever able to reach the case $`2\theta =\pi /2`$ we would have reached an infinite energy state for the vibrational mode at a finite driving strength. Clearly this is not possible and to understand why it is useful to reconsider the semiclassical dynamics for this model. The adiabatic approximation requires that we vary the driving strength on a time scale slower than all other time scales in the system. The key time scale for the ground state variation is just the time scale associated with the hyperbolic unstable fixed point, $`(N^2/4\chi ^2)^{1/2}`$, which goes to infinity as we approach $`\frac{2E}{\mathrm{\Omega }N}=1`$. Thus the adiabatic increase of the driving must proceed infinitely slowly, that is it must be switched to the finite value $`E=\frac{\mathrm{\Omega }N}{2}`$ in an infinite amount of time. This pumps an infinite amount of energy into the system and results in infinite squeezing in the centre of mass vibrational mode. Obviously in practice this cannot be achieved so the totally squeezed ground state is not possible. However it will still be possible to achieve some squeezing of the vibrational mode at smaller values of the driving. This would make an interesting observation for current ion trap experiments even with only a few ions. The squeezing of the vibrational mode can be observed using the dynamical method of reference
In current ion trap experiments, laser cooling techniques allow the centre of mass mode to be prepared in the ground state. Unfortunately it does not stay there. Heating due to a variety of sources, including fluctuating linear potentials, lead to an irreversible evolution away from the ground state. If such heating is present during the coupling of the electronic and vibrational motions, irreversible dynamics will be spread to the collective spin degrees of freedom as well.
As an example we consider what happens if we use the Tavis-Cummings interaction (excitation on first red sideband) in the presence of strong heating. Heating of the centre-of-mass mode due to fluctuating liner potentials may be described in the interaction picture by the master equation,
$$\frac{dW}{dt}=i\mathrm{\Omega }[a\widehat{J}_++a^{}\widehat{J}_{},W]+\frac{\gamma }{2}\left(𝒟[a]+𝒟[a^{}]\right)W$$
(35)
where $`W`$ is the density operator for the spin and vibrational degrees of freedom and the superoperator $`𝒟`$ is defined by $`𝒟[A]\rho =2A\rho A^{}A^{}A\rho \rho A^{}A.`$ The irreversible term corresponds to two point processes in which phonons are removed or added from centre of mass mode at the rates $`\gamma a^{}a`$ and $`\gamma aa^{}`$ respectively. This does not change any first order moments, however it does lead to a diffusion in energy as $`\frac{da^{}a}{dt}=\gamma .`$ The effect of heating can be included in the semiclassical analysis by adding an appropriate stochastic term. In the Ito calculus the effect is to add to the equations for $`X,Y`$ terms of the form
$`dX`$ $`=`$ $`(\mathrm{})+\sqrt{\gamma }dW_x(t)`$ (36)
$`dY`$ $`=`$ $`(\mathrm{})+\sqrt{\gamma }dW_y(t)`$ (37)
where $`dW_i(t)`$ are independent Wiener processes. If the heating rate is small enough these terms can be neglected. However if they are large new steady states can occur in the semiclassical and quantum descriptions which will be described in a future publication.
We would like to thank Howard Carmichael for useful discussions.
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# On the Axiomatics of the 5-dimensional Projective Unified Field Theory of Schmutzer
## 1 Introduction
As it is well known the 5-dimensional idea of a unified field theory goes back to the works of Kaluza and Klein . The pioneers of the projective approach to this theory were Veblen and van Dantzig . Later this approach was developed further by many other authors.
An essential progress in this projective type of theories was done by Jordan who first took into consideration the occurring scalar field which inevitably appears in this theory. However, the field equations used by him were unacceptable.
A basically different approach to a projective field theory was proposed by E.Schmutzer who (according to the requirements of a unified field theory) developed further and applied a basis vectors formalism initiated by Hessenberg, Schouten and others in the theory of manifolds. He had no longer considered the scalar field mentioned above to be an auxiliary one. On the contrary he associated this field with a new phenomenon of nature being on one the level of gravitation and electromagnetism. In 1980 he introduced the new term “scalarism” for this phenomenon. For this hypothetically used scalar field with fundamental importance in the PUFT Schmutzer introduced the term “scalaric field” in order to distinguish it from the various other scalar fields in physics. The most interesting and important results of application of PUFT are presented in the Appendix.
Beside the projective relativistic theory many authors were actively developing further the initial Kaluza - Klein theory aiming at a unified field theory of elementary particles. Here we only refer to the monographs of Wesson and Vladimirov , where one can find references to the historical, mathematical and physical literature on this subject.
Concluding this introduction we would like to mention the new 5-dimensional original field theory by Wesson recently appeared and offered for discussion.
In the first two versions of PUFT (see and respectively) Schmutzer used 5-dimensional Einstein-like field equations.<sup>1</sup><sup>1</sup>1Greek indices run from 1 to 5, Latin indices from 1 to 4; the signatures are: of the 5-dimensional metric $`(++++)`$, of the space-time metric $`(+++)`$. Comma means the partial and semicolon the covariant derivative, respectively.
$`\underset{\mu \nu }{\overset{5}{R}}{\displaystyle \frac{1}{2}}g_{\mu \nu }\stackrel{5}{R}+\mathrm{\Lambda }_{\mu \nu }=\kappa _0\mathrm{\Theta }_{\mu \nu }`$ (1)
Afterwards with the help of a special projection procedure (details can be found in the papers quoted above) a system of 4-dimensional field equations describing gravitation, electromagnetism and scalarism was derived. Here $`\kappa _0={\displaystyle \frac{8\pi G}{c^4}}`$ is Einstein’s gravitational constant, $`\theta _{\alpha \epsilon }`$ is the so-called energy projector of the non-geometrized matter named “substrate”, $`\stackrel{5}{R}_{\alpha \epsilon }`$ is the 5-dimensional Ricci tensor, $`\stackrel{5}{R}`$ is the 5-dimensional curvature invariant, and
$`\text{a)}\mathrm{\Lambda }_{\mu \nu }=\lambda _0g_{\mu \nu }\text{resp.}\text{b)}\mathrm{\Lambda }_{\mu \nu }=\lambda _0(g_{\mu \nu }+s_\mu s_\nu )S`$ (2)
are analogs of the cosmological terms in Version I and Version II, respectively. Here $`g_{\mu \nu }`$ is the metric tensor. In a special frame $`\left\{X^\mu \right\}`$ the unit vector $`s^\mu `$ has (see the next section for details) the following form $`s^\mu ={\displaystyle \frac{X^\mu }{S}}`$, where $`S=\sqrt{X_\mu X^\mu }=S_0e^\sigma `$ ($`S_0`$ is an arbitrary constant of the dimension of length).
The 5-dimensional Ricci tensor and the 5-dimensional curvature invariant, both mentioned above, are defined as follows: $`\stackrel{5}{R}{}_{\alpha \epsilon }{}^{}=\stackrel{5}{R}{}_{\alpha \epsilon \tau }{}^{\tau },\stackrel{5}{R}=\stackrel{5}{R}{}_{\alpha }{}^{\alpha },`$ where
$`\stackrel{5}{R}{}_{\mu \nu \epsilon }{}^{\alpha }\left\{\begin{array}{c}\alpha \\ \mu \epsilon \end{array}\right\}_{,\nu }\left\{\begin{array}{c}\alpha \\ \mu \nu \end{array}\right\}_{,\epsilon }+\left\{\begin{array}{c}\tau \\ \mu \epsilon \end{array}\right\}\left\{\begin{array}{c}\alpha \\ \tau \nu \end{array}\right\}\left\{\begin{array}{c}\tau \\ \mu \nu \end{array}\right\}\left\{\begin{array}{c}\alpha \\ \tau \epsilon \end{array}\right\}.`$ (3)
For physical reasons in the following the Gauss system of units is chosen.
Version I. The 4-dimensional field equations (without a cosmological term: $`\mathrm{\Lambda }_{\mu \nu }=0`$) have the following form:
$`\stackrel{4}{R}{}_{mn}{}^{}{\displaystyle \frac{1}{2}}\stackrel{4}{g}{}_{mn}{}^{}\stackrel{4}{R}=\kappa (E_{mn}+S_{mn}+\mathrm{\Theta }_{mn}).`$ (4)
These equations are generalized 4-dimensional equations for the gravitational field, where $`\kappa =\kappa _0e^\sigma `$,
$`E_{mn}={\displaystyle \frac{1}{4\pi }}(B_{mk}H_n^k+{\displaystyle \frac{1}{4}}g_{mn}B_{jk}H^{jk})`$ (5)
is the electromagnetic energy–momentum tensor and
$`S_{mn}={\displaystyle \frac{1}{2\kappa }}(\sigma _{,m}\sigma _{,n}+\sigma _{,m;n}g_{mn}(\sigma _{,k}\sigma ^{,k}+\sigma _{;k}^{,k})`$ (6)
is the energy–momentum tensor of the scalaric field $`\sigma `$. Further following field equations hold:
$`a)H_{;n}^{mn}={\displaystyle \frac{4\pi }{c}}j^m,b)B_{[mn,k]}=0,c)H_{mn}=e^{3\sigma }B_{mn},`$ (7)
$`\sigma _{;k}^{,k}=\kappa _0\left({\displaystyle \frac{2}{3}}\vartheta +{\displaystyle \frac{1}{8\pi }}B_{kj}H^{kj}\right).`$ (8)
These are the electromagnetic field equations and the field equation for the scalaric field $`\sigma `$. Here the following notations were used: $`\stackrel{4}{R}_{mn}`$, $`\stackrel{4}{R}`$ are the Ricci tensor and the curvature invariant in the 4-dimensional space-time, respectively. $`H^{mn}`$, $`B^{mn}`$ are the electromagnetic induction tensor and the electromagnetic field strength tensor, respectively. The quantity $`\vartheta `$ being one of the sourses of the scalaric field is called scalaric substrate energy density.
The idea of developing the Version II was to remove the second order derivatives in the energy – momentum tensor (6) of the scalaric field. By means of a modified projection formalism it became possible to obtain a system of equations being slightly different from the analogous system of the version I, given by the equations (6), (7c) and (8), namely:
$`\stackrel{4}{R}{}_{mn}{}^{}{\displaystyle \frac{1}{2}}\stackrel{4}{R}\stackrel{4}{g}{}_{mn}{}^{}+\lambda _0S_0\stackrel{4}{g}{}_{mn}{}^{}=\kappa _0(\theta _{mn}+E_{mn}+S_{mn}),`$ (9)
$`S_{mn}={\displaystyle \frac{3}{2\kappa _0}}(\sigma _{,m}\sigma _{,n}{\displaystyle \frac{1}{2}}g_{mn}\sigma _{,k}\sigma ^{,k}),`$ (10)
$`a)H_{;n}^{mn}={\displaystyle \frac{4\pi }{c}}j^m,b)B_{[mn,k]}=0,c)H_{mn}=e^{3\sigma }B_{mn},`$ (11)
$`\sigma _{;k}^{,k}=\kappa _0\left({\displaystyle \frac{2}{3}}\vartheta +{\displaystyle \frac{1}{8\pi }}B_{kj}H^{kj}\right).`$ (12)
However, the new projection formalism led to other problems, particularly in the spinor theory. Therefore approximately in 1994 E.Schmutzer left this version II and offered version III.
In the version III by deeply founded considerations on the level of the Lagrange-Hamilton formalism the following new 5-dimensional field equations were found :
$`R_{\mu \nu }{\displaystyle \frac{1}{2}}g_{\mu \nu }\stackrel{5}{R}{\displaystyle \frac{1}{S}}S_{,\mu ;\nu }+{\displaystyle \frac{K_0\kappa _0}{S^2}}S_{,\mu }S_{,\nu }{\displaystyle \frac{1}{S}}s_\mu s_\nu [2(1{\displaystyle \frac{1}{2}}K_0\kappa _0)S_{;\tau }^{,\tau }`$
$`{\displaystyle \frac{3}{S}}(1{\displaystyle \frac{1}{2}}K_0\kappa _0)S_{,\tau }S^{,\tau }+{\displaystyle \frac{3\lambda _0}{S}}{\displaystyle \frac{S}{2}}\stackrel{5}{R}]+{\displaystyle \frac{1}{S}}g_{\mu \nu }[S_{;\tau }^{,\tau }`$
$`{\displaystyle \frac{1}{S}}(1+{\displaystyle \frac{1}{2}}K_0\kappa _0)S_{,\tau }S^{,\tau }+{\displaystyle \frac{1}{S}}\lambda _0]=\kappa _0\mathrm{\Theta }_{\mu \nu }`$ (13)
($`\lambda _0`$ is a kind of cosmological constant. $`K_0`$ is a free constant, where Schmutzer preferred the choice $`K_0=2`$). Compared with the Einstein-like field equation (1) this is a rather complicated equation, but it fulfils important physical demands mentioned in the . From (1) one obtains
$`R_{mn}{\displaystyle \frac{1}{2}}g_{mn}R+{\displaystyle \frac{\lambda _0}{S_0^2}}e^{2\sigma }g_{mn}=\kappa _0(E_{mn}+S_{mn}+\mathrm{\Theta }_{mn}),`$ (14)
$`E_{mn}={\displaystyle \frac{1}{4\pi }}(B_{mk}H_n^k+{\displaystyle \frac{1}{4}}g_{mn}B_{jk}H^{jk}),`$ (15)
$`S_{mn}=K_0(\sigma _{,m}\sigma _{,n}{\displaystyle \frac{1}{2}}g_{mn}\sigma _{,k}\sigma ^{,k}),`$ (16)
$`a)H_{;n}^{mn}={\displaystyle \frac{4\pi }{c}}j^m,b)B_{[mn,k]}=0,c)H_{mn}=e^{2\sigma }B_{mn},`$ (17)
$`\sigma _{;k}^{,k}={\displaystyle \frac{1}{K_0}}\left(\vartheta +{\displaystyle \frac{1}{8\pi }}B_{kj}H^{kj}\right)+{\displaystyle \frac{2\lambda _0}{\kappa _0K_0S_0^2}}e^{2\sigma }.`$ (18)
In the present paper we introduce a new geometrical axiomatics for the Schmutzer theory. By means of this axiomatics we can give a new geometrical interpretation of physical results obtained in the PUFT.
## 2 Projection formalism
As it is well known, the physical basis of the 5-dimensional Projective Unified Field Theory is supported by the following mathematical theorem: The semidirect product of the Abelian group of gauge transformations (electromagnetism) and of the group of the general 4-dimensional coordinate transformations (gravitation) corresponds to the group being homomorphic to the group of all 5-dimensional homogeneous of degree 1 coordinate transformations
$`X^\mu ^{}=X^\mu ^{}(X^\nu )={\displaystyle \frac{1}{\alpha }}X^\mu ^{}(\alpha X^\nu )(\alpha =\mathrm{const}).`$ (19)
This mathematical theorem allows us to assume that the geometry, constructed on this group, can be a basis for the geometrization of the electromagnetic, the gravitational and the scalaric field.
From the equation (19) and Euler’s theorem on homogeneous functions follows that these special coordinates $`X^\mu `$ in the 5-dimensional space $`_5`$ are transformed as the components of a vector:
$`X^\mu ^{}=X_{,\nu }^\mu ^{}X^\nu .`$ (20)
Further the vector
$`𝓡=X^\mu 𝑬_\mu ,𝑬_\mu ={\displaystyle \frac{}{X^\mu }}`$ (21)
can be regarded as 5-dimensional radius vector. This was a very important starting point of Schmutzer in 1957. Also in the following this vector field $``$ plays a fundamental role. Of course, it is possible to introduce in the space $`_5`$ arbitrary coordinates $`y^\mu =y^\mu (X^\alpha )`$.
In context with the theorem mentioned above we should remark that the 4-dimensional coordinates $`\{x^i\}`$ in the space-time should satisfy the equation
$`x_{,\nu }^iX^\nu =0`$ (22)
(for details see ).
In order to construct a projection formalism, let us consider the congruence
$`y^\nu =y^\nu (x^i,\tau )`$ (23)
of integral curves of the vector field $``$, where $`\tau `$ is a continuous parameter specified along each curve ($`x^i=\mathrm{const}`$) of this congruence.
The congruence (23) is the starting point of our consideration. In general the quantities $`x^i`$ are not the first four coordinates of a 5-dimensional coordinate system.
Hereinafter we will consider the 4-dimensional hypersurface $`\tau (y^\nu )=\mathrm{const}`$ to be the 4-dimensional space-time. Moreover, the parameter $`\tau `$ should be chosen to make tangent vectors
$`\xi ^\nu (x^i,\tau ){\displaystyle \frac{}{\tau }}y^\nu (x^i,\tau )`$ (24)
coinciding with the vectors $`X^\nu `$:
$`\xi ^\nu (x^i,\tau )=X^\nu .`$ (25)
It is important to point out that the equation (25) is only valid in the frame $`\left\{X^\nu \right\}`$; but it can always be rewritten in an arbitrary frame $`\{y^\nu \}`$: $`\xi ^\nu (x^i,\tau )=^\nu `$, where $`^\nu `$ are components of the vector $``$ in the coordinate basis $`𝒆_\nu ={\displaystyle \frac{}{y^\nu }}`$. Let us emphasize that all equations containing vectors $`X^\mu `$ are only valid in the special frame. Henceforth it will not be specially accentuated.
According to (22) we postulate equality to zero of the Lie derivative with respect to $``$ for any 4-dimensional quantity (i.e. quantity which depends only on 4-dimensional coordinates). It is quite natural to extend the introduced postulate on all 5-dimensional vectors and tensors which further will be associated with the 4-dimensional quantities :
$`\mathrm{\pounds }_𝓡𝑻=0(\mathrm{\pounds }_𝓡T_\beta \mathrm{}^\alpha \mathrm{}=0).`$ (26)
In the coordinate basis (21) one can rewrite the last equation (26) in the form
$`T_{\nu _1\mathrm{}\nu _m,\lambda }^{\mu _1\mathrm{}\mu _n}X^\lambda =(nm)T_{\nu _1\mathrm{}\nu _m}^{\mu _1\mathrm{}\mu _n}.`$ (27)
The geometrical quantities satisfying the projector condition (27) are called projectors . By applying the projector condition (26) to a metric tensor $`g`$ we obtain
$`\mathrm{\pounds }_𝓡𝒈=0(g_{\mu \nu ,\epsilon }X^\epsilon =2g_{\mu \nu }).`$ (28)
From the last equation follows that the 5-dimensional radius vector $``$ is a Killing vector. Thus the congruence (23) is a Killing congruence (see for example ).
In order to study geometrical properties of this congruence we introduce a unit vector $`s`$, i.e.
$`𝒔={\displaystyle \frac{𝓡}{S}}\left(s^\mu ={\displaystyle \frac{X^\mu }{S}}\right),`$ (29)
where $`S=\sqrt{𝒈(𝓡,𝓡)}=\sqrt{g_{\mu \nu }X^\mu X^\nu }`$. From the definition (29) it is clear that $`s`$ is the unit tangential vector field to the lines of the congruence (23).
In order to provide a description of geometrical properties of the congruence (23) we introduce, as usually, the following quantities:
$`a)G^\mu s^\nu s_{;\nu }^\mu ,b)\omega _{\mu \nu }P_\mu ^\tau P_\nu ^\epsilon s_{[\tau ,\epsilon ]}c)D_{\mu \nu }P^\tau _\mu P^\epsilon _\nu s_{(\tau ;\epsilon )},`$ (30)
where
| $`G^\mu `$ | $``$ | | the first curvature vector of the lines of the congruence; |
| --- | --- | --- | --- |
| $`\omega _{\mu \nu }`$ | $``$ | | the angular velocity tensor of the congruence; |
| $`D_{\mu \nu }`$ | $``$ | | rate-of-strain tensor of the congruence. |
The quantity
$`P_\mu ^\tau =\delta _\mu ^\tau s^\tau s_\mu `$ (31)
is the projection tensor. The semicolon means the Riemannian covariant derivative ($`\stackrel{R}{}`$):
$`\underset{𝒆_\tau }{\overset{R}{}}𝒆_\alpha =\left\{\begin{array}{c}\epsilon \\ \alpha \tau \end{array}\right\}𝒆_\epsilon \left(𝒆_\alpha ={\displaystyle \frac{}{x^\alpha }}\right)`$ (32)
with $`\left\{\begin{array}{c}\epsilon \\ \alpha \tau \end{array}\right\}{\displaystyle \frac{1}{2}}g^{\epsilon \sigma }\left(g_{\sigma \alpha ,\tau }+g_{\tau \sigma ,\alpha }g_{\alpha \tau ,\sigma }\right)`$. If we take into account that the vector field $``$ is Killingian, we obtain
a) $`G^\mu `$ $`={\displaystyle \frac{1}{2S}}X^{\epsilon \mu }s_\epsilon ={\displaystyle \frac{S^{,\mu }}{S}},`$
b) $`\omega _{\mu \nu }`$ $`={\displaystyle \frac{1}{2S}}P_\mu ^\epsilon P_\nu ^\tau X_{\tau \epsilon },`$
c) $`D_{\mu \nu }`$ $`=0,`$ (33)
where the following abbreviation was used:
$`X_{\mu \nu }=X_{\nu ,\mu }X_{\mu ,\nu }.`$ (34)
From the equations (30) and (2) we obtain the following important relations:
$`s_{\mu ;\nu }=D_{\mu \nu }+\omega _{\mu \nu }+G_\mu s_\nu ,`$ (35)
$`{\displaystyle \frac{1}{2S}}X_{\nu \mu }=\omega _{\mu \nu }+{\displaystyle \frac{1}{S}}(s_\mu S_{,\nu }s_\nu S_{,\mu }).`$ (36)
From the relation (2b) follows that in general a holonomic hypersurface orthogonal to the given congruence does not exist. (The case $`\omega _{\mu \nu }=0`$ is physically not interesting, since further the angular velocity of the congruence will be associated with the electromagnetic field). Therefore in contrast to Schmutzer’s orthogonality approach of space-time (based on the basis vector formalism) here we want to offer an alternative version of this problem: we shall identify space-time with a 4-dimensional hypersurface in the 5-dimensional space abandoning the requirement of orthogonality of this hypersurface to the congruence.
Let us consider some hypersurface $`\tau (X^\alpha )=\mathrm{const}`$. As far as a parameter $`\tau `$ cannot univalently be derived from the equations (24) and (25), then hypersurfaces $`\tau (X^\alpha )=\mathrm{const}`$ are not defined univalently either. Therefore we can choose in $`_5`$ an arbitrary hypersurface which we shall identify with hypersurface $`\tau (X^\alpha )=0`$. This hypersurface should only satisfy the condition $`X^\alpha \tau _{,\alpha }0`$. With an exponential map we can extend it along the lines of congruence (23) to a finite region in $`_5`$. Thus we receive a one-parametric set of hypersurfaces. Hence from the equations (24) and (25) follows that
$`<\mathrm{d}\tau ,𝝃>=1(\xi ^\epsilon \tau _{,\epsilon }=X^\epsilon \tau _{,\epsilon }=1)`$ (37)
and
$`\mathrm{\pounds }_𝓡\mathrm{d}\tau =0(X^\epsilon \tau _{,\alpha ,\epsilon }+X_{,\alpha }^\epsilon \tau _{,\epsilon }=0).`$ (38)
From the last relation we can conclude that the one-form $`\mathrm{d}\tau `$ which further we also shall denote by $`\zeta `$ satisfies the projector condition (27):
$`\mathrm{\pounds }_𝓡𝜻=0,𝜻\mathrm{d}\tau (\zeta _{\mu ,\tau }X^\tau =\zeta _\mu ).`$ (39)
The unit one-form $`𝝂=\mathrm{\Lambda }𝜻`$ also fulfills this condition:
$`\mathrm{\pounds }_𝓡𝝂=0(\nu _{\tau ,\mu }X^\mu =\nu _\tau ),`$ (40)
where $`\mathrm{\Lambda }=<𝝂,𝓡>=\nu _\epsilon X^\epsilon `$ and $`\nu _\epsilon \nu ^\epsilon =1.`$
Above we introduced the projection tensor $`P_\epsilon ^\alpha `$. However, the hypersurface $`\tau (X^\alpha )=0`$ (we also shall denote it by $`_4`$) is not orthogonal to the congruence (23). Therefore it is possible to define two more projection tensors:
$`\text{a)}b_{\alpha \epsilon }g_{\alpha \epsilon }\nu _\alpha \nu _\epsilon ,\text{b)}h_\epsilon ^\alpha g_\epsilon ^\alpha \xi ^\alpha \zeta _\epsilon .`$ (41)
All these projection tensors satisfy the projector condition (27):
$`P_{\epsilon ,\nu }^\alpha X^\nu =0,h_{\epsilon ,\nu }^\alpha X^\nu =0,b_{\alpha \epsilon ,\nu }X^\nu =2b_{\alpha \epsilon }.`$ (42)
The projection tensor $`b_{\alpha \epsilon }`$ sometimes is called the first fundamental form of $`_4`$ or the induced metric on $`_4`$. (In the following we shall define the induced metric on $`_4`$ in a somewhat different way). The tensor $`\chi _{\alpha \epsilon }`$ defined on the hypersurface $`\tau =0`$ by
$`\chi _{\alpha \epsilon }b_\alpha ^\mu b_\epsilon ^\nu \nu _{(\mu ;\nu )}={\displaystyle \frac{1}{2\mathrm{\Lambda }}}\mathrm{\pounds }_𝝀b_{\alpha \epsilon },`$ (43)
is called the second fundamental form or the exterior curvature of $`\tau =0`$. Here the following abbreviations were used:
$`\lambda ^\epsilon \mathrm{\Lambda }\nu ^\epsilon =X^\epsilon 𝒳^\epsilon ,𝒳^\epsilon b_\alpha ^\epsilon X^\alpha .`$ (44)
The above introduced projection tensors in general differ from each other. Therefore the question, which of them should be used for the projection of 5-dimensional vectors and tensors into the 4-dimensional hypersurface, is not trivial. In order to give an answer to this question, we consider the map $`\varphi `$:
$`\varphi :_5\stackrel{\varphi }{}_4.`$ (45)
The map $`\varphi `$ should be defined in such a way to make mapped quantities not depending on the parameter $`\tau `$, i.e. on the “fifth coordinate” (cylinder condition). This requirement means that all points laying on the same line of the congruence are mapped to the same point on the hypersurface $`\tau (X^\nu )=\mathrm{const}`$. The elementary map of this type is an exponential map (see Fig. 1).
The coordinates of the point $`P_1`$ satisfy the relation
$`X^\alpha (P_1)=X^\alpha (x_0^m,\tau _1)=X^\alpha (x_0^m,0)\mathrm{exp}(\tau _1),`$ (46)
where $`X^\alpha (x_0^m,0)=X^\alpha (P^{^{}})`$. Using the equations (24) and (25) one can obtain
$`\varphi _\tau :X^\alpha (P)=\mathrm{exp}(\tau )X^\alpha (\varphi _\tau (P)).`$ (47)
Now we have to discuss how the vectors and tensors are transformed by the map $`\varphi _\tau `$.
Let $`V`$ be a tangent vector to the curve $`\lambda (t)`$ at the point $`P_1`$, having the following form in local coordinates in a neighborhood of the point $`P_1`$:
$`X^\alpha (\lambda (t))=X^\alpha (P_1)+tV^\alpha ,`$ (48)
where $`V^\alpha `$ are the components of the vector $`V`$ ($`𝑽={\displaystyle \frac{}{t}}`$) in the coordinate basis $`𝑬_\alpha `$, i.e.
$`𝑽=V^\alpha {\displaystyle \frac{}{X^\alpha }}.`$ (49)
Comparing the equation (48) with the following series expansion:
$`X^\alpha (\lambda (t))`$ $`=`$ $`X^\alpha (x^m(\lambda (t)),\tau (\lambda (t)))=`$
$`=`$ $`X^\alpha (x_0^m,\tau _1)+\left(X_{,m}^\alpha {\displaystyle \frac{\mathrm{d}x^m}{\mathrm{d}t}}+{\displaystyle \frac{X^\alpha }{\tau }}{\displaystyle \frac{\mathrm{d}\tau }{\mathrm{d}t}}\right)|_{P_1}t`$
$`+`$ $`O(t^2),`$ (50)
we obtain
$`V^\alpha |_{P_1}=\left(X_{,m}^\alpha {\displaystyle \frac{\mathrm{d}x^m}{\mathrm{d}t}}+\xi ^\alpha {\displaystyle \frac{\mathrm{d}\tau }{\mathrm{d}t}}\right)|_{P_1}\left(\xi ^\alpha ={\displaystyle \frac{\mathrm{d}X^\alpha }{\mathrm{d}\tau }}|_{x^m=\mathrm{const}}\right).`$ (51)
The curve $`\lambda (t)`$ can be projected by the exponential map $`\varphi _{\tau (\lambda )}`$ onto the hypersurface $`\tau =0`$. The notation $`\varphi _{\tau (\lambda )}`$ should accentuate that each point of the curve $`\lambda (t)`$ is mapped by the proper exponential map $`\varphi _\tau `$ ($`\tau `$ depends on $`t`$ ). We denote the mapped curve obtained by this procedure by $`\gamma (t)`$:
$`\varphi (\lambda (t))=\gamma (t),`$ (52)
where $`\varphi `$ means $`\varphi _{\tau (\lambda (t))}`$.
Further we shall consider only vector fields $`V`$ commuting with $`{\displaystyle \frac{}{\tau }}`$, i.e. the vector fields being projectors. In this case the maps of curves $`\lambda (t)`$ and $`\lambda ^{^{}}(t)`$ ($`\lambda ^{^{}}(0)=P^{^{}}`$) coincide:
$`\varphi (\lambda (t))=\varphi (\lambda ^{^{}}(t))=\gamma (t),`$ (53)
where $`\lambda ^{^{}}(t)=\varphi _{\tau _1}(\lambda (t))`$. Therefore, without any further restriction we may consider only such curves whose initial points $`P_1`$ ($`P_1=\lambda (0)`$) belong to the hypersurface $`\tau =0`$, i.e. $`P_1=P^{^{}}_4`$.
For the mapped curve $`\gamma (t)`$ following expansion is valid
$`X^\alpha (\gamma (t))`$ $`=`$ $`X^\alpha [\left(x_0^m+{\displaystyle \frac{\mathrm{d}x^m}{\mathrm{d}t}}+O(t^2)\right),0]`$
$`=`$ $`X^\alpha (x_0^m,0)+\left(X_{,m}^\alpha {\displaystyle \frac{\mathrm{d}x^m}{\mathrm{d}t}}\right)|_Pt+O(t^2).`$ (54)
From the equation (51) follows
$`V^\alpha =\left({\displaystyle \frac{\mathrm{d}x^m}{\mathrm{d}t}}e_m+{\displaystyle \frac{\mathrm{d}\tau }{\mathrm{d}t}}\xi \right)(X^\alpha ),`$ (55)
where the vectors are defined by
$`𝒆_m={\displaystyle \frac{}{x^m}},𝝃={\displaystyle \frac{}{\tau }}.`$ (56)
The equation (55) implies that the following relation for the vector field $`V`$ is fulfilled:
$`𝑽={\displaystyle \frac{\mathrm{d}x^m}{\mathrm{d}t}}e_m+{\displaystyle \frac{\mathrm{d}\tau }{\mathrm{d}t}}\xi .`$ (57)
Thus at the point $`P_1`$ ($`P_1_4;\gamma (0)=\lambda ^{^{}}(0)=P_1`$) the 4-dimensional $`T_{P_1}`$ and 5-dimensional $`T_{\varphi (P_1)}`$ vector spaces can be constructed as follows :
$`𝑽T_{P_1},\varphi _{}𝑽T_{\varphi (P_1)}.`$ (58)
Here we used the abbreviations (compare with ):
$`\varphi _{}𝑽\left({\displaystyle \frac{}{t}}\right)_\gamma |_{\varphi (P_1)}={\displaystyle \frac{\mathrm{d}x^m}{\mathrm{d}t}}|_{P_1}𝒆_m,`$ (59)
where
$`{\displaystyle \frac{\mathrm{d}x^m}{\mathrm{d}t}}=x_{,\alpha }^mV^\alpha .`$ (60)
It is necessary to note that $`\varphi (P_1)=P_1`$, $`\gamma (0)=\lambda ^{^{}}(0)=P_1`$.
It is easy to show that the one-form $`𝒆^m\mathrm{d}x^m`$ and the vectors $`𝒆_m={\displaystyle \frac{}{x^m}}`$ satisfy the equations
$`<𝜻,𝒆_m>=0,<𝒆^m,𝝃>=0.`$ (61)
These equations imply that one can rewrite the projector $`h_\epsilon ^\alpha `$ in the form
$`h_\epsilon ^\alpha =g_m^\alpha g_\epsilon ^m=g_\epsilon ^\alpha \xi ^\alpha \zeta _\epsilon ,`$ (62)
where we used the definitions
$`\text{a)}g_\epsilon ^m=<𝒆^m,𝒆_\epsilon >=x_{,\epsilon }^m,\text{b)}g_m^\epsilon =<𝒆^\epsilon ,𝒆_m>=X_{,m}^\epsilon .`$ (63)
Apart from that it is easy to show that between the quantities $`𝒆_m`$, $`𝒆^m`$, $`𝒆_\epsilon `$ and $`𝒆^\epsilon `$ the following relation is valid:
$`\text{a)}𝒆_\epsilon =g_\epsilon ^m𝒆_m+\zeta _\epsilon 𝝃,\text{b)}𝒆^\epsilon =g_m^\epsilon \mathrm{d}x^m+\xi ^\epsilon \mathrm{d}\tau .`$ (64)
The last relation and the definition (59) lead us to
$`\stackrel{}{𝑽}\varphi _{}𝑽=(g_\alpha ^mV^\alpha )𝒆_m=\stackrel{\epsilon }{\stackrel{}{V}}𝒆_\epsilon ,`$ (65)
where we used the abbreviation
$`\stackrel{\epsilon }{\stackrel{}{V}}h_\alpha ^\epsilon V^\alpha .`$ (66)
Thus in the tangent vector space $`T_P`$ it is possible to define a 4-dimensional subspace $`\underset{P_1}{\overset{}{T}}`$ ($`\underset{P_1}{\overset{}{T}}T_{P_1}`$):
$`\underset{P}{\overset{}{T}}=\{\stackrel{}{𝑽}:\stackrel{}{𝑽}\underset{P}{\overset{}{T}},\varphi _{}\stackrel{}{𝑽}=\stackrel{}{𝑽}\}.`$ (67)
The equation $`\varphi _{}𝑽=\stackrel{\epsilon }{\stackrel{}{V}}𝒆_\epsilon `$ should be interpreted in the following way:
$`\varphi _{}𝑽=(x_{,\alpha }^mV^\alpha )𝒆_mT(_4).`$ (68)
In a vector space $`T(_5)`$ it is possible to construct a 4-dimensional vector space formed by the vectors of the type $`(h_\alpha ^\epsilon V^\alpha )𝒆_\epsilon `$. The spaces $`T(_4)`$ and $`\stackrel{}{T}(_5)`$ are isomorphic:
$`(x_{,\alpha }^mV^\alpha )𝒆_m(h_\alpha ^\epsilon V^\alpha )𝒆_\epsilon .`$ (69)
The map $`\varphi _{}`$, namely
$$T_P(_5)\stackrel{\varphi _{}}{}T_{\varphi (P)}(_4),$$
naturally induces the map $`\varphi ^{}`$ for the one-forms:
$$T_{\varphi (P)}^{}(_4)\stackrel{\varphi ^{}}{}T_P^{}(_5),$$
where for all $`𝝎T_{\varphi (P_1)}^{}`$ and for all $`𝑽T_{P_1}`$ the next relation is valid:
$`<\varphi ^{}𝝎,𝑽>|_{P_1}=<𝝎,\varphi _{}𝑽>|_{\varphi (P_1)}.`$ (70)
The set of all one-forms satisfying the relation
$`\varphi ^{}\stackrel{}{𝝎}=\stackrel{}{𝝎}`$ (71)
forms a linear space $`\stackrel{}{\stackrel{}{T}}(_5)T^{}(_5)`$, where $`T^{}(_5)`$ is the space of all one-forms. From (70) follows that for any one-form from $`\stackrel{}{\stackrel{}{T}}(_5)`$ holds
$`<\stackrel{}{𝝎},𝝃>=0,(\underset{\alpha }{\overset{}{\omega }}X^\alpha =0).`$ (72)
There are two possible ways to associate elements of $`T^{}(_5)`$ with elements of $`\stackrel{}{\stackrel{}{T}}(_5)`$:
$$\text{a)}\omega _\alpha \underset{\alpha }{\overset{}{\omega }}=h_\alpha ^\epsilon \omega _\epsilon ,$$
$$\text{b)}\omega _\alpha \underset{\alpha }{\overset{}{\omega }}=P_\alpha ^\epsilon \omega _\epsilon .$$
In both cases the quantities $`\underset{\alpha }{\overset{}{\omega }}`$ satisfy the equation (72) automatically. Hereinafter quantities with tilde will be associated with physical quantities in the space-time. However, the definition a) cannot be accepted, as in this case the following relations would be valid:
$$\stackrel{\alpha }{\stackrel{}{V}}h_\epsilon ^\alpha V^\epsilon \stackrel{\alpha \epsilon }{\stackrel{}{g}}\underset{\epsilon }{\overset{}{M}},\underset{\alpha }{\overset{}{M}}h_\alpha ^\epsilon V_\epsilon \underset{\epsilon \alpha }{\overset{}{g}}\stackrel{\epsilon }{\stackrel{}{V}},$$
where
$$\underset{\alpha \beta }{\overset{}{g}}=h_\alpha ^\epsilon h_\beta ^\sigma g_{\epsilon \sigma },\stackrel{\alpha \sigma }{\stackrel{}{g}}=h_\nu ^\alpha h_\mu ^\sigma g^{\nu \mu }.$$
On the contrary, the definition b) is consistent. In this case the following relations will be valid:
$`\text{a)}\underset{\alpha \beta }{\overset{}{g}}`$ $``$ $`P_\alpha ^\epsilon P_\beta ^\sigma g_{\epsilon \sigma }=P_{\alpha \beta },`$
$`\text{b)}\stackrel{\alpha \sigma }{\stackrel{}{g}}`$ $``$ $`h_\nu ^\alpha h_\mu ^\sigma g^{\nu \mu }=g^{\alpha \sigma }2X^{(\alpha }\zeta ^{\sigma )}+{\displaystyle \frac{1}{\mathrm{\Lambda }^2}}X^\alpha X^\sigma ,`$
$`\text{c)}\underset{\nu }{\overset{\mu }{\stackrel{}{g}}}`$ $``$ $`h_\epsilon ^\mu P_\nu ^\alpha g_\alpha ^\epsilon =h_\nu ^\mu .`$ (73)
Using (2), for an arbitrary vector $`V`$ we obtain:
$`\stackrel{\alpha }{\stackrel{}{V}}=\stackrel{\alpha \epsilon }{\stackrel{}{g}}\underset{\epsilon }{\overset{}{V}},\underset{\alpha }{\overset{}{V}}=\underset{\alpha \epsilon }{\overset{}{g}}\stackrel{\epsilon }{\stackrel{}{V}}.`$ (74)
The last results can be summarized in the sentence: The 5-dimensional tensors are to be projected onto hypersurface $`\tau (X^\alpha )=0`$ (projected quantities are denoted by a tilde) with the help of the procedure:
$`T_\mathrm{}\nu ^\mu \mathrm{}\stackrel{\varphi }{}\stackrel{~}{T}_\mathrm{}\nu ^\mu \mathrm{}h_\sigma ^\mu \mathrm{}P_\nu ^\tau \mathrm{}T_\mathrm{}\tau ^\sigma \mathrm{}.`$ (75)
The quantities $`x^i`$ being introduced as parameters earlier and parametrizing the congruence (23) can be used as coordinates in $`_4`$. Let us point out that it is necessary to require a certain continuity for the quantities $`x^i`$. Apart from that these quantities are defined accurately within the following transformation: $`x^ix^i^{}=x^i^{}(x^j)`$. In this case the vectors $`𝒆_m={\displaystyle \frac{}{x^m}}`$ and the one-forms $`𝒆^m=\mathrm{d}x^m`$ form a basis in $`T(_4)`$ and $`T^{}(_4)`$, either. These bases satisfy the following relations:
$`<𝒆^m,𝒆_n>=\delta _n^m,[𝓡,𝒆_m]=0,<𝝂,𝒆_m>=0.`$ (76)
Using the equations (44), (63) and (64), we can find several important relations:
$`\text{a)}𝒆_\alpha h_\epsilon ^\alpha =g_\epsilon ^m𝒆_m,\text{b)}𝒆^\alpha h_\alpha ^\epsilon =g_m^\epsilon 𝒆^m;`$ (77)
$`\text{a)}𝒆_\alpha P_\epsilon ^\alpha =g_\epsilon ^m𝒆_m+(\zeta _\epsilon s_\epsilon )𝒔,\text{b)}𝒆^\alpha P_\alpha ^\epsilon =(g_m^\alpha P_\alpha ^\epsilon )𝒆^m;`$ (78)
$`\text{a)}𝒆_\alpha b_\epsilon ^\alpha =(g_\alpha ^mb_\epsilon ^\alpha )𝒆_m,\text{b)}𝒆^\alpha b_\alpha ^\epsilon =g_m^\epsilon 𝒆^m+𝒳^\epsilon \mathrm{d}\tau .`$ (79)
We already mentioned that the tangent spaces $`T(_4)`$ and $`\stackrel{}{T}(_5)`$ are isomorphic. Therefore one can write:
$`T^mg_\epsilon ^m\stackrel{\epsilon }{\stackrel{}{T}},\stackrel{\epsilon }{\stackrel{}{T}}=g_m^\epsilon T^m,`$ (80)
$`\omega _mg_m^\epsilon \stackrel{}{\omega },\underset{\epsilon }{\overset{}{\omega }}=g_\epsilon ^m\omega _m.`$ (81)
Thus the projection procedure from $`T(_5)`$ into $`T(_4)`$ is defined as follows:
$`T^\alpha T^m=\underset{\epsilon }{\overset{m}{\stackrel{}{g}}}T^\epsilon ,\omega _\alpha \omega _m=\underset{m}{\overset{\epsilon }{\stackrel{}{g}}}\omega _\epsilon ,`$ (82)
where we used the abbreviation
$`\underset{m}{\overset{\epsilon }{\stackrel{}{g}}}P_\alpha ^\epsilon g_m^\alpha g_m^\epsilon ,\underset{\epsilon }{\overset{m}{\stackrel{}{g}}}h_\epsilon ^\alpha g_\alpha ^m=g_\epsilon ^m.`$ (83)
The metric induced on the hypersurface $`_4`$ will be denoted further by $`\stackrel{}{𝒈}`$ (in the theory of surfaces one understands under the induced metric the quantity $`b_{\mu \nu }`$ defined by means of (41a) ). This metric satisfies the following relations:
$`\text{a)}\underset{mn}{\overset{}{g}}\underset{m}{\overset{\alpha }{\stackrel{}{g}}}\underset{n}{\overset{\beta }{\stackrel{}{g}}}g_{\alpha \beta },\text{b)}\stackrel{mn}{\stackrel{}{g}}\underset{\alpha }{\overset{m}{\stackrel{}{g}}}\underset{\beta }{\overset{n}{\stackrel{}{g}}}g^{\alpha \beta },`$
$`\text{c)}\underset{n}{\overset{m}{\stackrel{}{g}}}\underset{\alpha }{\overset{m}{\stackrel{}{g}}}\underset{n}{\overset{\epsilon }{\stackrel{}{g}}}g_\epsilon ^\alpha =\delta _n^m,\text{d)}\stackrel{mn}{\stackrel{}{g}}\underset{mk}{\overset{}{g}}=\delta _k^n.`$ (84)
It is necessary to accentuate that the relation
$`𝒈(𝒆_m,𝒆_n)\underset{mn}{\overset{}{g}},\underset{mn}{\overset{}{g}}=\stackrel{}{𝒈}(𝒆_m,𝒆_n)`$ (85)
is fulfilled, where
$`𝒈g_{\alpha \epsilon }\mathrm{d}X^\alpha \mathrm{d}X^\epsilon ,\stackrel{}{𝒈}\underset{mn}{\overset{}{g}}\mathrm{d}x^m\mathrm{d}x^n.`$ (86)
At the end we have to present further two important relations that immediately follow from (76):
$`\text{a)}x_{,\alpha }^mX^\alpha =0,\text{b)}X_{,m}^\epsilon \zeta _\epsilon =0.`$ (87)
## 3 Field Equations
In the introduction we mentioned that in the course of four decades three neighboring versions of Schmutzer’s 5-dimensional Projective Unified Field Theory came into being. All these versions are based on the following 5-dimensional field equations:
$`𝒢_{\alpha \epsilon }=\kappa _0\theta _{\alpha \epsilon }.`$ (88)
The explicit expression of the symmetric tensor $`𝒢_{\alpha \epsilon }`$ for the versions II and III of PUFT can be found, using the equations (1) and (1), respectively. In order to obtain a 4-dimensional field equations these 5-dimensional equations have to be projected onto the 4-dimensional space-time. The equation (88) can always be written in the following form:
$`\underset{\mu \nu }{\overset{}{𝒢}}+2\underset{(\mu }{\overset{}{𝒢}}s_{\nu )}+𝒢s_\mu s_\nu =\kappa _0(\underset{\mu \nu }{\overset{}{\stackrel{5}{\theta }}}+2\underset{(\mu }{\overset{}{\stackrel{5}{\theta }}}s_{\nu )}+\stackrel{}{\stackrel{5}{\theta }}s_\mu s_\nu ),`$ (89)
where the abbreviations are given by
$`\begin{array}{ccc}\text{}\text{a)}\underset{\mu \nu }{\overset{}{𝒢}}P_\mu ^\alpha P_\nu ^\beta 𝒢_{\alpha \beta },\hfill & \text{b)}\underset{\mu }{\overset{}{𝒢}}P_\mu ^\alpha s^\beta 𝒢_{\alpha \beta },\hfill & \text{c)}𝒢s^\alpha s^\beta 𝒢_{\alpha \beta },\hfill \\ \text{d)}\underset{\mu \nu }{\overset{}{\stackrel{5}{\theta }}}P_\mu ^\alpha P_\nu ^\beta \theta _{\alpha \beta },\hfill & \text{e)}\underset{\mu }{\overset{}{\stackrel{5}{\theta }}}P_\mu ^\alpha s^\beta \theta _{\alpha \beta },\hfill & \text{f)}\stackrel{5}{\theta }s^\alpha s^\beta \theta _{\alpha \beta }.\hfill \end{array}`$ (92)
It is easy to see that the equation (89) is equivalent to the following set of equations:
$`\begin{array}{c}\text{}\text{a)}\underset{\mu \nu }{\overset{}{𝒢}}=\kappa _0\underset{\mu \nu }{\overset{}{\stackrel{5}{\theta }}},\hfill \\ \text{b)}\underset{\mu }{\overset{}{𝒢}}=\kappa _0\underset{\mu }{\overset{}{\stackrel{5}{\theta }}},\hfill \\ \text{c)}𝒢=\kappa _0\stackrel{5}{\theta }\hfill \end{array}`$ (96)
Further we can see that the following correspondence is valid:
| | equation (96a) | $``$ | generalized Einstein equations, |
| --- | --- | --- | --- |
| | equation (96b) | $``$ | generalized Maxwell equations, |
| | equation (96c) | $``$ | Field equation of the scalaric field. |
Henceforth the 4-dimensional hypersurface $`_4`$ will be identified with the space-time. The physical metrics $`\stackrel{4}{𝒈}`$ of the space-time can be defined in different ways (all space-time quantities will be denoted by an index “4”). For example, we can identify the 4-dimensional physical metric $`\stackrel{4}{𝒈}`$ with the metric $`\stackrel{}{𝒈}`$ induced on the hypersurface $`_4`$:
$`\stackrel{\mathrm{𝟒}}{𝒈}=\stackrel{}{𝒈}(\stackrel{4}{g}{}_{\mu \nu }{}^{}=\underset{\mu \nu }{\overset{}{g}},\stackrel{4}{g}{}_{}{}^{\mu \nu }=\stackrel{\mu \nu }{\stackrel{}{g}}).`$ (97)
In this case we obtain version I or III of PUFT if we use the 5-dimensional equations (1) or (1), respectively. However, it is physically possible to connect these metrics $`\stackrel{4}{𝒈}`$ and $`\stackrel{}{𝒈}`$ by a conformal transformation:
$`\stackrel{\mathrm{𝟒}}{𝒈}=e^\sigma \stackrel{}{𝒈}(\stackrel{4}{g}{}_{\mu \nu }{}^{}=e^\sigma \underset{\mu \nu }{\overset{}{g}},\stackrel{4}{g}{}_{}{}^{\mu \nu }=e^\sigma \stackrel{\mu \nu }{\stackrel{}{g}}).`$ (98)
In this case the 5-dimensional Einstein-like equations (1) lead to the system of equations of version II of PUFT. In order to consider both these cases simultaneously we rewrite the equations (97) and (98) in the form:
$`\stackrel{\mathrm{𝟒}}{𝒈}=e^{ϵ\sigma }\stackrel{}{𝒈}\text{with}ϵ=\{\begin{array}{ccc}0\hfill & & \text{Version I}\text{+}\text{Version III}\hfill \\ 1\hfill & & \text{Version II}\hfill \end{array}`$ (101)
The projection formalism can be simplified by using a non-Riemannian connection in the 5-dimensional space and considering the Riemannian connection in the 4-dimensional space-time as induced.
### 3.1 Connection on $`_5`$
Let us introduce an induced (affine) connection on the hypersurface $`_4`$ denoted hereinafter as $`\underset{}{\overset{4}{}}`$. The induced connection $`\underset{}{\overset{4}{}}`$ and the connection on $`_5`$ (denoted as $`\underset{}{\overset{5}{}}`$) are connected in the following way:
$`\underset{\epsilon }{\overset{4}{}}\underset{\nu \mathrm{}}{\overset{\mu \mathrm{}}{\stackrel{}{T}}}=h_\sigma ^\mu \mathrm{}P_\epsilon ^\alpha P_\nu ^\lambda \mathrm{}\underset{\lambda \mathrm{}\left|\right|\alpha }{\overset{\sigma \mathrm{}}{\stackrel{}{T}}},`$ (102)
where
$`T_{\mathrm{}\left|\right|\alpha }^{\mathrm{}}=\underset{\alpha }{\overset{5}{}}T_{\mathrm{}}^{\mathrm{}},\underset{\alpha }{\overset{5}{}}\underset{𝒆_\alpha }{\overset{5}{}},\underset{\alpha }{\overset{4}{}}\underset{𝒆_\alpha }{\overset{4}{}}.`$ (103)
Henceforth we assume that the connection on $`_4`$ is Riemannian, i.e. metrical and symmetric:
$`\text{a)}\underset{\epsilon }{\overset{4}{}}\underset{\mu \nu }{\overset{4}{g}}=0,\text{b)}\underset{\alpha }{\overset{4}{}}\underset{\epsilon }{\overset{4}{}}f=\underset{\epsilon }{\overset{4}{}}\underset{\alpha }{\overset{4}{}}f.`$ (104)
Since the 4-dimensional covariant derivative is defined only for the projected vectors (see (102)), the function $`f`$ should satisfy the condition (26): $`\mathrm{\pounds }_𝓡f=f_{,\alpha }X^\alpha =0`$. The 4-dimensional covariant derivative (with respect to $`\underset{}{\overset{4}{}}`$) in the direction of the basis vectors $`𝒆_m`$ ($`𝒆_m={\displaystyle \frac{}{x^m}}`$) will be denoted by a semicolon:
$`\underset{𝒆_m}{\overset{4}{}}\stackrel{k}{\stackrel{}{T}}\underset{m}{\overset{4}{}}\stackrel{k}{\stackrel{}{T}}\underset{;m}{\overset{k}{\stackrel{}{T}}}=g_\alpha ^kg_m^\epsilon \underset{\epsilon }{\overset{4}{}}\stackrel{\alpha }{\stackrel{}{T}},`$ (105)
where $`\stackrel{}{𝑻}=𝒆_m\stackrel{m}{\stackrel{}{T}}=𝒆_\alpha \stackrel{\alpha }{\stackrel{}{T}}.`$
As it is well known, the Riemannian connection is completely defined by means of a metric. Therefore the relations (104) are in fact conditions for the 5-dimensional connection $`\underset{}{\overset{5}{}}`$. In particular, from (104) follows that the 5-dimensional connection $`\underset{}{\overset{5}{}}`$ has to satisfy the following relations:
$`\begin{array}{cc}\text{}\text{a)}\hfill & \underset{\epsilon \alpha \beta }{\overset{}{Q}}P_\epsilon ^\tau P_\alpha ^\nu P_\beta ^\mu Q_{\tau \nu \mu }=ϵ\sigma _{,\epsilon }\underset{\alpha \beta }{\overset{}{g}}=ϵ\sigma _{,\epsilon }P_{\alpha \beta },\hfill \\ \text{b)}\hfill & \underset{\alpha \beta }{\overset{\tau }{\stackrel{}{S}}}P_\alpha ^\nu P_\beta ^\mu h_\epsilon ^\tau S_{\nu \mu }^\epsilon =0,\hfill \\ \text{c)}\hfill & h_\alpha ^\mu P_\beta ^\nu X_{\left|\right|\nu }^\alpha =0,\hfill \end{array}`$ (109)
where the usual definitions
$`\text{a)}g_{\alpha \beta \left|\right|\epsilon }=Q_{\epsilon \alpha \beta },\text{b)}g_{\left|\right|\epsilon }^{\alpha \beta }=Q_\epsilon ^{\alpha \beta }\text{c)}S_{\alpha \beta }^\gamma =\mathrm{\Gamma }_{[\alpha \beta ]}^\gamma `$ (110)
are used. One can easily verify that the 5-dimensional connection in general is nonsymmetric and nonmetrical.
For this reason we write the 5-dimensional connection coefficients $`\mathrm{\Gamma }_{\mu \nu }^\epsilon `$ in the following form:
$`\mathrm{\Gamma }_{\mu \nu }^\epsilon =\left\{\begin{array}{c}\epsilon \\ \mu \nu \end{array}\right\}+\sigma _{\mu \nu }^\epsilon ,`$ (111)
where
$`\sigma _{\mu \nu }^\epsilon =(S_{\nu \mu }^\epsilon +S_{\nu \mu }^\epsilon S_{\mu }^{}{}_{\nu }{}^{\epsilon })+{\displaystyle \frac{1}{2}}(Q_{\nu \mu }^\epsilon +Q_{\mu }^{}{}_{\nu }{}^{\epsilon }Q_{\nu \mu }^\epsilon ).`$ (112)
The 5-dimensional connection cannot be found uniquely from the demands (109). However, the 5-dimensional field equations (see (1) and (1)) only contain Riemannian covariant derivatives, and therefore, the 5-dimensional connection $`\underset{}{\overset{5}{}}`$ is only an auxiliary quantity. Thus within some restrictions, the 5-dimensional connection coefficients $`\mathrm{\Gamma }_{\mu \nu }^\epsilon `$ can be chosen arbitrarily. Therefore we choose the 5-dimensional connection on $`_5`$ in a certain way to make calculations as simple as possible. First let us in general assume:
$`\underset{\xi }{\overset{5}{}}e_\epsilon =\mathrm{{\rm Y}}_\epsilon ^\nu 𝒆_\nu (\mathrm{\Gamma }_{\mu \nu }^\epsilon X^\nu =\mathrm{{\rm Y}}_\mu ^\epsilon ),`$ (113)
where $`\mathrm{{\rm Y}}_\mu ^\epsilon `$ is an arbitrary projector. Taking into account the relation
$`\left\{\begin{array}{c}\epsilon \\ \mu \nu \end{array}\right\}X^\mu =g_\nu ^\epsilon +{\displaystyle \frac{1}{2}}X_\nu ^\epsilon ,`$ (114)
which follows immediately from(28), we obtain
$`\sigma _{\mu \nu }^\epsilon X^\nu =\mathrm{\Sigma }_\mu ^\epsilon {\displaystyle \frac{1}{2}}X_\mu ^\epsilon ,`$ (115)
where we introduced the abbreviation
$`\mathrm{\Sigma }_\mu ^\epsilon \mathrm{{\rm Y}}_\mu ^\epsilon g_\mu ^\epsilon .`$ (116)
From the conditions (109b) and (109c) we obtain the following relation for the torsion tensor:
$`S_{\alpha \beta }^\gamma =A_{\alpha \beta }X^\gamma +{\displaystyle \frac{1}{2S^2}}h_\tau ^\gamma (X_\alpha \mathrm{\Sigma }_\beta ^\tau X_\beta \mathrm{\Sigma }_\alpha ^\tau ).`$ (117)
Here we used the abbreviation
$`A_{\epsilon \tau }=S_{\epsilon \tau }^\mu \zeta _\mu .`$ (118)
It is possible to show that the 5-dimensional connection has the simplest form if according to (109a) we put
$`Q_{\epsilon \alpha \beta }=ϵ\sigma _{,\epsilon }P_{\alpha \beta }.`$ (119)
In this case the quantity $`\mathrm{\Sigma }_{\mu \nu }`$ will be antisymmetric (with no other limitations):
$`\mathrm{\Sigma }_{\mu \nu }=\mathrm{\Sigma }_{\nu \mu }.`$ (120)
The tensors $`S_{\alpha \beta }^\gamma `$ and $`\sigma _{\alpha \beta }^\gamma `$ in this case are given by:
$`S_{\alpha \beta }^\gamma ={\displaystyle \frac{X^\gamma }{2S^2}}[X_{\alpha \beta }+(s_\alpha S_{,\beta }s_\beta S_{,\alpha })]+{\displaystyle \frac{1}{2S^2}}(X_\alpha \mathrm{\Sigma }_\beta ^\gamma X_\beta \mathrm{\Sigma }_\alpha ^\gamma ),`$ (121)
$`\sigma _{\lambda \mu \epsilon }=`$ $``$ $`{\displaystyle \frac{1}{2S}}[(X_{\lambda \mu }s_\epsilon X_{\epsilon \lambda }s_\mu +X_{\mu \epsilon }s_\lambda )`$
$`+`$ $`2s_\mu (s_\lambda S_{,\epsilon }s_\epsilon S_{,\lambda })]{\displaystyle \frac{s_\mu }{S}}\mathrm{\Sigma }_{\lambda \epsilon }`$
$``$ $`{\displaystyle \frac{ϵ}{2S}}(P_{\lambda \mu }S_{,\epsilon }P_{\mu \epsilon }S_{,\lambda }P_{\lambda \epsilon }S_{,\mu }).`$ (122)
From the last relation follows that the connection on $`_5`$ has the simplest form if the quantity $`\mathrm{\Sigma }_{\mu \nu }`$ is defined according to (120) as follows:
$`\mathrm{\Sigma }_{\epsilon \nu }=G_\epsilon X_\nu G_\nu X_\epsilon ,`$ (123)
where the abbreviation (2a) was used. Substituting the last expression into the relations (121) and (3.1), we obtain:
$`S_{\alpha \beta }^\sigma =s^\sigma \omega _{\alpha \beta },`$ (124)
$`\sigma _{\epsilon \tau \nu }`$ $`=`$ $`\omega _{\epsilon \tau }s_\nu \omega _{\nu \epsilon }s_\tau +\omega _{\tau \nu }s_\epsilon `$
$``$ $`{\displaystyle \frac{ϵ}{2}}(G_\epsilon P_{\tau \nu }+G_\tau P_{\nu \epsilon }G_\nu P_{\tau \epsilon }).`$ (125)
At the end of this section we would like to point out once again that the connection on $`_5`$ is an intermediate quantity. Its choice does not lead to any physical consequences. It can be shown that for any choice of $`\mathrm{\Sigma }_{\epsilon \alpha }`$ and $`Q_\nu ^{\epsilon \alpha }`$ (these quantities have to satisfy the conditions (109) only) the 4-dimensional physical equations get the same form. However, in the general case all calculations become unwieldy. Therefore we don’t present them here fully; hereinafter we only will consider the case (119) and (123). Thus the torsion tensor $`S_{\alpha \beta }^\epsilon `$ and the tensor $`\sigma _{\alpha \beta }^\epsilon `$ take the simplest form, i.e. (124) and (3.1), respectively.
### 3.2 Projection of the Curvature Tensor and Related Quantities
Now we have to analyse the equation (96). In order to do it, we can use the general relation
$`\underset{\epsilon \left|\right|\mu \left|\right|\lambda }{\overset{}{T}}\underset{\epsilon \left|\right|\lambda \left|\right|\mu }{\overset{}{T}}=\underset{\nu }{\overset{}{T}}G_{\epsilon \mu \lambda }^\nu +2\underset{\epsilon \left|\right|\alpha }{\overset{}{T}}S_{\lambda \mu }^\alpha ,`$ (126)
where
$`G_{\beta \gamma \delta }^\alpha =\mathrm{\Gamma }_{\beta \delta ,\gamma }^\alpha \mathrm{\Gamma }_{\beta \gamma ,\delta }^\alpha +\mathrm{\Gamma }_{\varrho \gamma }^\alpha \mathrm{\Gamma }_{\beta \delta }^\varrho \mathrm{\Gamma }_{\varrho \delta }^\alpha \mathrm{\Gamma }_{\beta \gamma }^\varrho .`$ (127)
To project the equation (126) onto space-time we need the following two relations:
$`\begin{array}{ccc}\text{a)}\underset{\beta }{\overset{4}{}}\underset{\alpha }{\overset{4}{}}\underset{\sigma }{\overset{}{T}}\hfill & =\hfill & \underset{\delta \left|\right|\epsilon \left|\right|\nu }{\overset{}{T}}P_\beta ^\nu P_\sigma ^\delta P_\alpha ^\epsilon ,\hfill \\ \text{b)}\underset{\beta }{\overset{4}{}}\underset{\alpha }{\overset{4}{}}\stackrel{\sigma }{\stackrel{}{T}}\hfill & =\hfill & \underset{\left|\right|\epsilon \left|\right|\nu }{\overset{\delta }{\stackrel{}{T}}}P_\beta ^\nu P_\alpha ^\epsilon h_\delta ^\sigma .\hfill \end{array}`$ (130)
The relation (130a) follows immediately from the equation
$`\underset{\beta }{\overset{4}{}}\underset{\alpha }{\overset{4}{}}\underset{\sigma }{\overset{}{T}}=(\underset{\gamma \left|\right|\tau }{\overset{}{T}}P_\delta ^\gamma P_\epsilon ^\tau )_{\left|\right|\nu }P_\beta ^\nu P_\sigma ^\delta P_\alpha ^\epsilon ,`$ (131)
in which the covariant derivatives
$`\text{a)}s_{\left|\right|\nu }^\mu =G^\mu s_\nu ,\text{b)}s_{\mu \left|\right|\nu }=G_\mu s_\nu `$ (132)
are substituted. The equation (130b) can be similarly proved. Proceeding further, we suppose that the vector $`\underset{\epsilon }{\overset{}{T}}`$ satisfies the projector condition (27). The equation
$`s^\delta \underset{\epsilon \left|\right|\delta }{\overset{}{T}}=(G^\lambda \underset{\lambda }{\overset{}{T}})s_\epsilon .`$ (133)
is fulfilled in this case. Using the relations (124), (132) and (133) we obtain the interesting equality
$`\underset{\beta }{\overset{4}{}}\underset{\alpha }{\overset{4}{}}\underset{\sigma }{\overset{}{T}}\underset{\alpha }{\overset{4}{}}\underset{\beta }{\overset{4}{}}\underset{\sigma }{\overset{}{T}}=\underset{\sigma \alpha \beta }{\overset{\nu }{\stackrel{}{G}}}\underset{\nu }{\overset{}{T}},`$ (134)
where according to (75) we used the abbreviation
$`\underset{\sigma \alpha \beta }{\overset{\nu }{\stackrel{}{G}}}h_\mu ^\nu P_\sigma ^\gamma P_\alpha ^\delta P_\beta ^\varrho G_{\gamma \delta \varrho }^\mu .`$ (135)
The equation (134) being written in the basis $`𝒆_n`$ (see (105) ) is given by:
$`T_{s;a;b}T_{s;b;a}=T_nG_{sab}^n,`$ (136)
where
$`T_n=g_n^\alpha \underset{\alpha }{\overset{}{T}},G_{sab}^n=g_\nu ^ng_s^\sigma g_a^\alpha g_b^\beta \underset{\sigma \alpha \beta }{\overset{\nu }{\stackrel{}{G}}}.`$ (137)
As the equation (136) is correct for all space-time vectors, the following relation is valid:
$`\stackrel{4}{R}{}_{sab}{}^{n}=G_{sab}^n.`$ (138)
Here $`\stackrel{4}{R}_{sab}^n`$ is the 4-dimensional Riemannian curvature tensor
$`\stackrel{4}{R}{}_{mnk}{}^{a}\underset{,n}{\overset{4}{\left\{\begin{array}{c}a\\ mk\end{array}\right\}}}\underset{,k}{\overset{4}{\left\{\begin{array}{c}a\\ mn\end{array}\right\}}}+\stackrel{4}{\left\{\begin{array}{c}t\\ mk\end{array}\right\}}\stackrel{4}{\left\{\begin{array}{c}a\\ tn\end{array}\right\}}\stackrel{4}{\left\{\begin{array}{c}t\\ mn\end{array}\right\}}\stackrel{4}{\left\{\begin{array}{c}a\\ tk\end{array}\right\}},`$ (139)
where we used the usual definition
$`\stackrel{4}{\left\{\begin{array}{c}k\\ at\end{array}\right\}}{\displaystyle \frac{1}{2}}\stackrel{4}{g}{}_{}{}^{ks}(\stackrel{4}{g}{}_{sa,t}{}^{}+\stackrel{4}{g}{}_{ts,a}{}^{}\stackrel{4}{g}{}_{at,s}{}^{}).`$ (140)
Obtaining the equation (136) we applied the following relation:
$`\underset{\nu }{\overset{4}{}}\underset{\mu }{\overset{}{T}}g_m^\mu g_n^\nu =T_{m;n}\text{(}T_m=g_m^\nu \underset{\nu }{\overset{}{T}}\text{)}.`$ (141)
Let $`T_\mu `$ be an arbitrary one-form (covariant vector). According to the projection formalism developed above we can project this one-form onto the hypersurface $`_4`$: $`T^\mu \underset{\mu }{\overset{}{T}}P_\mu ^\alpha T_\alpha `$. Then the equations (111) and (130) imply:
$`\underset{\lambda }{\overset{4}{}}\underset{\epsilon }{\overset{4}{}}\underset{\mu }{\overset{}{T}}=P_\mu ^\alpha P_\epsilon ^\beta P_\lambda ^\gamma (\underset{\alpha ;\beta }{\overset{}{T}}\sigma _{\alpha \beta }^\nu \underset{\nu }{\overset{}{T}})_{\left|\right|\gamma }.`$ (142)
Substituting the relations (2), (35), (3.1), (132), (133) as well as
$`a)P_{\mu \nu \left|\right|\epsilon }=ϵG_\epsilon P_{\mu \nu }s_\epsilon (G_\mu s_\nu +G_\nu s_\mu ),b)P^\nu _{\mu \left|\right|\epsilon }=s_\epsilon (G^\nu s_\mu +G_\mu s^\nu ),`$ (143)
$`s^\delta \underset{\epsilon ;\delta }{\overset{}{T}}=s^\delta \underset{\delta ;\epsilon }{\overset{}{T}}=(\omega _\epsilon ^\lambda s_\epsilon G^\lambda )\underset{\lambda }{\overset{}{T}},`$ (144)
$`G_{\left|\right|\epsilon }^\varrho =ϵG_\epsilon G^\varrho +g^{\varrho \lambda }G_{\lambda \left|\right|\epsilon }`$ (145)
into the last formula, we obtain the result
$`\underset{\lambda }{\overset{4}{}}\underset{\epsilon }{\overset{4}{}}\underset{\mu }{\overset{}{T}}\underset{\epsilon }{\overset{4}{}}\underset{\lambda }{\overset{4}{}}\underset{\mu }{\overset{}{T}}`$ $`=`$ $`\underset{\varrho }{\overset{}{T}}\{\underset{\mu \epsilon \lambda }{\overset{\varrho }{\stackrel{}{\stackrel{5}{R}}}}2\omega _{\epsilon \lambda }\omega _\mu ^\varrho +\omega _{\mu \epsilon }\omega _\lambda ^\varrho \omega _{\mu \lambda }\omega _\epsilon ^\varrho `$
$``$ $`{\displaystyle \frac{ϵ}{2}}[\left(\underset{\lambda }{\overset{4}{}}\underset{\mu }{\overset{}{G}}\right)P_\epsilon ^\varrho +P_{\mu \epsilon }g^{\varrho \alpha }\underset{\alpha \left|\right|\gamma }{\overset{}{G}}P_\lambda ^\gamma +\left(\underset{\epsilon }{\overset{4}{}}\underset{\mu }{\overset{}{G}}\right)P_\lambda ^\varrho `$
$``$ $`P_{\mu \lambda }g^{\varrho \alpha }\underset{\alpha \left|\right|\gamma }{\overset{}{G}}P_\epsilon ^\gamma )]+{\displaystyle \frac{ϵ^2}{4}}[G^\varrho (G_\epsilon P_{\mu \lambda }G_\lambda P_{\mu \epsilon })`$
$`+`$ $`G_\epsilon G_\mu P_\lambda ^\varrho G_\lambda G_\mu P_\epsilon ^\varrho `$
$``$ $`G^\nu G_\nu (P_{\mu \epsilon }P_\lambda ^\varrho P_{\mu \lambda }P_\epsilon ^\varrho )]\}.`$ (146)
Now we are able to analyse the equation (96a). However, before doing it, let us summarize some formulas which are related to the projection formalism.
In the 5-dimensional space $`_5`$ the basis vectors and basis one-forms were denoted by $`𝒆_\mu `$ and $`𝒆^\mu `$, respectively. The 4-dimensional holonomic hypersurface $`_4`$ in the 5-dimensional space $`_5`$ is identified with the 4-dimensional space-time. The quantities projected onto the hypersurface $`_4`$ were denoted by a tilde (see (75)). The quantities $`x^i`$ parametrizing curves of the congruence (23) can be used as coordinates in the space-time. The tangent vectors $`𝒆_i`$ to the coordinate lines ($`𝒆_i={\displaystyle \frac{}{x^i}}`$) of this 4-dimensional coordinate system form a 4-dimensional vector space $`\underset{P}{\overset{}{T}}`$ (see (67). Between the basis vectors $`𝒆_i`$ and $`𝒆_\mu `$ exists consistency (64). Similar relations are valid for the dual basis $`𝒆^i`$ ($`𝒆^i=dx^i`$), too. Thus the 4-dimensional vectors and one-forms can be rewritten in the following form:
$`\begin{array}{cc}a)\hfill & \stackrel{}{𝑽}=\stackrel{\alpha }{\stackrel{}{V}}𝒆_\alpha =V^i𝒆_i(\stackrel{\alpha }{\stackrel{}{V}}=g_i^\alpha V^i,V^i=g_\alpha ^i\stackrel{\alpha }{\stackrel{}{V}})\hfill \\ b)\hfill & \stackrel{}{𝝎}=\underset{\alpha }{\overset{}{\omega }}𝒆^\alpha =\omega _i𝒆^i(\underset{\alpha }{\overset{}{\omega }}=g^i\omega _\alpha ^i,\omega _i=g_i^\alpha \underset{\alpha }{\overset{}{\omega }}).\hfill \end{array}`$ (149)
Let us remember that on the hypersurface $`_4`$ we introduced two metrics: the induced metric $`\stackrel{}{𝒈}`$ and the physical metric $`\stackrel{4}{𝒈}`$. These metrics are connected by means of the relation (101). As the physical metric differs in general from the induced one, one should be careful in defining 4-dimensional physical quantities.
Using the abbreviation
$`\stackrel{4}{\omega }{}_{mn}{}^{}\omega _{mn}=g_m^\mu g_n^\nu \underset{\mu \nu }{\overset{}{\omega }}=g_m^\mu g_n^\nu \omega _{\mu \nu }`$ (150)
we obtain from (149) and (101) the following relations
$`\text{a)}\omega _m^ng_n^\mu g_\nu ^n\underset{\mu }{\overset{\nu }{\stackrel{}{\omega }}}=e^{ϵ\sigma }\stackrel{4}{\omega }{}_{m}{}^{n},\text{b)}\omega ^{mn}g_\mu ^mg_\nu ^n\stackrel{\mu \nu }{\stackrel{}{\omega }}=e^{2ϵ\sigma }\stackrel{4}{\omega }{}_{}{}^{mn},`$ (151)
where $`\stackrel{4}{\omega }{}_{m}{}^{n}=\stackrel{4}{g}{}_{}{}^{nk}\stackrel{4}{\omega }_{mk}`$ and $`\stackrel{4}{\omega }{}_{}{}^{mn}=\stackrel{4}{g}{}_{}{}^{mk}\stackrel{4}{g}{}_{}{}^{nl}\stackrel{4}{\omega }_{kl}`$. In a similar way we deduce from (74) and (2a) the equations
$`a)g_m^\mu G_\mu =g_m^\mu \sigma _{,\mu }=\sigma _{,m},b)g^m_\mu \stackrel{}{G}^\mu =\stackrel{}{g}^{\mu \nu }g^m_\nu \sigma _{,\mu }=e^{ϵ\sigma }\sigma ^{,m},`$ (152)
where $`\sigma ^{,m}\stackrel{4}{g}{}_{}{}^{mn}\sigma _{,n}^{}`$.
Let us note that the space-time indices are to be moved with the help of the space-time metric $`\stackrel{4}{𝒈}`$.
It can be shown that for the arbitrarily projected quantities $`\stackrel{}{𝝎}`$ and $`\stackrel{}{𝑽}`$ the relation
$`\underset{\tau }{\overset{}{\omega }}\stackrel{\tau }{\stackrel{}{V}}=\omega _nV^n\text{(}\omega _mg_m^\mu \underset{\mu }{\overset{}{\omega }},V^mg_\tau ^m\stackrel{\tau }{\stackrel{}{V}}\text{)}`$ (153)
is true. From (134) and (138) follows
$`g_s^\sigma g_a^\alpha g_b^\beta (\underset{\beta }{\overset{4}{}}\underset{\alpha }{\overset{4}{}}\underset{\sigma }{\overset{}{T}}\underset{\alpha }{\overset{4}{}}\underset{\beta }{\overset{4}{}}\underset{\sigma }{\overset{}{T}})=\stackrel{4}{R}{}_{sab}{}^{n}\underset{n}{\overset{4}{T}}\text{(}\underset{n}{\overset{4}{T}}g_n^\sigma \underset{\sigma }{\overset{}{T}}\text{).}`$ (154)
Using the relations (3.2) and further the relations (152) to (154), we obtain the final result
$`\underset{mkl}{\overset{a}{\stackrel{}{\stackrel{5}{R}}}}`$ $``$ $`g_\alpha ^ag_m^\mu g_k^\epsilon g_l^\lambda R_{\mu \epsilon \lambda }^\alpha =\text{}`$ (155)
$`=`$ $`\stackrel{4}{R}{}_{mkl}{}^{a}+e^{ϵ\sigma }(2\stackrel{4}{\omega }{}_{kl}{}^{}\stackrel{4}{\omega }{}_{m}{}^{a}\stackrel{4}{\omega }{}_{mk}{}^{}\stackrel{4}{\omega }{}_{l}{}^{a}+\stackrel{4}{\omega }{}_{ml}{}^{}\stackrel{4}{\omega }{}_{k}{}^{a})\text{}`$
$`+`$ $`{\displaystyle \frac{ϵ}{2}}(\stackrel{4}{g}{}_{}{}^{a}{}_{}{}^{k}\sigma _{,m;l}^{}\stackrel{4}{g}{}_{mk}{}^{}\sigma _{;l}^{,a}\stackrel{4}{g}{}_{l}{}^{a}\sigma _{,m;k}^{}+\stackrel{4}{g}{}_{ml}{}^{}\sigma _{;k}^{,a})\text{}`$
$``$ $`{\displaystyle \frac{ϵ^2}{4}}[\sigma ^{,a}(\sigma _{,k}\stackrel{4}{g}{}_{ml}{}^{}\sigma _{,l}\stackrel{4}{g}{}_{mk}{}^{})+\stackrel{4}{g}{}_{l}{}^{a}\sigma _{,k}^{}\sigma _{,m}\text{}`$
$``$ $`\stackrel{4}{g}{}_{k}{}^{a}\sigma _{,m}^{}\sigma _{,l}(\sigma _{,c}\sigma ^{,c})(\stackrel{4}{g}{}_{mk}{}^{}\stackrel{4}{g}{}_{l}{}^{a}\stackrel{4}{g}{}_{ml}{}^{}\stackrel{4}{g}{}_{k}{}^{a})].`$
In order to find the projection of the 5-dimensional Ricci tensor onto space-time $`_4`$ let us consider the relation
$`\underset{mn}{\overset{}{\stackrel{5}{R}}}g_m^\mu g_n^\nu \underset{\mu \nu }{\overset{}{\stackrel{5}{R}}}=g_m^\mu g_n^\nu \left[\stackrel{}{\stackrel{5}{R}}{}_{\mu \nu \varrho }{}^{\varrho }+P_\mu ^\alpha P_\nu ^\beta (s_\varrho s^\sigma \stackrel{5}{R}{}_{\alpha \beta \sigma }{}^{\varrho })\right],`$ (156)
where
$`\underset{\mu \nu }{\overset{}{\stackrel{5}{R}}}P_\mu ^\alpha P_\nu ^\sigma \underset{\alpha \sigma }{\overset{}{\stackrel{5}{R}}},\stackrel{}{\stackrel{5}{R}}{}_{\mu \nu \sigma }{}^{\varrho }h_\lambda ^\varrho P_\mu ^\alpha P_\nu ^\beta P_\sigma ^\tau \stackrel{5}{R}{}_{\mu \nu \sigma \tau }{}^{\lambda }.`$ (157)
From (134), (138) and (3.2) we find
$`g_m^\mu g_k^\nu \stackrel{}{\stackrel{5}{R}}_{\mu \nu \varrho }^\varrho `$ $`=`$ $`\stackrel{4}{R}{}_{mk}{}^{}+3e^{ϵ\sigma }\stackrel{4}{\omega }{}_{ka}{}^{}\stackrel{4}{\omega }{}_{m}{}^{a}{\displaystyle \frac{ϵ}{2}}(\sigma _{,m;k}+{\displaystyle \frac{1}{2}}\stackrel{4}{g}{}_{mk}{}^{}\sigma _{;a}^{,a})\text{}`$ (158)
$`+`$ $`{\displaystyle \frac{ϵ^2}{2}}[\stackrel{4}{g}{}_{mk}{}^{}(\sigma ^{,a}\sigma _{,a})\sigma _{,k}\sigma _{,m}].`$
The second term on the right hand side of the equation (156) can be calculated in the simplest way using the formulas (114) and (36). The result is
$`\stackrel{5}{R}{}_{\alpha \beta \varrho }{}^{\sigma }s_{\sigma }^{}s^\varrho =\sigma _{,\beta ;\alpha }+\sigma _{,\alpha }\sigma _{,\beta }\omega _{\lambda \beta }\omega _\alpha ^\lambda +\omega _{\lambda [\beta }s_{\alpha ]}\sigma ^{,\lambda }s_\alpha s_\beta \sigma ^{,\lambda }\sigma _{,\lambda }.`$ (159)
By substituting the expressions (158) and (159) into (156) we get the formula
$`\stackrel{5}{R}_{mn}`$ $`=`$ $`\stackrel{4}{R}{}_{mn}{}^{}+2e^{ϵ\sigma }\stackrel{4}{\omega }{}_{na}{}^{}\stackrel{4}{\omega }{}_{m}{}^{a}{\displaystyle \frac{ϵ}{2}}\stackrel{4}{g}{}_{mn}{}^{}\sigma _{;a}^{,a}+(1+ϵ{\displaystyle \frac{ϵ^2}{2}})\sigma _{,m}\sigma _{,n}`$ (160)
$`+`$ $`(1ϵ)\left[\sigma _{,m;n}{\displaystyle \frac{ϵ}{2}}\stackrel{4}{g}{}_{mn}{}^{}\sigma _{;a}^{,a}\right].`$
Here we used the equation
$`P_\mu ^\alpha P_\nu ^\beta \sigma _{,\beta ;\alpha }=\underset{\mu }{\overset{4}{}}\underset{\nu }{\overset{}{G}}+{\displaystyle \frac{ϵ}{2}}[2\underset{\nu }{\overset{}{G}}\underset{\mu }{\overset{}{G}}(G^\varrho \underset{\varrho }{\overset{}{G}})P_{\mu \nu }].`$ (161)
By means of the expression
$`\stackrel{5}{R}=(P^{\mu \nu }+s^\mu s^\nu )\stackrel{5}{R}_{\mu \nu }`$ (162)
and taking into account the intermediate formulas
$`P^{\mu \nu }\stackrel{5}{R}{}_{\mu \nu }{}^{}=e^{ϵ\sigma }\stackrel{4}{g}{}_{}{}^{mn}\underset{mn}{\overset{}{\stackrel{5}{R}}},`$ (163)
$`s^\mu s^\nu \stackrel{5}{R}{}_{\mu \nu }{}^{}=\sigma _{;\alpha }^{,\alpha }\omega _{\lambda \alpha }\omega ^{\lambda \alpha },`$ (164)
$`\sigma _{;\alpha }^{,\alpha }=e^{ϵ\sigma }\left[\sigma _{;i}^{,i}+(1ϵ)\sigma ^{,i}\sigma _{,i}\right],`$ (165)
we find the result
$`\stackrel{5}{R}=e^{ϵ\sigma }(\stackrel{4}{R}+e^{ϵ\sigma }\stackrel{4}{\omega }{}_{mn}{}^{}\stackrel{4}{\omega }{}_{}{}^{mn}(23ϵ)\sigma _{;m}^{,m}+2(1ϵ+{\displaystyle \frac{3}{4}}ϵ^2)\sigma ^{,m}\sigma _{,m}).`$ (166)
Now we immediately can obtain the 4-dimensional field equations of PUFT being restricted to the versions II and III of PUFT. Today the version I of PUFT has only historical value.
### 3.3 Version II
By projecting the 5-dimensional field equations (1) onto the 4-dimensional space-time with the help of the projection formalism developed above we obtain the 4- dimensional field equations of PUFT. As the space-time metric $`\stackrel{\mathrm{𝟒}}{g}`$ is connected with the induced metric $`\stackrel{}{𝒈}`$ on the hypersurface $`_4`$ by means of (98), in case of the version II of PUFT it is necessary to put:
$`\begin{array}{ccc}G_{\mu \nu }\hfill & =& \underset{\mu \nu }{\overset{5}{R}}\frac{1}{2}g_{\mu \nu }\stackrel{5}{R}+\lambda _0S_0e^\sigma \left(g_{\mu \nu }+s_\mu s_\nu \right)\hfill \\ ϵ\hfill & =& 1\hfill \end{array}\}.`$ (169)
#### 3.3.1 Generalized gravitational field equation
Using the last results obtained from the equation (96a) within the framework of(169), the field equations read:
$`\stackrel{4}{R}{}_{mn}{}^{}{\displaystyle \frac{1}{2}}\stackrel{4}{g}{}_{mn}{}^{}\stackrel{4}{R}+\lambda _0S_0\stackrel{4}{g}{}_{mn}{}^{}=\kappa _0\underset{mn}{\overset{4}{T}},`$ (170)
where
$`\underset{mn}{\overset{4}{T}}=\underset{mn}{\overset{4}{\theta }}+{\displaystyle \frac{1}{\kappa _0}}[2e^\sigma (\stackrel{4}{\omega }{}_{ma}{}^{}\stackrel{4}{\omega }{}_{n}{}^{a}+{\displaystyle \frac{1}{4}}\stackrel{4}{g}{}_{mn}{}^{}\stackrel{4}{\omega }{}_{ab}{}^{}\stackrel{4}{\omega }{}_{}{}^{ab})\text{}`$
$`{\displaystyle \frac{3}{2}}(\sigma _{,m}\sigma _{,n}{\displaystyle \frac{1}{2}}\stackrel{4}{g}{}_{mn}{}^{}\sigma _{}^{,a}\sigma _{,a})]`$ (171)
and
$`\underset{mn}{\overset{4}{\theta }}g_m^\mu g_n^\nu \underset{\mu \nu }{\overset{}{\stackrel{5}{\theta }}}.`$ (172)
#### 3.3.2 Generalized electromagnetic field equations (Maxwell equations)
Comparing the relation (3.3.1) with the expression (5), we find that the angular velocity of the congruence (23) $`\omega _{\alpha \sigma }`$ is connected with the electromagnetic strength tensor in the following way:
$`B_{\alpha \sigma }\underset{\alpha \sigma }{\overset{}{B}}=B_0e^{a\sigma }\omega _{\alpha \sigma },`$ (173)
where the constant $`B_0`$ depends on the system of units. We choose the constant $`a`$ in order to fulfill the next equation
$`B_{<\mu \nu \left|\right|\alpha >}=0.`$ (174)
It is easy to see that the relation
$`B_{<\mu \nu \left|\right|\alpha >}=B_{<\mu \nu ;\alpha >}`$ (175)
holds. Using the expression (36) and the equation
$`X_{<\mu \nu ;\alpha >}=0`$ (176)
we find
$`B_{<\mu \nu ;\alpha >}=B_0e^{a\sigma }(1+a)\omega _{<\alpha \mu }\sigma _{,\nu >}.`$ (177)
This implies $`a=1`$ and
$`B_{\alpha \sigma }=B_0e^\sigma \omega _{\alpha \sigma },B_{mn}g_m^\mu g_n^\nu B_{\mu \nu }=B_0e^\sigma \stackrel{4}{\omega }{}_{mn}{}^{}.`$ (178)
It is obvious (see (102) and (105)) that the electromagnetic field strength tensor satisfies the cyclic Maxwell system
$`B_{<mn;k>}=0.`$ (179)
By substituting (178) in the expression (3.3.1) we find that the electromagnetic induction tensor is to be defined as follows:
$`H_{mn}=e^{3\sigma }B_{mn},`$ (180)
and the constant $`B_0`$ can be chosen as
$`B_0=\pm \sqrt{{\displaystyle \frac{8\pi }{\kappa _0}}}.`$ (181)
In this case the electromagnetic part of the energy-momentum tensor $`\underset{mn}{\overset{4}{T}}`$ (3.3.1) takes its usual form (in the Gaussian system of units):
$`E_{mn}={\displaystyle \frac{1}{4\pi }}(B_{mk}H_n^k+{\displaystyle \frac{1}{4}}\stackrel{4}{g}{}_{mn}{}^{}B_{jk}^{}H^{jk}).`$ (182)
It is easy to see that the one-form
$`A_\mu B_0S_0P_\mu ^\sigma \zeta _\sigma =B_0S_0P_\mu ^\sigma \tau _{,\sigma }`$ (183)
has the following properties:
$`A_{\mu \left|\right|\nu }A_{\nu \left|\right|\mu }=B_{\nu \mu }\text{and}A_{m,n}A_{n,m}=B_{nm}\text{(}A_m=g_m^\sigma A_\sigma \text{)}.`$ (184)
Thus the orthogonal vector, projected into the hypersurface $`_4`$ in an appropriate way, is the electromagnetic vector potential.
Now we are ready to expound the equation (96b). The result is
$`H_{;n}^{mn}={\displaystyle \frac{4\pi }{c}}j^m,`$ (185)
where the abbreviations we used are given by
$`j^m={\displaystyle \frac{\kappa _0B_0c}{4\pi }}\stackrel{4}{\theta }{}_{}{}^{m}\text{and}\stackrel{4}{\theta }{}_{}{}^{m}\stackrel{4}{g}{}_{mn}{}^{}g_{n}^{\mu }\underset{\mu }{\overset{}{\stackrel{5}{\theta }}}.`$ (186)
#### 3.3.3 Field equation of the scalaric field $`\sigma `$
Using the Relations (164) and (92) we can rewrite the equation (96c) in the form
$`\sigma _{;m}^{,m}={\displaystyle \frac{\kappa _0}{8\pi }}B_{mn}H^{mn}+{\displaystyle \frac{2}{3}}\kappa _0\vartheta ,`$ (187)
where the following definition was used:
$`\vartheta \stackrel{5}{\theta }e^\sigma {\displaystyle \frac{1}{2}}\stackrel{4}{g}{}_{}{}^{mn}\underset{mn}{\overset{4}{\theta }}.`$ (188)
### 3.4 Version III
In the version III of PUFT the space-time metric $`\stackrel{\mathrm{𝟒}}{g}`$ coincides with the on the hypersurface $`_4`$ induced metric $`\stackrel{}{𝒈}`$ (see (97)). This makes the projection formalism a little bit easier, because $`ϵ=0`$. On the contrary the 5-dimensional field equations (1) are more complicated. Further investigations have shown that the case $`\kappa _0K_0=2`$, already mentioned above, is of a particular interest. Hence for the version III we find the following equation
$`\begin{array}{ccc}G_{\mu \nu }\hfill & =& R_{\mu \nu }\frac{1}{2}g_{\mu \nu }\stackrel{5}{R}\frac{1}{S}S_{,\mu ;\nu }\frac{2}{S^2}S_{,\mu }S_{,\nu }\hfill \\ & & \frac{1}{S}s_\mu s_\nu (4S_{;\tau }^{,\tau }\frac{6}{S}S_{,\tau }S^{,\tau }+\frac{3\lambda _0}{S}\frac{S}{2}\stackrel{5}{R})+\frac{1}{S}g_{\mu \nu }(S_{;\tau }^{,\tau }+\frac{\lambda _c}{S})\hfill \\ ϵ\hfill & =& 0\hfill \end{array}\}.`$ (192)
Following the above introduced procedure of deducing the field equations, one obtains the system of equations listed below .
#### 3.4.1 Generalized gravitational field equation
$`\stackrel{4}{R}{}_{mn}{}^{}{\displaystyle \frac{1}{2}}\stackrel{4}{g}{}_{mn}{}^{}\stackrel{4}{R}+{\displaystyle \frac{\lambda _0}{S_0^2}}e^{2\sigma }\stackrel{4}{g}{}_{mn}{}^{}=\kappa _0\underset{mn}{\overset{4}{T}},`$ (193)
where $`\underset{mn}{\overset{4}{T}}=\theta _{mn}+E_{mn}+S_{mn}`$ with
$`E_{mn}={\displaystyle \frac{1}{4\pi }}(B_{mk}H_n^k+{\displaystyle \frac{1}{4}}\stackrel{4}{g}{}_{mn}{}^{}B_{jk}^{}H^{jk}),S_{mn}={\displaystyle \frac{2}{\kappa _0}}(\sigma _{,m}\sigma _{,n}{\displaystyle \frac{1}{2}}g_{mn}\sigma _{,k}\sigma ^{,k})`$ (194)
holds. $`\underset{mn}{\overset{4}{\theta }}g_m^\mu g_n^\nu \underset{\mu \nu }{\overset{}{\stackrel{5}{\theta }}}`$ is the energy-momentum tensor of the substrate.
#### 3.4.2 Generalized electromagnetic field equations (Maxwell equations)
$`a)H_{;n}^{mn}={\displaystyle \frac{4\pi }{c}}j^m,b)B_{[mn,k]}=0,c)H_{mn}=e^{2\sigma }B_{mn},`$ (195)
where we used the abbreviations
$`\text{a)}B_{mn}g_m^\mu g_n^\nu B_{\mu \nu }=B_0e^\sigma \stackrel{4}{\omega }{}_{mn}{}^{},\text{b)}j^m={\displaystyle \frac{\kappa _0B_0c}{4\pi }}\stackrel{4}{\theta }{}_{}{}^{m}e_{}^{\sigma }`$ (196)
and $`\stackrel{4}{\theta }{}_{}{}^{m}\stackrel{4}{g}{}_{mn}{}^{}g_{n}^{\mu }\underset{\mu }{\overset{}{\stackrel{5}{\theta }}}`$, $`B_0=\pm \sqrt{\frac{8\pi }{\kappa _0}}`$.
#### 3.4.3 Field equation of the scalaric field $`\sigma `$
$`\sigma _{;m}^{,m}={\displaystyle \frac{\kappa _0}{16\pi }}B_{mn}H^{mn}{\displaystyle \frac{\kappa _0}{2}}\vartheta {\displaystyle \frac{\lambda _0}{S_0^2}}e^{2\sigma }\text{with}\vartheta \stackrel{5}{\theta }\stackrel{4}{g}{}_{}{}^{mn}\underset{mn}{\overset{4}{\theta }}.`$ (197)
## 4 Concluding Notes
Now let us summarize the basic ideas of the new geometrical approach to the axiomatics of Schmutzer’s 5-dimensional Projective Unified Field Theory.
The mathematical basis for the 5-dimensional Projective Unified Field Theory forms the group of all 5-dimensional homogeneous coordinate transformations of degree one (19). The 5-dimensional geometry, constructed on this group, supposes the existence of a Killing vector field. The integral curves of this vector field form a Killing congruence (23) which is the basis of the projection formalism developed here. The angular velocity $`\omega _{\mu \nu }`$ of this congruence is interpreted as the electromagnetic field strength tensor (see (178) and (195)). It is well known that, if $`\omega _{\mu \nu }=0`$ holds, a hypersurface, holonomic and orthogonal to the congruence exists. There are two possibilities to construct an axiomatics of PUFT: abandoning either holonomicity or orthogonality. The first of the two possibilities was investigated in detail in numerous papers by Schmutzer (see and there quoted papers). The second possibility was considered in the present paper. In this case it is possible to say that the holonomicity of space-time and the non-orthogonality of the given congruence with respect to the space-time hypersurface are embodied in the basis of the axiomatics offered here. In this way PUFT has got a new geometrical interpretation.
Inner curvature of space-time ($`_4`$) identified with the hypersurface $`\tau =\text{const}`$ (see page 2) describes the gravitation. The norm of the Killing vector field $`\xi `$ (24) is connected with the new scalaric field $`\sigma `$: $`\sqrt{\xi ^\nu \xi _\nu }=S_0e^\sigma `$. The tensor of the angular velocity $`\omega _{\mu \nu }`$ (30) of the Killing congruence (23) describes the electromagnetic field. Thus the orthogonal vector projected in an appropriate way onto the hypersurface $`\tau =\text{const}`$ ($`_4`$) is the electromagnetic vector potential (see (183) and (184)). The relation (183) implies that the electromagnetic potential vanishes if the hypersurface $`_4`$ is orthogonal to congruence (23).
It is easy to show that physically the 4-dimensional field equations in the version II of PUFT differ slightly from the corresponding equations in the version III physically slightly. Here we won’t dwell on this problem, therefore let us just remark that the cosmological term in the equation (1) can be accepted in the form $`\mathrm{\Lambda }_{\mu \nu }=\lambda _0e^\sigma P_{\mu \nu }`$ ($`\mathrm{\Lambda }_{\nu ;\mu }^\mu =0`$) as well. In this case additional terms containing a cosmological constant in the equations (170) and (187) take the form $`\lambda _0e^{2\sigma }\stackrel{4}{g}_{mn}`$ and $`\frac{4}{3}\lambda _0e^{2\sigma }`$, respectively.
In conclusion let us emphasize that the axiomatics constructed here leads to the same 4-dimensional field equations which formerly were obtained by E.Schmutzer in a different way.
## 5 Appendix. Results of Application of PUFT
Since 1995 a series of papers by E.Schmutzer on a closed homogeneous isotropic cosmological model of the universe and on the influence of the expansion of such a model on cosmogony and astrophysics appeared ( see where further literature is quoted) or are in press . Let us mention some main results:
* In order to be in agreement with the equivalence principle the usual concept of mass is basically changed: mass depends on the cosmological scalaric field. Hence follows a considerable change of the cosmological situation at the start of the universe (fulfilling of certain aspects of Mach’s principle).
* The big bang singularity does not exist. The ”big start” (Urstart) of the universe begins softly and is (using a certain physically motivated choice of parameters) characterized by a kind of oscillations: expansion interrupted by small contractions.
* The cosmological scenario appears to be divided into a short repulsion (antigravitational) era (duration of 128 years) and a cosmologically long attraction era (age of the universe = 18 billions of years).
* The Hubble factor (”constant”) is $`75\text{km}/s\text{Mpc}`$.
* Maxima and minima in the curves for the temporal behaviour of the cosmological mass density and the temperature could be interesting for the explanation of cosmogonic activities (birth of galaxies and stars).
* The equation of motion of a body is in full agreement with the Einstein effects (periastron motion, deflection and frequency shift of electromagnetic waves).
* Further consequences of the equation of motion are:
+ Time dependence of the ”effective gravitational constant” with the present relative value: $`3.510^{11}/\text{year}`$.
+ For an orbiting body around a center: positive value of the angular (secular) acceleration, negative values of the time derivatives of the orbital radius (decrease), revolution period (decrease), excentricity (transition from elliptic to cyclic orbits).
+ Heat production in a moving body with application to the moon, planets, sun, galaxy etc. with remarkably interesting numerical results.
The author would like to thank Professor Ernst Schmutzer for numerous helpful discussions on the axiomatics of PUFT.
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# The Propagation of Magneto-Centrifugally Launched Jets: I
## 1 INTRODUCTION
Highly collimated supersonic jets are a ubiquitous phenomena occurring in many astrophysical environments. These jets are observed propagating from sources as diverse as Active Galactic Nuclei (AGN, Leahy (1991)), Young Stellar Objects (YSOs, Reipurth (1997)) and Planetary Nebulae (PNe, Soker & Livio (1994)). While considerable progress has been made in understanding the nature of jets from AGN and YSOs, there remains considerable debate concerning the nature of the more recently discovered PNe jets.
The ubiquity of jets in astrophysics has made them a popular subject for study. They are excellent laboratories for the study of basic astrophysical processes (shocks, instabilities, etc.). Their long dynamical or “look-back” times, $`t_{dyn}=L_j/V_j`$, also make them ideal astrophysical fossils for studying the evolution of the obscured and often unobservable central sources, i.e., there is the hope in jet studies that the physics of the central engine can be revealed by studying the exhaust. Given the diversity of jet producing environments there also exists the hope that an underlying unity can be be found in terms of the fundamental processes which create jets. Articulating these processes is one of the critical issues facing astrophysical jet studies.
Accretion disks are believed to play a key role in the physics of both YSOs and AGN. In-falling, rotating matter is stored in these disks until dissipation allows material to spiral inward and feed the central, gravitating object. Both YSO and AGN disks are believed to support strong, well ordered magnetic fields. The current consensus holds that these fields are the agents for producing jets in a process known as Magneto-centrifugal launching. In this mechanism, plasma in the disk is loaded on to co-rotating field lines. If conditions in the disk are favorable (i.e., field strength and orientation) the plasma is centrifugally flung outward along the field lines. Strong toroidal field components are generated in the flow as the field is dragged backwards by the plasma inertia leading to collimation of the wind into into a narrow jet. We note, however, that the external medium might also help focus the outflow. Magneto-centrifugal launching has been studied in detail by many authors both analytically (Heyvaerts & Norman (1989), Pudritz (1991), Shu et al. (1994), ostriker (1997), Lery et al. (1999)) and through numerical simulations (Ouyed & Pudritz (1997), Romanova et al. (1998), Kudoh et al. (1998)).
In the YSO community two principle flavors of the Magneto-centrifugal launching model exist. The first is a pure disk wind model (Pudritz (1991)) in which the jet is generated at the surface of a Keplerian disk. The second, called “X-Winds” (Shu et al. (1994)), produces a jet from the boundary layer between the disk and the central star’s magnetosphere. Other models exist as well (Goodson et al. (1997)) and there remains considerable debate as too which mechanisms are obtained in real YSO flows.
While there is an exhaustive literature concerning jet launching and collimation, there has also been considerable study of jet propagation. Propagation studies focus on scales many orders of magnitude larger (Reipurth (1997)) than the region where collimation occurs. For example in the work of (Ouyed & Pudritz (1997)) the collimation of the jet was followed out to a height above the disk of $`H=80R_i`$ where $`R_i`$ is the inner disk radius. Since $`R_i10R_{}`$ ($`R_{}`$ is the stellar radius, Hartmann (1998)), the scale of the simulation was at least 10 times smaller than the smallest scales on which jets have been resolved and at least $`10^3`$ times smaller than the typical scale of observational jet studies. Much of the propagation work has been numerical and for both YSOs and AGN much of it has been have been purely hydrodynamic. For YSOs only a handful of MHD studies of jet propagation have been carried out to date (Todo et al. (1992), Cerqueira et al. (1998), Frank et al. (1998), Cerqueira et al. (1999), Gardiner et al. (1999), Stone & Hardee (1999)). If, however, strong magnetic forces produce the jets then these forces should effect their propagation downstream. Unless the fields are somehow removed, Maxwell stresses should alter at least some characteristics of the jet’s propagation. Recently Frank et al. (1999) have shown that ambipolar diffusion may be operative in YSO jets in some part of the flow. However the time-scales involved are such that changes in jet magnetic fields will only occur for parsec-scale jets. Flows on time-scales less than $`\tau 10^3`$ y will not lose their fields. In the case of AGN, the ambipolar time-scales are even larger. Thus, a proper accounting for the MHD forces in the propagation of both YSO and AGN jets is needed. In this paper we focus mainly on YSOs but our results will be applicable to AGN jets as well.
To date all radiative MHD jet simulations of steady, constant density “top-hat” jets have been performed using simple field geometries. Cerqueira et al. (1999) showed that jets with purely poloidal $`\stackrel{}{B}=B_z𝐤`$ topologies did not have propagation characteristics which differed significantly from pure hydrodynamic jets. Gardiner et al. (1999) have also found similar results for pulsed “top-hat” MHD jets with poloidal fields. Frank et al. (1998) however, found that if the field had a strong toroidal ($`B_\varphi `$) component then the jet head could be strongly effected by the Maxwell stresses leading to the production of so-called “nose-cones”. Nose-cones form when post-shock gas is restricted from lateral expansion by the axially directed “hoop stresses” associated with strong toroidal fields. Instead of back-flowing to form a cocoon, the shocked gas is confined to the head of the beam in the region downstream of the jet-shock.
The hoop stresses lead to a conical streamlined configuration for the head i.e. a nose-cone. Such structures were also seen in the early MHD simulations of AGN jets (Lind et al. (1989)). In Frank et al. (1998) the addition of radiative losses, appropriate for YSO jets, caused the nose-cones to narrow significantly. In a more extensive set of calculations Stone & Hardee (1999) found that MHD effects on jet propagation is strongly dependent on initial field topology.
While these results were promising, there still remains considerable distance to be traveled in the study of MHD jets. The principle issue that must be addressed is that all the simulations carried out to date is the use of ad-hoc field topologies. Unless a force-free configuration is adopted, $`J\times 𝐁=0`$, Maxwell stresses will act on the jet beam independent of propagation effects. Thus some effort must be expended in developing equilibrium configurations for MHD jet simulation initial conditions. With little to guide them, all modelers have chosen simple topologies which allow for a simple specification of the required equilibrium. Frank et al. (1999) used a pure toroidal geometry. Gardiner et al. (1999) used a pure poloidal geometry. Cerqueira et al. (1999) used both toroidal and poloidal as well as force free helical configurations which had to extend throughout the entire computational domain (jet + ambient medium). Stone & Hardee (1999) used helical pressure matched beams in a variety of configurations. None of the configurations used in these papers deviated from the simple constant velocity, constant density model for the jet beam. These efforts were necessary for articulating the basic role of MHD forces in jets, but they do not help establish a connection between conditions in the jet and the protostellar source (a protostar and rotating magnetized accretion disk). What is needed for use by the broader community is to begin the simulations with jet cross-sections derived directly from magneto-centrifugal flow models. That is the goal of the work presented here.
In what follows we present models of MHD jet propagation with initial configurations in the jet taken directly from the solution of force balance perpendicular (the Grad-Shafranov equation) and parallel (the Bernoulli equation) to magnetic surfaces generated by a magnetized rotator. Our simulations follow the evolution of jets composed of helical fields embedded in hypersonic plasmas whose density and velocity vary with radius. Thus our models constitute a further step towards realism in the theoretical description of magnetized astrophysical jets. The goal of this paper is to articulate the basic physics which can occur in these kinds of jets and to look for differences between the propagation of jets forming from different kinds of rotators. We note that the parameter space of solutions is quite large and in this paper we present only the first results of this project. In future papers we will present a more systematic exploration of parameter space.
The plan of the paper is as follows. In section 2 we describe the methods used to construct the initial equilibria and numerically simulate the flows. In section 3 we present results of our simulations focusing on adiabatic, isothermal and radiative cases. The next section compares the results with observations. Finally, in section 5 we present and discuss our conclusions.
## 2 Numerical Methods and Initial Equilibria
### 2.1 Basic Equations
We numerically integrate the equations of ideal magnetohydrodynamics (MHD), modified to include the loss of thermal energy due to optically thin radiative losses. In cylindrical coordinates these equations take the following form,
$`{\displaystyle \frac{\rho }{t}}+{\displaystyle \frac{1}{r}}{\displaystyle \frac{}{r}}(r\rho v_r)+{\displaystyle \frac{}{z}}(\rho v_z)`$ $`=`$ $`0`$ (1)
$`{\displaystyle \frac{\rho v_r}{t}}+{\displaystyle \frac{1}{r}}{\displaystyle \frac{}{r}}(r\rho v_r^2rB_r^2)+{\displaystyle \frac{}{r}}(p^{})+{\displaystyle \frac{}{z}}(\rho v_rv_zB_rB_z)`$ $`=`$ $`{\displaystyle \frac{(\rho v_\varphi ^2B_\varphi ^2)}{r}}`$ (2)
$`{\displaystyle \frac{\rho v_\varphi }{t}}+{\displaystyle \frac{1}{r}}{\displaystyle \frac{}{r}}(r\rho v_\varphi v_rrB_\varphi B_r)+{\displaystyle \frac{}{z}}(\rho v_\varphi v_zB_\varphi B_z)`$ $`=`$ $`{\displaystyle \frac{(B_rB_\varphi \rho v_\varphi v_r)}{r}}`$ (3)
$`{\displaystyle \frac{\rho v_z}{t}}+{\displaystyle \frac{1}{r}}{\displaystyle \frac{}{r}}(r\rho v_zv_rrB_zB_r)+{\displaystyle \frac{}{z}}(\rho v_z^2B_z^2+p^{})`$ $`=`$ $`0`$ (4)
$`{\displaystyle \frac{B_r}{t}}+{\displaystyle \frac{}{z}}(v_zB_rv_rB_z)`$ $`=`$ $`0`$ (5)
$`{\displaystyle \frac{B_\varphi }{t}}+{\displaystyle \frac{1}{r}}{\displaystyle \frac{}{r}}(rv_rB_\varphi rv_\varphi B_r)+{\displaystyle \frac{}{z}}(v_zB_\varphi v_\varphi B_z)`$ $`=`$ $`{\displaystyle \frac{(v_rB_\varphi v_\varphi B_r)}{r}}`$ (6)
$`{\displaystyle \frac{B_z}{t}}+{\displaystyle \frac{1}{r}}{\displaystyle \frac{}{r}}(rv_rB_zrv_zB_r)`$ $`=`$ $`0`$ (7)
$`{\displaystyle \frac{E}{t}}+{\displaystyle \frac{1}{r}}{\displaystyle \frac{}{r}}[r(E+p^{})v_rrB_r(\stackrel{}{B}\stackrel{}{v})]+{\displaystyle \frac{}{z}}[(E+p^{})v_zB_z(\stackrel{}{B}\stackrel{}{v})]`$ $`=`$ $`\left({\displaystyle \frac{\rho }{\mu }}\right)^2\mathrm{\Lambda }(T)`$ (8)
The total energy and pressure are given by
$`p^{}=p+{\displaystyle \frac{1}{2}}B^2`$ (9)
$`E={\displaystyle \frac{1}{2}}\rho v^2+{\displaystyle \frac{p}{\gamma 1}}+{\displaystyle \frac{1}{2}}B^2`$ (10)
where $`\mu `$ is the mean molecular weight and $`B^2=\stackrel{}{B}\stackrel{}{B}`$. Equation 1 and 8 represents conservation of mass and energy respectively. Equations 2 \- 4 represent conservation of momentum. Equations 5 \- 7 represent the induction equation. The energy conservation equation includes a source term, $`n^2\mathrm{\Lambda }(T)`$, (where $`n=\rho /\mu `$ is the number density), which models radiative losses in the optically thin limit. We use the Dalgarno-McCray “coronal” cooling curve (Dalgarno et al. (1972)). A “floor” temperature of $`T_f=10^4`$ K is set such that gas can not cool to lower values. The fluid is assumed to be an ideal gas, where $`\gamma `$ is the ratio of specific heats. Other relevant quantities are the sound speed $`c=\sqrt{\gamma p/\rho }`$, the Alfvén speed parallel to the magnetic field $`v_a=\sqrt{B^2/\rho }`$ and the plasma beta parameter, $`\beta =2p/B^2`$.
In addition to the hyperbolic equations represented above an additional constraint is imposed via the condition of flux conservation,
$$\frac{1}{r}\frac{}{r}B_r+\frac{}{z}B_z=0.$$
(11)
Using these equations we model the propagation of a magnetized jet through a constant density, constant pressure magnetized ambient medium. The initial conditions for the jet, i.e., its cross sectional distribution of $`\rho ,p,\stackrel{}{v}`$ and $`\stackrel{}{B}`$, are calculated via the Given Geometry Method of (Lery et al. (1998), Lery et al. (1999), Lery & Frank (1999)). We describe this method and the equilibrium MHD jet solutions it produces in the next section.
#### 2.1.1 The Model
The jets we inject into the computational grid are taken directly from a (simplified) model of the magneto-centrifugal launching/collimation process. The model, known as the Given Geometry Method (GGM: Lery et al. (1998), Lery et al. (1999), Lery & Frank (1999)) allows asymptotic MHD jet equilibria to be linked directly to the properties of a rotating source. The GGM assumes a time-independent, axisymmetric flow. It further simplifies the problem of magneto-centrifugal launching/collimation by assuming that the nested magnetic flux surfaces defining the flow (labeled by the variable $`a`$) possess a shape which is known a priori inside the fast critical surface. The fast surface defines the locus of points beyond which the flow is kinetic energy dominated. The flux surfaces are assumed to be conical and, as an additional simplification, an equilibrium across the surfaces is assumed at the Alfvén point which yields an equation referred to as the Alfvén regularity condition. This condition is not a criticality condition since the Alfvèn point is not strictly a critical point.
The flow properties must be determined by solving for the equilibrium of forces parallel and perpendicular to the magnetic surfaces (the former described by using the Bernoulli equation for a polytropic equation of state and the latter is solved via the Grad-Shafronov equation). The equilibrium parallel to the surfaces takes the form of criticality conditions at the two other (fast and slow) MHD critical points. This corresponds to differential form of the Bernoulli equation on constant $`a`$ with respect to $`\rho `$ and $`r`$ vanishing at the critical points.
In the general case the GGM yields five integrals of motion that are preserved on any axisymmetric magnetic surface $`a`$. Two of the integrals are given as boundary conditions in the model. These are the angular velocity $`\mathrm{\Omega }(a)`$ and an entropy (or ploytropic) factor $`Q(a)`$. These are supplied as a model for the source rotator. We note that the entropy paramter Q(a) can be decribed as follows: The density $`\rho `$ is related to the pressure $`p`$ by a polytropic equation of state, $`p=Q(a)\rho ^\gamma `$ where $`\gamma `$ is the polytropic index and Q the polytropic constant that is related to the entropy. This assumption replaces consideration of energy balance and is meant to simply represent more complex heating and cooling processes (See, for example, Vlahakis & Tsinganos (1998) for more general equations of state). Changing Q(a) changes the local thermal energy balance in the flow.
The Alfvén regularity condition together with the criticality conditions then determine the three other unknown integrals: namely the specific energy $`E(a)`$; the specific angular momentum $`L(a)`$; the mass to magnetic flux ratio $`\alpha (a)`$. Far from the source (large z) the flow becomes cylindrically collimated. In this asymptotic regime the jet is assumed to be in pressure equilibrium with an external medium. The pressure matching condition along with with the Grad-Shafronov and Bernoulli equations are all solved in the asymptotic cylindrically collimated regime.
We note that in the GGM the source (e.g. the accretion disk) is not explicitly described since it is point-like. Instead the shape of the magnetic field lines defined by the flux function a(r,z) is specified out to the fast magnetosonic point, but not its angular distribution. The rotation rate $`\mathrm{\Omega }(a)`$ and the polytropic parameter $`Q(a)`$ are then specified on the field lines. The shape of the magnetic field lines from the fast magnetosonic point to the fully collimated region is not specified. The strong toroidal component which develops in the wind and the fully collimated jet develops mainly due to differential rotation and the interia of mass on the field lines. This is similar to other disk wind models Ouyed & Pudritz (1997).
#### 2.1.2 Numerical Solutions
Inside the fast critical surface, the variables calculated in the numerical procedure are the energy $`E`$ and the radii and densities at the three critical surfaces, ($`r_s`$, $`r_f`$, $`r_A`$, $`\rho _s`$, $`\rho _f`$, $`\rho _A`$). In the asymptotic cylindrically collimated regime, the jet is entirely defined by this set of $`r`$ and $`\rho `$ and all other physical quantities can be derived from them. For the numerical calculations, the equations have been reformulated as ODEs or converted from algebraic conditions into ODEs as functions of the flux surfaces $`a`$. The system consists of eight differential equations and the numerical solutions are obtained by initiating the integration of the system from the axis. Given the input parameters $`Q(a)`$, $`\mathrm{\Omega }(a)`$, $`\alpha _0`$, $`\gamma `$, and $`\rho _0`$, all the critical positions and densities can be numerically obtained using analytical formulae (see Lery et al. (1998)). We further constrain the solution to be super-Alfvénic and super-fast-magnetosonic on the axis in the asymptotic region.
#### 2.1.3 Classes of Jet Equilibria
In our approach the most important aspect of the source rotator is is angular rotation profile. We focus on this aspect of the source because it most clearly connects to different scenarios of magneto-centrifugal launching/collimation. Profiles of angular velocity of the source rotator considered in this paper are shown in Fig. The Propagation of Magneto-Centrifugally Launched Jets: I. The pure Keplerian rotator (dashed line) starts with a constant rotation close to the axis, as in the rigid body case (dot-dashed lines), but then follows a Keplerian profile. The Multi-component (solid line) case also starts with a rigid rotation corresponding, for example, to an axial ordinary wind. The angular velocity then doubles its value in order to model a jet rotating more rapidly than the star in an intermediate region between the ordinary wind and the Keplerian disc wind that follows. Note that the angular velocity is always sub-Keplerian in the intermediate region. For all the rotation laws, the axial value of the angular velocity $`\mathrm{\Omega }_0`$ is set to unity in the Figure, and the radius is normalized to the size of the jet. For reference we also show the profile of a solid body rotator though we will not consider this model in the simulations. Fig. The Propagation of Magneto-Centrifugally Launched Jets: I also shows that a return poloidal electric current flows back inside the jet for both the Keplerian and the Multi-component jets.
As shown in Lery & Frank (1999), it is possible to derive an approximate analytical solutions of the model in the cylindrical region. It has been found that the density can be expressed as a function of $`r`$ and of the first integrals as $`\rho (r)C\alpha (r)/\mathrm{\Omega }(r)r^2`$ where $`C`$ is a constant. The asymptotic poloidal velocity of the flow can also be derived and is given by $`v_z(r)\mathrm{\Omega }(r)r/C`$. Therefore the velocity increases with the angular velocity while the density decreases. This explains why the density drops when $`\mathrm{\Omega }`$ is important in the inner part of the jet for the Multi-component case, while it increases afterwards in the Keplerian rotation regime. The velocity roughly follows an opposite behavior with respect to the angular velocity. More detailed analysis of these equilibria are given by Lery & Frank (1999).
#### 2.1.4 Identification of Basic Features
The quantities that define the jet in the cylindrically collimated regime, are plotted in Fig. The Propagation of Magneto-Centrifugally Launched Jets: I for pure Keplerian, Multi-component and constant rotations. The Keplerian and Multi-component models will be used as input for the numerical simulations. The $`z`$ and $`\varphi `$ components of velocity and magnetic field are represented together with the density $`\rho `$ and the net electric current $`I_C`$, as functions of the relative radius (normalized to the jet radius). The length scale is the jet radius, the density is normalized to its value on the jet axis $`\rho _0`$, and the non-dimensional velocities refer to the fast magnetosonic velocity $`v_f^2=c_s^2+v_A^2`$ on the axis, $`c_s`$ being the sound speed. The magnetic field is normalized to $`\sqrt{\rho _0}v_f`$.
The most important features in these graphs are the variations of the toroidal component of the magnetic field and of the density. Note that the region near $`R=.3R_j`$ is dominated by the magnetic pinching force, or hoop stress, where $`B_\varphi `$ is maximum. The gas pressure is important at radii less than this value in order to maintain the equilibrium and is the origin of the large density gradients in this region. We denote the high density region centered on the axis as the core, and the lower density outer regions as the collar. Note that the bulk of the jet’s momentum resides in the core. Hence we expect this portion of the beam to penetrate more easily into the ambient medium during the jet’s propagation while the collar will be more strongly decelerated. Fig. The Propagation of Magneto-Centrifugally Launched Jets: I also shows that a return poloidal electric current flows back inside the jet for both the Keplerian and the Multi-component jets. More detailed analysis of these equilibria are given by Lery & Frank (1999).
Thus Keplerian and Multi-component jets are characterized by a dense, current-carrying core, carrying most of the momentum, and are surrounded by a collar carrying an internal return current.
#### 2.1.5 Scaling for the Simulations
The input parameters of the model can be selected so as to qualitatively reproduce observed situations. Given the properties of the jet-emitting object, i.e., its radius $`R_{}`$, its temperature $`T_{}`$, the total mass loss rate $`\dot{M}_{}`$, the base density $`n_{}`$, the magnetic field $`B_{}`$, the factor $`Q_{}`$ and $`\gamma `$, it is possible to deduce the dimensionless parameters $`\overline{\mathrm{\Omega }}`$, $`\overline{Q}`$, $`\overline{\alpha }_0`$. The parameter $`\overline{\alpha }_0`$ can be a-posteriori related to the mass loss rate $`\dot{M}_{}`$, $`R_{}`$, and the magnetic field $`B_{}`$. So we define $`Q_{}2kT_{}n_{}/(m_pn_{})^\gamma `$, $`\alpha _{}\dot{M}_{}/4\pi R_{}^2B_{}`$, and $`\mathrm{\Omega }_{}\sqrt{GM_{}/R_{}^3}`$. All those quantities are non-dimensionalized to reference values by setting $`\overline{Q}Q_{}/Q_{ref}`$, $`\overline{\alpha }_0\alpha _{}/\alpha _{ref}`$ and $`\overline{\mathrm{\Omega }}\mathrm{\Omega }_{}/\mathrm{\Omega }_{ref}`$. The entropy $`\overline{Q}(a)`$ is assumed to be constant across the jet. In the present paper, we have chosen to model YSO jets with different rotation laws using typical values for TTauri stars as presented by Bertout et al. (1988). At the base of flow, we deduce the corresponding dimensionless input parameters: $`\overline{Q}=0.87`$, $`\overline{\mathrm{\Omega }}=2`$, $`\overline{\alpha }_0=0.7`$, and $`\rho _0=5.10^7`$. Major quantities of reference are then given (in CGS) by $`R_{ref}=10^{15}cm`$, $`n_{ref}=250cm^3`$, $`v_{ref}=10^7cms^1`$ for Young Stellar Objects. $`V_ref`$ is simply a canonical speed for YSO jets which we use to set the scales in the simulations.
#### 2.1.6 Comparisons with other models
A detailed comparison of the present model with previously published studies has been given by Lery & Frank (1999). Here we report only the most important conclusions. As with Ferreira (ferr3 (1997)), it is possible to show that the GGM yields a minimum mass loss rate injected in the jet which has a lower limit and can not be arbitrarily small. These results also agree with Ostriker (ostriker (1997)) and Lery et al. (1999b ) who conclude that the optical jet may represent only the densest part of the total outflow. We obtain a fast magnetosonic Mach number, (which also corresponds to the Alfvénic Mach number on the axis), between 2 and 4. This range corresponds to what (Camenzind (1997)) has found for his model for low-mass protostellar object. The corresponding jets have low fast magnetosonic Mach-numbers $`M_A2`$. By taking into account an accretion disc around the stellar magnetosphere, Fendt & Camenzind (fendt (1996)) also find a fast magnetosonic Mach-number to be $`2.5`$. This results does not, however, appear to be a general statement about MHD jets since there exist models with larger values (Sauty et al. (sautytsing (1994)) and Trussoni et al. (trussoni (1997))). Finally, the analytical results given by Shu et al. (shu95 (1995)) agree with those of the GGM model in terms of jet structure. We note however that despite the similarity of the analytical results, neither the Multi-component (or the Keplerian) case can be seen as equivilent to the X-wind model. Note in particular that our model does not describe the physical processes occurring at the source itself, i.e., at the surface of the disk or the disk-star boundary.
### 2.2 Numerical Method and Implementation
A detailed description of the numerical code can be found in references given below. Here we simply state the code’s most salient features. Specifically, the method we use to solve equations 1-8 is explicit, finite element (volume), up-winded, conservative, $`2^{\mathrm{nd}}`$ order accurate, and total variation diminishing (TVD). TVD stands for Total Variation Diminishing. This refers to the ability of the code to caputure strong discontinuities in the flow without producing spurious oscillations. TVD methods are part of a general class of ”High-Resolution” codes which solve the hydro or MHD equations in conservative form by using a Gudinov method (ie solving the Riemann problem at every grid interface) and including sophisticated algorithms for limiting the fluxes through cell boundaries to keep the solution monotone. More detail concerning High Resolution methods can be found in Leveque (1998). The code is conservative up to machine accuracy, ensuring that it will accurately capture shock strengths and speeds. It has been well tested in standard 1-D shock tube tests as well as multi-dimensional stability calculations. Various manifestations of the code have been reported in the literature including its 1-dimensional (1-D) cartesian form (Ryu & Jones (1995)), its 2-D cartesian form (Ryu et al. 1995a ), and its 2-D axisymmetric (cylindrical coordinates) form (Ryu et al. 1995b ). The TVD property is ensured in the same way as was done originally by Harten for the Euler equations in (Harten (1983)). In the 2-D versions of the code, multidimensionality is handled through the use of Strang splitting (Strang (1968)). The cooling is applied in a first order fashion. Finally, the crucial and problematic issue of maintaining $`\stackrel{}{}\stackrel{}{B}=0`$ is accomplished with a staggered grid approach (Ryu et al. (1998)).
Each simulation was carried out in cylindrical coordinates $`(r,\varphi ,z)`$ with axisymmetry and inversion symmetry across the $`z=0`$ plane. We follow a quarter meridional plane ($`r0`$, $`z0`$, $`\varphi =0`$) with $`512\times 2048`$ grid cells. The jet radius spans $`64`$ grid cells. Thus our simulations follow the jet propagation for $`Z=32R_j`$. In Table 1 we present a set of important physical parameters for the simulations presented below. While we will use dimentionless variables in most of the description that follows Table 1 allows the reader to compare the physical scales in the simulations with observations.
In what follows we express all distances in terms of $`R_j`$ and all times in terms of the magnetosonic crossing time $`\tau =R_j/M_f`$ where $`M_f=\sqrt{v_a^2+c^2}`$ (note: for our initial conditions the magnetosonic speed and the fast mode speed are identical). We utilize outflow boundary conditions at the outermost radial and axial boundaries. In cylindrical coordinates the $`r=0`$ line is necessarily a reflecting boundary. During the tests which we have run we have found little evidence of incorrect reflections from the $`r=0`$ line due to the coordinate singularity, though it must be admitted that this problem plagues all numerical codes in cylindrical coordinates with axisymmetry. Inversion symmetry through the $`z=0`$ plane dictates the use of reflecting boundary conditions at the $`z=0`$ plane. In some simulations we have found that waves propagating inward from the outermost radial boundary (caused by the bow-shock propagating off the grid) led to a slow compression of field at the base of the jet at late times. We wished to avoid the use of logarithmic grids hence we suppressed the inflow of material at the outer radial boundaries problem by injecting a slow ($`M_s1.2`$) wide angle flow at at the base of the grid. This advected the material off the grid and kept the field from being overly compressed near the jet inlet at the base of the grid and the imposed flow had no effect of the propagation of the jet far downstream.
In each simulation we inject the jet into the computational domain via 2 layers of “ghost-zones” below the base $`z=0`$ of the grid. The jet properties are read into the grid from data files provided by the Given Geometry Model described above. The density and pressure in the ambient medium are copied from values in the last radial zone of the jet: $`\rho _a=\rho _j(R_j),p_a=p_j(R_j)`$. In addition the ambient medium is given a pure poloidal magnetic field ($`\stackrel{}{B}_a=B_{a,z}\widehat{k}`$) whose magnitude is also taken from the last radial zone of the jet $`B_{a,z}=B_{j,z}(R_j)`$. Since $`B_{j,\varphi }(R_j)=0`$ our jets are in radial pressure balance the with the ambient medium.
In order to articulate the basic dynamics inherent to the flows we have run three classes of model for both the Keplerian and Multi-component jet. In our Adiabatic models we have set $`\gamma =5/3`$ and turned off the cooling source term. In our Isothermal models we have set $`\gamma =1.001`$ and turned off the cooling term. In our Radiative models $`\gamma =5/3`$ and the cooling source term was turned on. We have run both radiative and isothermal models as consistency checks as well as to allow us to model jets with different Mach numbers. For reasons explained above the equilibria provided by the Given Geometry model yielded jets of low magnetosonic number ($`3<M_f<5`$). We wish to model jets with lower temperatures. This could be accomplished by scaling down both the pressure and magnetic field such that the force balance was maintained. In this way we were able to model jets with Mach numbers of order ($`6<M_f<9`$).
The cross sectional variation of $`\stackrel{}{B}_j(r)`$ in the jet presents a problem in terms of initial conditions. This is a general difficulty which all attempts to model the evolution of magnetized jets must confront. Flux conservation, $`\stackrel{}{}\stackrel{}{B}=0`$, demands that any discontinuities in the field must be associated with current sheets (which are the cause of field kinks at MHD shocks). Attempts to initialize simulations of magnetized jets propagating into magnetized ambient media must deal with the likely mismatch of field topologies and magnitudes at the head and sides of the jet when the simulation is first switched on. The use of cylindrical coordinates eliminates the problem for the $`B_\varphi `$ component. While an initial discontinuity in the toroidal component may produce transients, it will not violate flux conservation. Thus we must only deal with the $`r`$\- and $`z`$-components of the field. In our simulations we solved the problem by continuing the $`z`$-component of the jet field into the ambient medium. Thus for $`r<R_j`$, $`B_{a,z}(r)=B_{j,z}(r)`$. Since $`B_z`$ is relatively weak, $`B_z^2/2p`$, the gradient in the ambient field produces little mass motion.
In order to test the effect of the initial conditions on the observed behavior (i.e., transients) we have run a series of models which began with the jet and ambient conditions joined smoothly via a hyperbolic tangent function. The smoothing length $`h`$ was varied from $`h=.5R_j`$ to $`h=9R_j`$. We found that the long term behavior of the jet was unaffected by the choice of h. We also note that this version of the code produces a relatively strong boundary layer at the jet/ambient gas interface at distances far from (well behind) the head of the jet. While such layers are to be expected due to unresolved instabilities (mainly Kelvin-Helmholtz modes) we found the effect was partially attributable to the treatment of transverse wave modes in the code. We performed a number of tests to ensure that changes in the flow variables in the boundary layer were not affecting the results.
## 3 Results
In the this section we present the results of the simulations. We provide a description of the behavior seen in the models along with attempts to understand the underlying physics.
### 3.1 Basic features
Several features are common to almost all of the simulations. In the input equilibria the core-collar structure is always present with gradients of jet variables between the jet core and collar. When the equilibrium jet encounters the external medium, the various elements of the equilibrium are shocked. This creates two bow-shocks and a cocoon. The bow-shock closest to the axis takes the form of a nose cone. Intrinsic instabilities develop in the inner part of the shocked core, as well as in the cocoon. The annotated Figure The Propagation of Magneto-Centrifugally Launched Jets: I sums up this section by showing the set of common features on a characteristic Multi-component jet. We now focus on the propagation characteristics of the two types of input equilibria.
### 3.2 Keplerian Rotator
To breifly review Keplerian rotators produce jets with a nearly constant velocity cross section. The mass density is stratified with a high density core surrounded by a lower density collar. The core-collar density ratio for the present jet is relatively low: $`\rho _{core}/\rho _{col}=1.67`$, (this is also the density ratio between the core and the ambient medium, $`\eta =\rho _j/\rho _a`$). The toroidal magnetic field in the jet reaches its maximum value just at the outer edge of the core. As we shall see, the coupling of higher density in the core with the strong magnetic stresses along the core/collar boundary dominates the propagation characteristics of the entire jet.
#### 3.2.1 Keplerian Rotator: Adiabatic Jet
##### Propagation
In Fig. The Propagation of Magneto-Centrifugally Launched Jets: I we present gray-scale maps of the density evolution of an adiabatic jet driven by a Keplerian rotator. In the first frame, taken at $`t=5.3\tau `$, the classic jet-shock/bow-shock pair are apparent. The bow-shock accelerates the ambient gas while the jet-shock decelerates the jet material. The speed of the jet head or bow-shock is $`v_h77`$ km $`\mathrm{s}^1`$. This speed is relatively constant throughout the simulations. Frank et al. (1998) derived a formula for the bow-shock speed which accounted for magnetic pressure in the beam. Using the familiar result for hydrodynamic jets, $`v_{ho}=\frac{v_j}{1+1/\sqrt{\eta \nu }}`$, ($`\nu `$ is the ratio of bow-shock and jet head radii $`R_h/R_j)`$ the MHD bow-shock speed is.
$$v_h=v_{ho}\frac{1\sqrt{\frac{1}{\eta \nu }\frac{p^{}}{\rho _jv_j^2}(1\frac{1}{\eta \nu })}}{1\frac{1}{\sqrt{\eta \nu }}},$$
(12)
If we take $`\nu =1`$ then this equation gives $`v_h70`$ km $`\mathrm{s}^1`$with magnetic pressure accounting for approximately 4% of the momentum flux driving the shock. As was noted in Frank et al. (1998), the higher velocity of the jet head seen the simulations can be attributed to the aerodynamic effect of streamlining the jet head via MHD hoop stresses (an $`\nu `$ *effect*). The nose-cone shape which develops reduces the drag on the jet head increasing its velocity relative to a more blunt jet head which would occur in a pure hydrodynamic simulation.
In the first frame of Fig. The Propagation of Magneto-Centrifugally Launched Jets: I we already see the effect of the core/collar structure on the jet propagation. The higher density core, confined by the magnetic hoop stresses, maintains its structural integrity on the downstream side of the jet-shock. It is noteworthy at this early time that the core appears to propagate ahead of the rest of the beam. This is to be expected purely from momentum considerations as equation 12 predicts a $`\mathrm{\Delta }v=10`$ km $`\mathrm{s}^1`$ difference in the speed of the core and collar bow-shocks. Detailed examination of the simulations also shows that at these early times there is little material flowing from the jet into the cocoon. This is apparent in Fig. The Propagation of Magneto-Centrifugally Launched Jets: I which shows the poloidal plane velocity vectors for the head of the jet. The origin of this effect lies, once again, in the relative strength of the toroidal field in the core and collar. What material does flow into the cocoon comes primarily from the outer most regions of the collar ($`r>.75R_j`$). Note that there is no transverse motion in the shocked core material. We attribute this to magnetic forces. In the region where material is flowing in the radial direction the magnetic field has an average value that is $`1/3`$ that at the core/collar interface. A better means of judging the relative strength of the field comes from examination of the plasma parameter $`\beta `$. Downstream of the jet-shock the core collar interface has $`\beta =3`$ while the collar/ambient interface has $`\beta =23`$. Thus fields can exert stronger stresses along the core restricting its lateral expansion.
The remaining frames of Fig. The Propagation of Magneto-Centrifugally Launched Jets: I show the distance between the jet and bow-shock continues to grow. Note that as the jet evolves the core never loses its identity. It acts, essentially, as a jet within a jet. At later times we see secondary shocks developing within the shocked core as well as vortex shedding at the head of the jet at radii consistent with the core/collar interface. In particular, by the second frame we see what appears to be a second jet-shock forming inside the core as it pushes through the ambient gas. Note also that at later times the shocked core material takes on the familiar nose-cone morphology seen in top-hat MHD jets with strong toroidal geometries. These features emphasize the apparent independence of the core’s propagation characteristics relative to the rest of the jet. Thus our results show that the core/collar dichotomy appears to control the main features of the jet beam via gradients in inertia and Maxwell stresses.
##### The Lateral Expansion
The second notable feature in the adiabatic Keplerian jet simulations is the large scale mass expulsion event which occurs in the third frame in Fig. The Propagation of Magneto-Centrifugally Launched Jets: I. The plasma driven laterally (in the $`r`$ direction) in this event is composed entirely of shocked collar material. At early times the cocoon is fed solely through the outer annuli of the jet as Fig. The Propagation of Magneto-Centrifugally Launched Jets: I demonstrated. The hoop stresses in the material in the collar at smaller radii are, however, too strong to allow plasma to stream transversely into the cocoon. To see this explicitly consider the radial momentum equation where terms involving $`v_\varphi `$, $`B_r`$ are ignored and we also ignore variations in $`z`$.
$$\frac{\rho v_r}{t}+\frac{1}{r}\frac{}{r}(r\rho v_r^2)=\frac{}{r}(p)\frac{}{r}\left(\frac{1}{2}B_z^2+B_\varphi ^2\right)\frac{(B_\varphi ^2)}{r}$$
(13)
The last term on the right is the hoop stress. The second to last term is the magnetic pressure. From Fig. The Propagation of Magneto-Centrifugally Launched Jets: I it is clear that for $`r>.3R_j`$ both the gas and magnetic pressure gradients are negative. Thus in these regions the hoop stress opposes the pressure forces and acts to constrain lateral expansion of the flow.
At later times however two features occur which alter the balance of forces. First as more material builds up immediately behind jet shock, both the gas pressure and magnetic pressure increase relative to the magnetic tension. Second, and most importantly, the jet-shock becomes distorted, tipping towards direction of jet propagation. The jet-shock becomes conical with the vertex of the cone pointing into the undisturbed beam. Thus the jet-shock becomes oblique relative to the un-shocked mass flux. The shock conditions for velocity lead to the following expression for the post-shock radial velocity
$$v_r\frac{1}{4}v_j\mathrm{cos}(\theta )\mathrm{sin}(\theta )$$
(14)
where $`\theta `$ is the angle between the jet-shock and the $`z`$ axis and $`v_j`$ is the jet velocity relative to the jet-shock. For $`\theta <90^o`$ the post-shock gas acquires a significant $`v_r`$ component. Material in the beam is refracted away from the axis. The momentum flux in the outward radial ram pressure, $`\rho v_r^2`$, is able to overwhelm the magnetic tension force leading to a large scale expulsion of material into the cocoon.
##### The Shocks
It is difficult to isolate the processes which cause the bending of the jet-shock into a conical shape. The dynamics at the jet head are highly non-linear and time-dependent and it is not obvious if the change in shock geometry is an amplification of events downstream where the flow pattern is quite complicated or if the distortion can be linked to events upstream. Close inspection of the simulations gives the impression that the distortion of the jet-shock occurs after a pinch wave reflects off the axis just upstream of the jet-shock in the un-shocked beam. The origin of the pinch appears to come from the slow expulsion of material into the cocoon prior to the third frame. In the early evolution the cocoon distorts the flow of ambient gas behind the bow-shock, i.e., the cocoon represents an obstacle which the post-bow-shock flow must stream around. As the shocked ambient material streams over the cocoon it becomes transonic. Its return to parallel streaming along the jet boundary can only occur via an additional shock. This feature is apparent at $`z8R_j`$ in frame 2 of Fig. The Propagation of Magneto-Centrifugally Launched Jets: I. Such flow patterns are well known to areodynamicists as they are common in aerofoil theory (Ramm (1990)). It appears that the pinch wave is generated just downstream of this additional shock and may be attributable to the higher pressures generated on the surface of the jet.
After the large mass expulsion event the jet-shock appears to relax to a configuration where it is perpendicular to the z-axis ($`\theta =90^o`$) and the flow into the cocoon is reduced. As the expelled material curls back towards the jet beam, however, it impinges on the jet surface and another, strong pinch is generated. At the end of the simulation (after the jet head has moved off the grid) we find the jet-shock becoming distorted yet again leading, perhaps, to a second mass-shedding event.
##### The Magnetic Effects
In Fig. The Propagation of Magneto-Centrifugally Launched Jets: I a we show the magnetic field structure in the adiabatic Keplerian jet. The Figure shows the poloidal ($`\stackrel{}{B}_p=B_r\widehat{e}_r+B_z\widehat{e}_z`$) magnetic field lines and the toroidal field ($`B_\varphi `$). The field clearly traces out the main features of the flow described above: the nose-cone at the jet head; the mass ejected behind the jet shock; the pinch wave occurring where the ejected mass is swept back onto the jet beam. It is noteworthy that it is the toroidal field which articulates these structures most clearly. This is appropriate as the toroidal field dominates in the jet providing much of the force which shapes the jet dynamics. Note that the Figure indicates that the magnetic pitch $`B_\varphi /|B_p|`$ increases behind shocks. This is a general feature of helical fields in jets. In a fast MHD shock only the component of the field parallel to the shock face is strengthened via compression (the parallel component will scale as $`\rho `$). Thus in a jet with a helical magnetic topology shock waves act to comb out the field leading to enhanced toroidal fields (and hoop stresses) in the post-shock regions (Gardiner & Frank (1999)). We note however that shear in the flow will also lead to significant strengthening of the field via stretching of field lines.
Our results for the Keplerian model show that the behavior of a jet with a more realistic initial density, pressure and magnetic field structure leads to propagation characteristics which have not been seen in previous hydrodynamic jet simulations. As we shall see this theme is repeated in all the simulations.
#### 3.2.2 Keplerian Rotator: Isothermal/Radiative Jet
##### The Propagation
In Fig. The Propagation of Magneto-Centrifugally Launched Jets: I we present gray-scale maps for the density evolution of the radiative Keplerian rotator jet. Recall that the radiative model begins with modified initial conditions compared with the adiabatic or isothermal Keplerian simulation. In order to keep the initial temperature at $`T=T_o10^4K`$ we scaled down all the radial distributions of all variables in the jet except $`v_z`$. This had the additional effect of producing a jet with a higher Mach number, $`M_f7`$. We note that the basic features seen in the radiative simulation are quite similar to those seen in the isothermal model ($`\gamma =1.001`$). Thus for brevity we do not present the isothermal results. The only notable difference between the two models is the width of the bow shock. This can be understood purely in terms of the opening angle $`\theta _c`$ of the Mach cone for a supersonic flow, $`\theta _c=\mathrm{sin}^1(1/M_f)`$. We see a wider opening angle for the lower Mach number isothermal flow as expected from the relation for $`\theta _c`$.
The dynamics of both the isothermal and radiative simulations are dominated by the loss of pressure support between the jet- and bow-shocks. In the isothermal model this occurs because $`PT_o`$ where $`T_o`$ is a constant equal to the ambient temperature. In the radiative model the gas behind both the bow- and jet-shocks are driven to temperatures of $`T=(3/16)(\mu /k)V_s^210^5K`$ where ($`s`$) refers to the shock speeds. The cooling time for a jet with $`n_j=100`$ $`\mathrm{cm}^3`$ is $`t_c=.25T/(n_j\mathrm{\Lambda }(T)).2(\tau )=5`$ years. Thus we expect the post shock gas at the head of the jet to cool effectively and for the dynamics to be, essentially isothermal. This is confirmed by consideration of the first frame in Fig. The Propagation of Magneto-Centrifugally Launched Jets: I which shows by $`t=1.38`$ the two shocks have already collapsed on to each other producing a thin shell. The densities and magnetic field strengths in the shell are high with $`1000`$ $`\mathrm{cm}^3`$ $`<n<6000`$ $`\mathrm{cm}^3`$ and $`100\mu G<B<300\mu G`$. Unlike purely hydrodynamic radiative shocks, MHD radiative shocks possess a theoretical limit for the post-shock compression. This occurs because of magnetic pressure exerted by the component of the field perpendicular to the shock normal. Equating ram pressure and magnetic pressure allows a simple form of the maximum post-shock density to be derived (Hollenbach & McKee (1979)),
$$n_m=\frac{(2m_h)^{1/2}(n_j)^{3/2}v_s}{B_{\varphi 0}}.$$
(15)
Note the above use our scaling for the field. For the Keplerian jet $`n_m1000`$ $`\mathrm{cm}^3`$ which is in good agreement with the simple prediction above. The higher densities achieved in the simulation come from the pinch force induced compression on the axis
##### The Shocks
The effect of the radial stratification in both the density and magnetic field (the core/collar structure) are already apparent at the first frame of the simulation just as in the adiabatic models. In the radiative simulation however, the bow-shock quickly assumes a pointed, cusp-like shape. This is due to the higher density in the jet core and the magnetic pinch forces from the strong toroidal field. Compared with the adiabatic simulations described above the head of the jet assumes what might described as a “bullet” shape rather than a nose-cone. In Frank et al. (1998) significant streamlining was observed in the radiative simulations compared with adiabatic ones. This was attributed to the loss of thermal energy and, hence, the increased effectiveness of magnetic stresses. Here we see a similar effect which is enhanced by the increased ram pressure in the jet core relative to the collar.
We also see a mass shedding event in these simulations though it is far weaker than what occurs in the adiabatic models. Frame 2 of Fig. The Propagation of Magneto-Centrifugally Launched Jets: I shows the initialization of the event. As in the adaibatic models a secondary bow shock wave generated by his event leads to pinching of the jet beam and a downstream distortion of the jet shock.
Note the cusp which appears in the jet-shock in the first frame of the simulation. While all features on the axis of an axisymmetric simulation must be taken with some suspicion, a close examination of the simulation data reveal a straight-forward explanation for this structure. The strongest post-shock field values in the jet head occur just downstream of the jet-shock at a radius where the pre-shock field is a maximum. Recall that this occurs just on the outside edge of the jet core $`r.3R_j`$. This is also where $`\beta `$ drops to its lowest value, $`\beta .5`$. Thus magnetic stresses dominate the plasma at this location in the jet head. At radial positions just inward of the point where $`\beta =\beta _{min}`$ we find $`v_r`$ obtaining is maximum inward (negative) value. This radially inward flow is apparent in Fig. The Propagation of Magneto-Centrifugally Launched Jets: I in which we present the velocity vectors at the head of the jet. Thus, at positions immediately downstream of the jet-shock magnetic forces squeeze and compress the jet core. Given the strong cooling, the relation $`P\rho `$ is approximately valid and the axial location of the pinch is a local pressure maximum. The downstream pinch can communicate upstream with the jet-shock face thus producing the bulge or cusp which faces into the oncoming material in the beam. Note that this feature was not seen in the adiabatic models because there the post-shock gas pressure was high enough to inhibit the strong pinch. In fact, $`\beta `$ will always increase across an adiabatic shock. If we write the post-shock compression as $`X=\rho _2/\rho _1`$ with the subscripts $`1`$ and $`2`$ corresponding to pre-and post-shock conditions respectively then (Priest (1986)),
$$\frac{\beta _2}{\beta _1}=\gamma M_f^2\left(1\frac{1}{X^3}\right)\frac{1}{\beta _1}\left(1\frac{1}{X^2}\right).$$
(16)
The equation above shows that for strong adiabatic shocks ($`M_f1`$, $`X4`$) $`\beta _2>\beta _1`$. It is only when the post-shock thermal energy is lost to radiation that the post-shock magnetic forces can dominate.
The magnetic field structure in the jet is shown in the bottom panel of Fig. The Propagation of Magneto-Centrifugally Launched Jets: I. The most prominent feature of the field configuration is the compact size and high field strengths in the jet head. The effect of the pinch wave is clearly apparent upstream of the jet head. The field also shows the effect of the weaker mass shedding event which occurs in this model. Note the isolated loops of $`B_\varphi `$ and strong distortion of $`B_{pol}`$ behind the cusp in the bow-shock. It is also notable that the turbulence and multiple instabilities which are seen in most radiative jet simulations do not occur here. As Gardiner & Frank (1999) have found for their pure poloidal simulations this is one of the principle effects of strong magnetic fields. Thus if YSO jets do contain strong embedded fields then one must consider what their effect on the morphology of the HH objects should be.
### 3.3 Multi-component Rotator
The structural differences between the Keplerian and Multi-component jet originate primarily in the differences in density and velocity cross sections. As we saw in Fig. The Propagation of Magneto-Centrifugally Launched Jets: I the Multi-component jet has three structural elements: a high density core; a low density inner collar; a moderate density outer collar. The ratio of the peak density in the core to minimum density in the inner collar is $`\rho _{max}/\rho _{min}=100`$. This is almost two orders of magnitude higher from what is obtained in the Keplerian jet. In addition, the velocity in the jet peaks in the inner core just at the point where the density drops with $`V_{max}/V_{min}=1.7`$. In the last section we discussed how the (milder) cross sectional variations in the Keplerian jet had important dynamical consequences for its propagation characteristics. Thus we expect the more extreme radial variations in the Multi-component jet will likely effect the dynamics in more extreme ways.
#### 3.3.1 Multi-component Rotator: Adiabatic Jet
##### Propagation and “Peel-off”
Fig. The Propagation of Magneto-Centrifugally Launched Jets: I shows the evolution of the Multi-component jet through four gray-scale maps of $`log`$ density taken at different times in the evolution of the simulation. The most prominent feature in the flow is what we have termed the “peel-off” of the jets’ outer collar. As the jet propagates down the grid, the outer collar develops a strong radial velocity component. As the outer collar expands sideways it is decelerated and develops into a large scale vortex. This is somewhat similar to what was seen in the mass expulsion event seen in the Keplerian case. The high density core of the jet continues its forward propagation driving through the ambient medium at high speed and quickly pulls away from the decelerated collar. By the end of the simulation the “naked” core has propagated far downstream where it encounters the ambient medium in a manner unaffected by the outer collar.
The origin of the peel-off appears to reside in the density stratification of the jet. First we note that since the peel-off occurs early in the simulation we must be suspicious of it as an artifact of the way the simulations are initiated. Given the nature of this study it is difficult to circumvent the need to begin our simulations with a fully formed jet as this is the point of the project. As was noted in section 2 experiments in which the smoothing length between the jet and ambient conditions was varied revealed no change in the propagation characteristics. Even when the initial head of the jet was smoothly joined with ambient medium over a length of many jet radii we found the development of the peel-off was only delayed. The outer layers always developed their transverse motion and the evolution was identical to models with shorter or no smoothing transition. Thus, while the development of this feature may be a transient, it is a highly robust one. Consideration of the dynamics inherent to stratified jets such as these however allows one to infer the mechanism driving the peel-off. We focus on the jet shock. The highest post-shock pressures occur behind the highest velocity regions of the jet. This occurs in the low density inner collar. Since $`\eta (r)<1`$ in this region the jet-shock is relatively strong and is pushed back into the jet deeper than in either the core or outer collar. The oblique geometry of the inner shock generates a strong transverse flow in both positive and negative radial directions. This can be seen in Fig. The Propagation of Magneto-Centrifugally Launched Jets: I which shows the poloidal flow vectors. In addition the variation of the the inner shock produces a finger of high pressure gas which reaches back into the jet in the low density inner collar. A strong radial pressure gradient is established which drives the outer collar away from the core much like splitting wood with an axe. Once the sideways expansion begins, the ram pressure of the ambient medium (in the frame of the jet) continues to divert the flow of the outer collar.
The difference in jet propagation speeds between the Multi-component and Keplerian jets is also dramatic. The velocity of the bow-shock at the end of the simulation is $`V100`$ km $`\mathrm{s}^1`$ which is a $`25\%`$ increase over the propagation speed of the Keplerian jet. This difference can be attributed two effects. First, the Multi-component jet has a higher value of $`\eta `$ in the core relative to the ambient medium ($`\eta 6.8`$ for the Multi-component jet). From equation 12 this translates into a relative propagation velocity difference of $`14\%`$. The excess in propagation speeds above this is likely to be attributable to a second effect - the streamlining of the jet head. Once the outer collar peels away from the core, the jet presents a smaller and more streamlined head to the ambient medium allowing it to propagate at higher speeds. This can be seen by comparing the bow-shock opening angles for the Multi-component and Keplerian jets. Note the streamlining of the head of the naked core also comes via the strong magnetic pinch forces at its outer radial edge and at late times the core also develops the familiar nose-cone morphology.
##### The Instabilities
The development of strong instabilities in the core of the Multi-component jet is another notable characteristic of the simulations. Once the core is exposed we see periodic pinches in the beam. As the instabilities evolve they expand radially and develop a arc-like shape. At later times these arcs become swept backwards by shear in the beam. Further evolution leads to a loss of their sharp edges and individual arcs begin to merge. From detailed consideration of animations of the simulations it appears that the instabilities first appear near the head of the naked core. Only at later times as the peel-off the outer layers continues do they appear immediately downstream of the peel-off region. This is most likely an indication of where the perturbations driving the instabilities occur.
We have performed a stability analysis of the Multi-component magnetic configuration (Lery Lery (1996),Lery & Frank (1999)). A global normal mode stability analysis was performed using the same method as in Appl, Lery & Baty (Appl, Lery & Baty (2000)). In these two papers the stability of magnetized astrophysical jets with respect to modes driven by the electric current density distribution was addressed. The results show that the current driven (CD) instabilities grow rapidly on time scales of order of the Alfvén crossing time in the jet frame and that they are likely to modify the magnetic structure of the jet. Since they are internal modes (see Appl, Lery & Baty (Appl, Lery & Baty (2000))) the CD instabilities should not disrupt the jet. In the present work, we have focused on the pinch mode because of the axisymmetric nature of our calculations. In 3-D it is likely that the kink mode may play a role as well but should not break the integrity of the jet. This is because the jet is super-fast and should see its boundary as a rigid wall. Consequently, the instabilities should be mainly internal, the jet would not be disrupted. The instabilities should certainly be expected to change the jets magnetic configuration drastically .
The analysis shows that the strong pinch, (or sausage), mode is mainly due to large gradients of the density and magnetic field. In Fig. The Propagation of Magneto-Centrifugally Launched Jets: I, we have plotted the dispersion relation for different values of the external pressure surrounding the jet. We have adopted the standard temporal approach where the axial wavenumber is real and the imaginary part of the complex frequency corresponds to growth rate. Wave-numbers are given in units of inverse jet radius and growth rate is normalized to the inverse Alfvén time. It has been found that the location of the peak mainly depends on the magnetic distribution in the jet. The short-$`k`$ cut-off is due to the finite size of the jet radius that has an external boundary that behaves as a rigid wall for Mach numbers larger than unity (see Appl, Lery & Baty (Appl, Lery & Baty (2000))). We find that pure magnetic instabilities driven by electric current develop on rapid Alfvén time-scales. Also, Fig. The Propagation of Magneto-Centrifugally Launched Jets: I clearly shows that when the external pressure increases the jet becomes more unstable. This is precisely what we observe when the core becomes naked downstream of the peel-off region. It also explains why the instabilities do not develop as rapidly for the Keplerian case where the density and pressure gradients are less important.
From the simulations it appears that at later times waves driven off the peel-off region seed the instabilities while at earlier times the seeds occur via shocks at the jet head. As can be seen in Fig. The Propagation of Magneto-Centrifugally Launched Jets: I the pinching instabilities on the axis have a wavelength of approximately $`\lambda .5R_j`$. Note however the presence of a second characteristic wavelength which runs along the surface of the core. This feature, which appears as an envelope encompassing the shorter wavelength modes, has $`\lambda 3R_j`$.
The stability analysis of Multi-component jets shows that the unstable modes that should grow the most rapidly have a wavelength of about 3 jet radii for the collar and half the jet radius for the core. These results have been reported by Lery & Frank (1999). They are in good agreement with the simulations, and can be also compared to observations. For example, the jet of HH34 presents a mean knot separation of $`3.4r_{jet}`$ as given by Burke et al. (burke (1988)). Thus, the present results suggest that these instabilities could be at the origin of the knotty structure of a large number of jets as seen, for example, in HL Tau, HH1, HH30 and HH34 (Ray et al. rayetal (1996)).
Finally consider the magnetic structure in the jet which is shown in Fig. The Propagation of Magneto-Centrifugally Launched Jets: I. The field structure is quite complicated as might be expected. Note the form of the bow-shock in the $`B_{pol}`$ component as well as its the relative absence in the peeled off outer collar which is dominated by toroidal fields. Examination of $`\beta `$ in these regions shows that the gas remains hydrodynamically dominated with $`\beta 1`$ in spite of amplification from both the shocks and radial stretching. Within the core the pinch modes are clearly seen in the toroidal component with specific islands in the beam corresponding to regions of strong pinch. Numerous islands of $`B_{pol}`$ are created by the instabilities indicating the presence of reconnection.
#### 3.3.2 Multi-component Rotator: Isothermal/Radiative Jet
In Fig. The Propagation of Magneto-Centrifugally Launched Jets: I we present gray-scale maps of the density evolution of a radiative jet driven by a Multi-component rotator. The radial distributions were scaled down for all variables except $`v_z`$ in the jet such that $`M_f9`$. The basic features of the simulation are similar to the isothermal model however there are some differences. To address these we also present in Fig. The Propagation of Magneto-Centrifugally Launched Jets: I a single frame from the isothermal simulation. Note first that, once again, the width of the bow-shock is reduced in both the radiative and isothermal case relative to adiabatic model. This can be attributed both to cooling and the increase in fast mode Mach number.
The most important point to notice in this simulation is that the outer collar still peels away from the jet core which then propagates ahead of the rest of the flow. This occurs even though the cooling is strong. As in the adiabatic model, the initiation of transverse flow in the outer layers occurs due to the shape of the jet shock. As the first frame of Fig. The Propagation of Magneto-Centrifugally Launched Jets: I demonstrates, with cooling included both the bow-shock and jet-shock effectively “drape” around the head of the jet. This feature occurs due to the loss of pressure support behind the shocks. As in the adiabatic model the peel-off appears to be primarily driven by the redirection of the flow behind the oblique jet-shock. Note that, in spite of cooling, the jet-shock in the low density collar (where $`\eta <1`$) must sink back into the body of the jet. The bow-shock follows suit and the result is a highly oblique section of the shock in the inner collar. When undisturbed beam material impinges on this shock it is either directed towards the axis forming a strong pinch in the core or it is shunted radially outward forcing the outer collar to peel away. Thus in both the adiabatic, isothermal and radiative cases the non-uniformity in the jet-shock drives a flow pattern which enhances the “jet within a jet” nature of the flow and the core always ends up propagating away from the rest of the beam.
Unlike the adiabatic model where the peel-off region had a low zed velocity, both the radiative and isothermal models show the point at which the core and collar separate moves with a speed that is a large fraction of the beam speed. The origin of this effect appears to be the lower pressures behind the jet-shock which causes less deceleration. The propagation of the separation point may also be due to the reduced width of the bow-shock and a smaller cocoon (both expected in non-adiabatic models).
The magnetic field structure shown in Fig. The Propagation of Magneto-Centrifugally Launched Jets: I is similar to what is seen in the adiabatic case. Note, however, the strong pinch which occurs at the point where the peel-off occurs. The loss of gas pressure support will also decrease $`\beta `$ implying that the toroidal field can now exert a stronger influence.
The principle difference between the radiative and isothermal models occurs in the core. First note that it is difficult to see the instabilities in the radiative model. A detailed inspection of the simulation frames shows they are present but they appear to diffuse more rapidly than in the isothermal case. The isothermal simulations do show the same form of the modes occurring as in the adiabatic models and with similar length scales. The difference between the isothermal and radiative solutions is likely to due the greater thermalization which occurs in the higher Mach number flow.
## 4 Comparisons with Observations
In recent observations of molecular outflows (Dutrey et al. (1997), Gueth & Guilloteau (1997), Gueth et al. (1998)) show small linear structures just ahead of the familiar bow-shaped shocks. Three examples of such features are presented in Fig. The Propagation of Magneto-Centrifugally Launched Jets: I. These structures point almost exactly away from the position of the protostellar condensation. These precursors of the bow-shock show a roughly conical shape. As such they could trace an underlying jet which is propagating beyond the bow-shock. The present simulations are suggestive offering an explanation for the observed structures. The fast “core”-jet propagates ahead of the collar and the surrounding molecular outflow.
The outer “collar” may be solely responsible for the larger bow-shock structure or it may itself be embedded in a larger wide angle wind. Thus molecular observations of conical precursors to the bow-shocks may be a signature of density and magnetic stratification discussed in this study. Therefore, the global evolution that we obtain for our jets, e.g., a core-collar structure could lead to common behavior for several YSO jets, and also may help in understanding the relation between jets and molecular outflows.
## 5 Discussion and Conclusions
We have carried out a series of simulations intended to address the issue of MHD jet propagation. Whereas previous studies have used ad-hoc initial conditions we inject flows into our computational grid derived from models of collimated jets driven by magneto-centrifugal launching. This strategy allows us to compare the propagation characteristics of jets driven by different types of outflows. In particular we have studied the propagation of jets driven by: (1) a purely Keplerian rotator (a disk) exterior to a solid body rotator (a star); (2) a Keplerian rotator with a sub-Keplerian boundary layer both of which are exterior to a solid body rotator. The former model we refer to as a Keplerian jet, the latter is called a Multi-component model. Our simulations follow the jets out to observable scales. In the Keplerian jet simulations the jet radius is $`R_j=1.5x10^{15}`$ cm making the grid extend out to $`3000AU`$. For the multi-component jet $`R_j`$ is almost a factor of ten larger and the the grid extends out $`.1pc`$. The width of the multi-component jet is interesting in that it yields a model with a very narrow, dense core (a jet) surrounded by a wider lower density outflow.
Both models were calculated under the assumption that the jets are launched under isothermal conditions. We have carried out simulations of the propagation of both Keplerian and Multi-component jets under adiabatic, isothermal and radiative conditions in order to determine the behavior of the resulting flows with, and without, radiative losses. We note again that our adiabatic and isothermal simulations have low magneto-sonic Mach numbers $`M_{ms}=24`$. While these values are small compared with the values used in previous numerical studies of MHD jets $`(M_{ms}>10`$, Stone & Hardee (1999)) they is quite similar to what has been obtained in other studies MHD collimation of jets ($`M_{ms}3`$,Camenzind (1997)).
Our simulations show significant differences in the propagation characteristics for the two types of rotators. In addition, features are seen in both classes of jet which have not been seen in previous models of either pure hydrodynamic or MHD jet propagation. In all cases it appears that the most important aspect of the flow behavior seen in the simulations can be traced back to the annular stratification of the jets. In particular, the radial distributions of density, velocity, and toroidal magnetic field appear to be the principle causes of the new behavior seen in the simulations. Both Keplerian and Multi-component jets exhibit a core/collar structure such that a high density core region exists near the axis surrounded by one or more lower density annuli (collar) extending out to the jet boundary. The strongest toroidal fields exist at the boundary between the core and collar.
Since the momentum in the core is higher than that in the collar the propagation characteristics of the jets are dominated by the core pulling ahead of the collar. The strong field surrounding the core ensures that the two regions remain fairly distinct in terms of their dynamics. As the jets propagate we see the core acting as a jet within a jet. In the Keplerian case the relatively low density contrast between core and collar keeps the two propagating at relatively similar velocities. The stratification of the magnetic fields produces strong dynamical differences between core and collar. All plasma flowing into the cocoon comes from the lower field strength regions of the collar. In the Multi-component case there exists an extremely low density inner collar (which also has higher velocity than the surrounding regions) and this leads to a complete separation of core and collar. The “peel-off” of the collar in the Multi-component models is quite dramatic and occurs in both the adiabatic and isothermal simulations.
Our results have bearing on a number of issues. The simplest conclusion that can be drawn is that the structure imposed on a YSO jet by the launching and collimation process can lead to fairly complex propagation characteristics. Thus our models build on and extend the previous works which utilized only “top-hat” jets as initial conditions. Our results also indicate that jets launched from different classes of rotators may have different propagation characteristics. It is likely that in real jet systems the dynamics is too complex to make an isomorphic identification of a given class of rotators with a set of observed jet morphologies. There is however the possibility that as these kinds of studies mature one might be able to distinguish between different classes of MHD launching models via consideration of the way the jets from these models would appear on the sky.
Finally we note that given the large parameter space of initial conditions for both the Given Geometry Model and for the jet propagation simulations, the work described here which focuses only on two instances must be seen as preliminary. It does however point to the fact that the jets produced by magnetized rotators are likely to be more complex in their structure and, furthermore, that this complexity will be reflected in the observed jet morphologies. In future studies we will attempt to build a larger catalog of jet propagation characteristics through a more thorough exploration of parameter space of the Given Geometry Model.
We wish to thank Guy Delemarter, Jack Thomas, Colin Norman and Lee Hartmann for their input and discussions leading to this paper. This work was supported by NSF Grant AST-0978765. DR was supported in part by KOSEF through grant 981-0203-0011-2.
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# 1 Introduction
## 1 Introduction
The non-perturbative computation of the running QCD coupling constant $`\alpha _s(p)`$ follows a two-sided goal: the large energy matching to perturbative asymptotic QCD formula turns out to be a most direct method to predict $`\mathrm{\Lambda }_{\overline{\mathrm{MS}}}`$ from QCD first principles . On the other hand, the moderate or low energy behavior of $`\alpha _s(p)`$ is extremely instructive about non-perturbative properties of QCD. In this paper we restrict ourselves to high and intermediate energies and consider power corrections ($`1/p^2`$) to the leading asymptotic behavior. As we shall see it turns out that there is no sharp separation between the asymptotic domain and the intermediate one. The power correction, beyond the lessons it contains by itself, greatly improves our asymptotic study and is never negligible up to $`10`$ GeV ! This surprising fact could only be revealed thanks to the high accuracy achieved in the present work.
The asymptotic approach has been recently followed in Ref. . It is worth remarking that this matching procedure has also been developed for the two-point Green function to the same goal , both matching programs leading to a consistent estimate of $`\mathrm{\Lambda }_{\overline{\mathrm{MS}}}`$. The almost two-sigma discrepancy between the last estimate and the one obtained from Schrödinger functional methods seems to imply that some source of systematic uncertainty remains uncontrolled. Furthermore the careful study of the asymptotic behavior carried out for the gluon propagator in Refs. stresses a danger: it is sometimes difficult to distinguish between the real asymptotic scaling and a behavior mimicking asymptoticity but with a certain effective “re-scaled” $`\mathrm{\Lambda }`$ parameter which differs significantly from the real asymptotic one. Consequently, and in spite of the agreement between both above-mentioned estimates of $`\mathrm{\Lambda }_{\overline{\mathrm{MS}}}`$ (matching two or three-point Green functions to perturbative formulae), we must inquire whether both are not biased by some sizable non-perturbative effect.
The operator product expansion (OPE) gives a standard procedure to parametrize non-perturbative QCD effects in terms of power corrections to perturbative results. In this framework, the powers involved in the expansion are expected to be uniquely fixed by the symmetries and the dimensions of the operators appearing in the product expansion. It should be noted that, due to the asymptotic nature of QCD perturbative expansions, power corrections are reshuffled between operators and coefficient functions in the OPE .
Since we work in Landau gauge, the gauge dependent local operator $`A_\mu A^\mu `$ is allowed in the OPE implying a dominant $`1/p^2`$ power correction. Indeed sizable $`1/p^2`$ corrections are present as we shall show at length in the next section.
In a less straightforward way our result is related to another hot issue (this is the spirit of the preliminary study in ): the possible presence of $`1/p^2`$ terms in gauge invariant quantities such as Wilson loops . Since no gluon local gauge invariant operator of dimension less than 4 exists it is expected from OPE that the dominant power correction should be $`1/p^4`$. This picture has however recently been challenged . It was pointed out that power corrections which are not a priori expected from the OPE may in fact appear in the expansion of physical observables. Such terms may arise from (UV-subleading) power corrections to $`\alpha _s(p)`$, corresponding to non-analytical contributions to the $`\beta `$-function. It is worth stressing the following. One knows that the perturbative analysis does not encode all the information on the coupling. Among all that is missed, a peculiar contribution to the coupling could result in a peculiar correction to physical observables. As a matter of fact, some evidence for an unexpected $`\frac{\mathrm{\Lambda }^2}{p^2}`$ contribution to the gluon condensate was obtained through lattice calculations in Ref. (see also for an early evidence of such a contribution, although the perturbative series involved was not managed up to high orders).
Let us insist: there is no direct relation between the $`1/p^2`$ corrections found in the present work to $`\alpha _s(p)`$ in a gauge dependent scheme and the power corrections advocated in the preceding paragraph to justify $`1/p^2`$ corrections to gauge independent quantities. Nobody knows how to relate a gauge dependent scheme to a gauge independent one beyond the perturbative regime. Still, it may be that these two phenomena are not completely unrelated and the large $`1/p^2`$ corrections found in the present work might trigger some further thoughts along this line.
The paper is organized as follows: in Section 2 we explain the meaning of the lattice data, our strategy for the analysis and report the results. In Section 3 we briefly review some theoretical arguments in support of power corrections to $`\alpha _s(p)`$, illustrating the special role of the $`\frac{\mathrm{\Lambda }^2}{p^2}`$ term. Finally, in Section 4 we draw our conclusions.
## 2 Lattice calculation of $`\alpha _s`$ from Green functions and power corrections
### 2.1 $`\alpha _s`$ on the lattice
Several methods for computing $`\alpha _s(p)`$ non-perturbatively on the lattice have been proposed in recent years . In most cases, the goal of such studies is to obtain an accurate prediction for $`\alpha _s(M_Z)`$, i.e. the running coupling at the $`Z`$ peak, which is a fundamental parameter in the standard model. For this reason, lattice parameters are usually tuned so as to allow the computation of $`\alpha _s(p)`$ at momentum scales of at least a few GeVs, where the two-loop asymptotic behavior is expected to dominate and power contributions are suppressed.
As we shall see, however, non-perturbative power corrections cannot be neglected at energy scales as large as 10 GeV which is a sufficient reason to consider them in the fit. As a bonus the knowledge of these power corrections provides us with a physically significant quantity as argued in the introduction.
For this purpose, the best method is one where one can measure $`\alpha _s(p)`$ in a wide range of momenta from a single Monte Carlo data set. Indeed, a narrow energy window does not allow to disentangle in a clear cut manner the power corrections from unknown higher perturbative orders, and these corrections can be mimicked by an effective $`\mathrm{\Lambda }_{\mathrm{QCD}}`$ different from the asymptotic one.
One method which fulfills the above criterion is the determination of the coupling from the renormalized lattice three-gluon vertex function . This is achieved by evaluating two- and three-point off-shell Green’s functions of the gluon field in the Landau gauge, and imposing non-perturbative renormalization conditions on them, for different values of the external momenta. By varying the renormalization scale $`p`$, one can determine $`\alpha _s(p)`$ for different momenta from a single simulation. Obviously the renormalization scale must be chosen in a range of lattice momenta such that both finite volume effects and discretization errors are under control. Both these constraints impose too strict limits on the energy range if only one lattice run is used. This is why the procedure used in Ref. combining several lattice simulations at different $`\beta `$’s is a necessity to get a larger range of lattice momenta. The use of different volumes will also help to control finite volume artifacts.
The definition of the coupling corresponds to a momentum-subtraction renormalization scheme in continuum QCD . It should be noted that in this scheme the coupling is a gauge-dependent quantity. As we already mentioned, one consequence of this fact is that $`1/p^2`$ corrections should be expected, based on OPE considerations. For full details of the method and the lattice calibration we refer the reader to Ref. .
### 2.2 Models for power corrections and construction of the data set
In the present work we shall not address the general problem of defining an optimal analytic form for the coupling at all scales to which we could fit our data. For the purpose of our investigation, we shall compare the non-perturbative data for $`\alpha _s`$ with simple models obtained by adding a power correction term of the form $`1/p^2`$ to the perturbative formula at a given order. This amounts to a quite raw separation between a perturbative versus a non–perturbative contribution, the major problem of course being the possible interplay between a description in terms of (non-perturbative) power corrections and our ignorance about higher orders of perturbation theory. As crude as it may be, our recipe allows for addressing this problem, as we shall see.
In order to identify a momentum interval where our ansatz could fit the data, one should keep in mind that the momentum range should start well above the location of the perturbative Landau pole, but it should nonetheless include low scales where power corrections may be large up to large enough momentum scales in order to be confident that the asymptotic regime (i.e. perturbative running) has become dominant. Our choice will be<sup>3</sup><sup>3</sup>3The range’s upper limit is determined by the condition $`ap\pi /2`$, where $`a`$ is the lattice spacing, applied to our $`24^4`$ lattice at $`\beta =6.8`$. 2.0 GeV $`p`$ 9.6 GeV. It will turn out that in the whole range both the perturbative three-loop contribution and the non perturbative $`1/p^2`$ one contribute by a sizable amount, although obviously the former becomes dominant at larger scales.
A data set which fulfills the above-mentioned requirements can be constructed by aggregating data computed at different $`\beta `$ values ($`\beta =`$6.0, 6.2, 6.4, 6.8) on a $`24^4`$ lattice. The fact that such a data set can be assembled is a very good a posteriori check that the expected scaling in the lattice cutoff $`a`$ takes place.
On the other hand, the physical volumes for these simulations are very different from each other. The appropriate matching of the data proves that finite-size effects remain controlled. The pattern for these volume effects is clear from fig. 1(a), where the whole set of data is plotted, including the points too much affected by finite volume artifacts and eventually rejected in our fits. As can be seen, the effects are negligible for large enough $`Lp`$, where L is the physical lattice length. Evidences for the same trend arise from the comparison of data obtained on two different lattices ($`16^4`$, $`24^4`$) at $`\beta =6.8`$, shown in fig. 1(b). In practice we will take as infrared cut-off at each $`\beta `$ the value $`p_{\mathrm{min}}`$ such that, including points below this value, the $`\chi _{d.o.f}^2`$ increases dramatically. This criterion leads to: $`p_{min}(6.0)=2.0`$ GeV, $`p_{min}(6.2)=2.5`$ GeV, $`p_{min}(6.4)=4.0`$ GeV and $`p_{min}(6.8)=6.0`$ GeV. It corresponds to $`Lp24`$ <sup>4</sup><sup>4</sup>4Almost all the data from our $`16^4`$ lattice at $`\beta =6.8`$ turn out to be excluded by such a requirement, $`Lp24`$ leading for example to an infra-red cut-off equal to 9. GeV. This is why we will not use at all the volume $`16^4`$.. Reversely, when we vary the infrared cut-offs above the values just mentioned, the $`\chi _{d.o.f}^2`$ stays stable, in the range 1.4 to 1.6. Even more striking, the resulting value for $`\mathrm{\Lambda }_{\overline{\mathrm{MS}}}`$ is very stable, varying no more than 2 MeV. Consequently the data set obtained with these cut-offs should be considered as IR safe and will be used in the following fits.
At each $`\beta `$ value one should of course also worry of data in the range of the larger values of momentum, which are affected by lattice artifacts of $`O(a^2p^2)`$ (UV discretization effects). This was taken care of by the sinus-improvement program which has already been described in , which basically amounts to the substitution of the lattice momenta $`\frac{2\pi n_\mu }{L}`$ with their $`O(a^2p^2)`$ analogues $`\frac{2}{a}\mathrm{sin}(\frac{ap_\mu }{2})`$. As already noticed in , a good rationale for this in our gauge-fixed situation is that the gauge fixing algorithm leads to the relation $`\frac{2}{a}\mathrm{sin}(\frac{ap_\mu }{2})A_\mu (p)=0`$, while $`p_\mu A_\mu `$ does not vanish<sup>5</sup><sup>5</sup>5One should keep in mind that we are imposing Landau gauge.. It should be noticed that without this prescription the quality of almost any fit degenerates. The relevant configurations were generated on a QH1 Quadrics system based in Orsay (see ).
### 2.3 Fitting data to our ansatz
We now proceed to fitting the data set to our ansatz, which, we recall, is the addition of a term proportional to $`\frac{1}{p^2}`$ to a given order of perturbation theory:
$`\alpha _s(p^2)=\alpha _s^{Pert}(p^2)+{\displaystyle \frac{c}{p^2}}.`$ (1)
Working at three loop level, the perturbative expression for the running coupling constant, $`\alpha _S^{Pert}(p^2)`$ requires to inverse either the unexpanded formula
$$\stackrel{~}{\mathrm{\Lambda }}=\stackrel{~}{\mathrm{\Lambda }}^{(c)}(\stackrel{~}{\alpha })\left(1+\frac{\beta _1\stackrel{~}{\alpha }}{2\pi \beta _0}+\frac{\stackrel{~}{\beta }_2\stackrel{~}{\alpha }^2}{32\pi ^2\beta _0}\right)^{\frac{\beta _1}{2\beta _0^2}}$$
$`\times \mathrm{exp}\left\{{\displaystyle \frac{\beta _0\stackrel{~}{\beta }_24\beta _1^2}{2\beta _0^2\sqrt{\mathrm{\Delta }}}}\left[\mathrm{arctan}\left({\displaystyle \frac{\sqrt{\mathrm{\Delta }}}{2\beta _1+\stackrel{~}{\beta }_2\stackrel{~}{\alpha }/4\pi }}\right)\mathrm{arctan}\left({\displaystyle \frac{\sqrt{\mathrm{\Delta }}}{2\beta _1}}\right)\right]\right\}`$ (2)
or the expanded one
$`\stackrel{~}{\mathrm{\Lambda }}=\stackrel{~}{\mathrm{\Lambda }}^{(c)}(\stackrel{~}{\alpha })\left(1+{\displaystyle \frac{8\beta _1^2\beta _0\stackrel{~}{\beta }_2}{16\pi ^2\beta _0^3}}\stackrel{~}{\alpha }\right)`$ (3)
where $`\stackrel{~}{\mathrm{\Lambda }}^{(c)}`$ denotes the conventional two loops formula:
$`\stackrel{~}{\mathrm{\Lambda }}^{(c)}p\mathrm{exp}\left({\displaystyle \frac{2\pi }{\beta _0\stackrel{~}{\alpha }(p^2)}}\right)\times \left({\displaystyle \frac{\beta _0\stackrel{~}{\alpha }(p^2)}{4\pi }}\right)^{\frac{\beta _1}{\beta _0^2}};`$ (4)
and $`\mathrm{\Delta }2\beta _0\stackrel{~}{\beta }_24\beta _1^2>0`$ (for our $`\stackrel{~}{\mathrm{MOM}}`$ scheme). In the previous formula, the use of $`\stackrel{~}{\mathrm{\Lambda }}`$, $`\stackrel{~}{\alpha }`$ and $`\stackrel{~}{\beta }`$’s stands for the $`\mathrm{\Lambda }`$ parameter, the running coupling constant and beta function coefficients in the particular $`\stackrel{~}{\mathrm{MOM}}`$ renormalization scheme. From now on we will systematically convert $`\stackrel{~}{\mathrm{\Lambda }}`$ into $`\mathrm{\Lambda }_{\overline{\mathrm{MS}}}`$ using
$`\mathrm{\Lambda }_{\overline{\mathrm{MS}}}=\mathrm{exp}(70/66)\stackrel{~}{\mathrm{\Lambda }}0.346\stackrel{~}{\mathrm{\Lambda }}`$ (5)
In Eqs. (2-3) the $`p^2`$ dependence of $`\stackrel{~}{\alpha }`$ has been omitted for simplicity.
No analytical expression can exactly inverse neither unexpanded eq. (2) nor expanded eq. (3). The following formula gives an approximated solution to the inversion of eq. (3):
$`\stackrel{~}{\alpha }(p^2)`$ $`=`$ $`{\displaystyle \frac{4\pi }{\beta _0t}}{\displaystyle \frac{8\pi \beta _1}{\beta _0}}{\displaystyle \frac{\mathrm{log}(t)}{(\beta _0t)^2}}`$ (6)
$`+{\displaystyle \frac{1}{(\beta _0t)^3}}\left({\displaystyle \frac{2\pi \stackrel{~}{\beta }_2}{\beta _0}}+{\displaystyle \frac{16\pi \beta _1^2}{\beta _0^2}}(\mathrm{log}^2(t)\mathrm{log}(t)1)\right)`$
where $`t=\mathrm{log}(p^2/\stackrel{~}{\mathrm{\Lambda }}^2)`$. On the other hand, an exact numerical inversion of Eq. (2) can be easily obtained and used to fit our data. Both ansätze, eq. (6) and the numerical inversion of eq. (2), should only differ by perturbative contributions higher than three loops. Thus, we will fit with the two ansätze, the difference between these two predictions being considered as an estimate of the systematic uncertainty coming from the neglected perturbative orders. To make more direct the comparison with the Schrödinger functional estimate of $`\mathrm{\Lambda }_{\overline{\mathrm{MS}}}`$ , the central value for our prediction of $`\mathrm{\Lambda }_{\overline{\mathrm{MS}}}`$ should be taken from the fit with the exact inversion of Eq. (2)<sup>6</sup><sup>6</sup>6The ALPHA collaboration estimate: $`\mathrm{\Lambda }_{\overline{\mathrm{MS}}}=238(19)`$ MeV comes from evaluating eq. (2) for a value of $`\alpha _S`$ obtained from the lattice at very high momenta. . This yields:
$`\mathrm{\Lambda }_{\overline{\mathrm{MS}}}=237\pm 3\mathrm{MeV},c=0.63\pm 0.03\mathrm{GeV}^2,\chi ^2/\mathrm{d}.\mathrm{o}.\mathrm{f}.=1.6;`$ (7)
in perfect agreement with the determination which uses a totally different technique. Using eq. (6) we obtain $`\mathrm{\Lambda }_{\overline{\mathrm{MS}}}=227(5)\mathrm{MeV}`$ and $`c=0.50(6)\mathrm{GeV}^2`$. Comparing the latter results with eq. (7) provides an estimate of the higher loop uncertainty of about 10 MeV for $`\mathrm{\Lambda }_{\overline{\mathrm{MS}}}`$ and $`0.1\mathrm{GeV}^2`$ for $`c`$. The size of the power correction will be discussed in sections 3.1 and 3.2.
Fig. 3 illustrates the effect of the power correction in a striking manner. The upper set of points shows $`\mathrm{\Lambda }`$, converted to $`\mathrm{\Lambda }_{\overline{\mathrm{MS}}}`$, computed through eq. (2) from $`\alpha _s`$ provided by the lattice at every value $`p`$. Scaling would imply a constancy of $`\mathrm{\Lambda }_{\overline{\mathrm{MS}}}`$ which is far from true. The lower set of points corresponds to the same formula applied to $`\alpha _s(p^2)^{\mathrm{lattice}}0.63/p^2`$, i.e. to what should fit the perturbative formula. Now the constancy of $`\mathrm{\Lambda }_{\overline{\mathrm{MS}}}`$ is very good within the statistical errors. It is now clear why the upper points show a trend to decrease: as the energy scale increases, the effect of the power corrections must decrease, and the upper points converge slowly towards the lower ones. The surprise is that above 9.0 GeV the merging of the two sets of points has not yet taken place contrarily to the general expectation that power corrections are negligible at such a scale. We will elaborate on this in the next section.
We note that using eq. (6) one can draw the same conclusion: imposing $`c=0`$ (i.e. fitting to a pure three loop formula) there is no good fit on the whole range of momenta and the best that one can obtain on a restricted interval yields a value for $`\mathrm{\Lambda }_{\overline{\mathrm{MS}}}`$ higher than expected (e.g. with respect to ) and a definitly worse $`\chi _{d.o.f}^2`$.
We now proceed to address the question of the stability of our previous result with respect to the inclusion of the fourth loop. To fit the value of $`\mathrm{\Lambda }`$, one would need $`\beta _3`$, the four loop coefficient in the perturbative expansion of the $`\beta `$-function. This coefficient is unknown but one can pin down a reasonable range of variation. Let us consider the ratio $`b_3/b_2`$ with $`b_3=\beta _3/(4\pi )^3`$ and $`b_2=\beta _2/(2\pi )^2`$. This ratio is larger than one in the only scheme $`(\overline{\mathrm{MS}})`$ for which it is known. On the other hand we expect that, for a reasonably small value of the coupling, the perturbative expansion of the $`\beta `$-function is still in the regime in which it seems to converge at four loops. As a working hypothesis, we then suppose that the contribution to the $`\beta `$-function coming from $`\beta _3`$ (that is, the one proportional to $`\alpha ^5`$) is not too much larger than the one coming from $`\beta _2`$ at a typical value of $`\alpha 0.4`$, implying $`b_3/b_2<2.5`$. Actually we conservatively vary $`b_3`$ from 0 up to $`5b_2`$ in the following exploration of a large range of values for $`b_3`$.
First of all we try a pure four loop fit (that is without any power correction). We observe that there is no good fit on the whole range of momenta. If one tries to add again a term proportional to $`\frac{1}{p^2}`$ to the four loop perturbative expression, the following should be noted. Good fits are recovered either on the whole range of momenta for $`b_3/b_2<\mathrm{\hspace{0.33em}1}`$ or by discarding momenta below $`3`$ GeV as $`b_3/b_2<\mathrm{\hspace{0.33em}4}`$. As far as the value of $`\mathrm{\Lambda }`$ is concerned, it is really stable when the fits are of good quality<sup>7</sup><sup>7</sup>7Even for $`b_3/b_2=5`$, one could obtain a good quality fit ($`\chi ^2=1.8`$) with $`\mathrm{\Lambda }_{\overline{\mathrm{MS}}}=238`$, if momenta below $`4`$ GeV are now discarded.. As for the coefficient of the non-perturbative term $`\frac{1}{p^2}`$, it is less stable. Results are summarized in Tab. 1. One can see how the “new player on the ground”, $`\beta _3`$, is strongly correlated to the coefficient of the power term. This coefficient is anyhow fully consistent with the value determined at three loop level, as long as $`b_3/b_2<\mathrm{\hspace{0.33em}2}`$ i.e. when the asymptotic behavior is not too dramatically perturbed by the four-loop contribution.
Again, much the same holds when fitting to a formula for $`\alpha _s`$ as a function of momentum (like eq. (6)) at the four loop level. In order to trust a four loop more than power corrections one would need both to discard lower momenta and to accept excedingly large values for $`b_3`$.
Of course, the perturbative knowledge of $`\beta _3`$ coefficient is unavoidable to get total confidence on our results versus higher loop orders inclusion. One should take this last analysis only as a preliminary check. Still it is another indication that things are working pretty well (that is, consistently) with respect to our theoretical prejudice. In any case, the value obtained for $`\mathrm{\Lambda }_{\overline{\mathrm{MS}}}`$ is almost insensitive to $`\beta _3`$.
We collect all the hints from these many counterproofs. To estimate the systematic errors we use two methods: first we compare the exact inversion of eq. (2) and the use of eq. (6) on our results, second we use the results on the whole energy range in table 1 with a reasonable $`\chi ^2`$. We take the three loop result with formula (2) as our central value
$`\mathrm{\Lambda }_{\overline{\mathrm{MS}}}=237\pm 3_{10}^{+0}\mathrm{MeV},c=0.63\pm 0.03_{0.13}^{+0.0}\mathrm{GeV}^2`$ (8)
The present analysis has improved over the one presented in by using a larger data set which provides a wider momentum window and taking into account a power correction. If one tried to repeat the fits only in the range of the lower momenta as we did in , there would be no clear cut indication for power corrections and the effective value for $`\mathrm{\Lambda }`$ would turn out to be higher. What the upper momenta data really do is to single out the asymptotic value for $`\mathrm{\Lambda }`$, while deviation from asymptotia in the lower momenta data asks for power correction rather than higher loops.
## 3 Discussion on the $`\frac{1}{p^2}`$ corrections to $`\alpha _s(p)`$
Power corrections to $`\alpha _s(p)`$ can be shown to arise naturally in many physical schemes . The occurrence of such corrections cannot be excluded a priori in any renormalization scheme. Even more so, as previously stated, in a gauge dependent renormalization scheme as the $`\stackrel{~}{\mathrm{MOM}}`$ discussed here. Clearly, the non-perturbative nature of such effects makes it very hard to assess their dependence on the renormalization scheme, which is only very weakly constrained by the general properties of the theory.
As discussed in the following, several arguments have been put forward in the past to suggest that a most likely candidate for a leading power correction to $`\alpha _s(p)`$ would be the same term of order $`\mathrm{\Lambda }^2/p^2`$ we found. Furthermore, it is worth emphasizing that this does not contradict the OPE expectation for a gauge dependent quantity.
### 3.1 Static quark potential and confinement
Consider the interaction of two heavy quarks in the static limit (for a more detailed discussion see ). In the one-gluon-exchange approximation, the static coulombic potential $`V(r)`$ can be written as
$$V(r)\alpha _sd^3k\frac{\mathrm{exp}^{i\stackrel{}{k}\stackrel{}{r}}}{|\stackrel{}{k}|^2}.$$
(9)
Using standard arguments of renormalon analysis, one may consider a generalization of (9) obtained by replacing $`\alpha _s`$ with a running coupling:
$$V(r)d^3k\alpha _s(|\stackrel{}{k}|^2)\frac{\mathrm{exp}^{i\stackrel{}{k}\stackrel{}{r}}}{|\stackrel{}{k}|^2}.$$
(10)
The presence in $`\alpha _s(k^2)`$ of a power correction term of the form $`c/k^2`$ would generate a linear confining piece $`Kr`$ in the potential $`V(r)`$. It would of course be totally unjustified to take seriously our power correction in order to derive the linear potential. If we nevertheless perform this exercise to have a feeling of the scales, we get $`K=2/3c0.4`$ GeV<sup>2</sup>, while the usual string tension is around $`0.2`$ GeV<sup>2</sup>. Note that a standard renormalon analysis of (10) (see for the details) reveals contributions to the potential containing various powers of $`r`$, but a linear contribution is missing. This is a typical result of renormalon analysis: renormalons can miss important pieces of non-perturbative information.
### 3.2 An estimate from another lattice method
The lattice community has been made aware for some time of the possibility of non-perturbative contributions to the running coupling; for a clear discussion see . Consider the “force” definition of the running coupling:
$$\alpha _{q\overline{q}}(Q)=\frac{3}{4}r^2\frac{dV(r)}{dr}(Q=\frac{1}{r}),$$
(11)
where again $`V(r)`$ represents the static interquark potential. By keeping into account the string tension contribution to $`V(r)`$, which can be measured in lattice simulations, one obtains a $`1/Q^2`$ contribution, whose order of magnitude is given by the string tension itself. Ironically, this term has been mainly considered as a sort of ambiguity, resulting in an indetermination in the value of $`\alpha (Q)`$ at a given scale. From a different point of view, such a term could be interpreted as a clue for the existence of a $`\frac{\mathrm{\Lambda }^2}{p^2}`$ contribution, and it also provides an estimate for the expected order of magnitude of it, at least in one (physically sound) scheme. The same naïve exercise than in the preceding subsection leads from eq. (11) to $`K=4/3c0.8`$ $`\mathrm{GeV}^2`$.
In order to make a deeper contact with what we are actually studying, one should of course keep in mind that all the preceding arguments would be physically sounder in the Coulomb gauge in which the notion of force has a clearer meaning. In Landau gauge such a picture is far from clear. Furthermore, it is known that the interquark potential becomes close to linear at distances larger than 0.5 fm, which corresponds to a momentum smaller than 0.4 GeV, while our $`1/p^2`$ fit only apply above 2.0 GeV i.e. at distances smaller than 0.1 fm. Still, a couple of considerations are in order at this point. We note that the rough order of magnitude of the power term we found is the same as inferred from the simple arguments from the static potential (some $`10^1\mathrm{GeV}^2`$), even if we are in a different (UV) regime. While this could turn out to be accidental, it is nevertheless intriguing to think of some sort of relation. In any case, one should nevertheless note that the power correction that we have obtained is rather large with respect to the common wisdom of non-perturbative effects being negligible at scales such as 10.0 GeV. A further study of the relation between the power correction and the confining potential is clearly needed.
### 3.3 On the relation with the lattice gluon condensate puzzle.
Having just stressed that the contribution of the power correction is rather large also at a scale such as $`10.0\mathrm{GeV}`$, this is a good point to go back to the argument referring to the unexpected result of . As already mentioned in the introduction, non-perturbative contribution to the running coupling can be advocated in order to explain the $`\mathrm{\Lambda }^2/Q^2`$ contribution to the gluon condensate. The argument runs as follows (see ). From general arguments one expects the condensate $`W`$ to be written in the form
$`{\displaystyle _0^{Q^2}}{\displaystyle \frac{p^2dp^2}{Q^4}}f({\displaystyle \frac{p^2}{\mathrm{\Lambda }^2}})`$ (12)
which is based on the fact that this condensate has dimension four and is renormalization group invariant, so that the function $`f(p^2/\mathrm{\Lambda }^2)`$, which is independent on $`Q`$ (for large $`Q`$), can be expressed as a function of a running coupling. This leads to consider the contribution coming from the large frequency part of the integral
$$_{\rho \mathrm{\Lambda }^2}^{Q^2}\frac{p^2dp^2}{Q^4}\alpha _s(\frac{p^2}{\mathrm{\Lambda }^2})$$
(13)
in which the function $`f(p^2/\mathrm{\Lambda }^2)`$ is taken proportional to the running coupling<sup>8</sup><sup>8</sup>8Note that higher powers would be subleading both for renormalons and for power corrections.. By taking into account the perturbative running coupling the IR renormalon contribution can be obtained (see for details). Again, one can consider also contributions coming from a non–perturbative correction to the coupling of the form $`c/p^2`$. From the UV limit of integration one then obtains a $`\mathrm{\Lambda }^2/Q^2`$ contribution to $`W`$. Let us insist: this is a contribution coming from a coupling with a $`c/p^2`$ correction in the UV region. While stressing again that one can draw no definite conclusion from the following consideration, still it is worth noting that our scheme provides an example of a coupling in which a $`c/p^2`$ correction is not negligible even at $`10.0\mathrm{GeV}`$. For a discussion of possible scenarios for the result of see also .
### 3.4 Landau pole and analyticity.
It is well known that perturbative QCD formulae for the running of $`\alpha _s`$ inevitably contain singularities, which are often referred to as the Landau pole. The details of the analytical structure depend on the order at which the $`\beta `$-function is truncated and on the particular solution chosen. The existence of an interplay between the analytical structure of the perturbative solution and the structure of non-perturbative effects has been advocated for a long time . To illustrate this idea, consider the one-loop formula for the running coupling $`\alpha _s(p)`$:
$$\alpha _s(p^2)=\frac{1}{b_0\mathrm{log}(\frac{p^2}{\mathrm{\Lambda }^2})}.$$
(14)
Here the singularity is a simple pole, which can be removed if one redefines $`\alpha _s(p)`$ according to the following prescription:
$$\alpha _s(p^2)=\frac{1}{b_0\mathrm{log}(\frac{p^2}{\mathrm{\Lambda }^2})}+\frac{\mathrm{\Lambda }^2}{b_0(\mathrm{\Lambda }^2p^2)},$$
(15)
where a power correction of the asymptotic form $`\frac{\mathrm{\Lambda }^2}{p^2}`$ appears. However, the sign of such a correction is the opposite of what one would expect from the results of and from the results in Section 2 (although the absolute value is of the right order), so that in the end one could envisage a more general formula for the regularized coupling:
$$\alpha _s(p^2)=\frac{1}{b_0\mathrm{log}(\frac{p^2}{\mathrm{\Lambda }^2})}+\frac{\mathrm{\Lambda }^2}{b_0(\mathrm{\Lambda }^2p^2)}+\eta \frac{\mathrm{\Lambda }^2}{p^2}.$$
(16)
The message from (16) is that the coefficient of the power correction is not constrained by the mere cancellation of the pole.
At higher perturbative orders one encounters multiple singularities, which include an unphysical cut. There are several ways to regularize them. In particular, the method discussed in combines a spectral-representation approach with the Renormalization Group. The method was originally formulated for QED, but it has recently been extended to the QCD case .
## 4 Conclusions
We have studied the strong coupling constant estimated non perturbatively on the lattice from Green functions in the Landau gauge using the $`\stackrel{~}{\mathrm{MOM}}`$ scheme. This has been performed with a large statistics of 1000 field configurations per run and running at $`\beta =6.0,6.2,6.4,6.8`$. Finite volume effects as well as finite lattice spacing effects have been carefully controlled.
We have parametrized the momentum dependence of $`\alpha _s`$ using the three-loop perturbative formula plus a $`c/p^2`$ term. We have obtain a good fit in the energy range from 2.0 GeV up to 10.0 GeV. As a result of this study we find
$`\mathrm{\Lambda }_{\overline{\mathrm{MS}}}=237\pm 3_{10}^{+0}\mathrm{MeV},c=0.63\pm 0.03_{0.13}^{+0.0}\mathrm{GeV}^2,`$ (17)
$`\mathrm{\Lambda }_{\overline{\mathrm{MS}}}`$ agrees perfectly well with the result of and the existence of sizable power corrections is convincingly established. The stability of our fit has been extensively checked.
The power correction turns out to be rather large, providing a 3% correction on $`\alpha _s`$ at 10.0 GeV, i.e. a 20 % correction on $`\mathrm{\Lambda }`$.
Having gathered evidences for a particular scheme, one needs to address the issue of the general relevance of such a finding in the spirit of the discussion of section 3, assessing the scheme dependence of our results. As already discussed, the non-perturbative nature of power corrections makes it very hard to formulate any theoretical procedure to estimate the impact of scheme dependence. The best one can do at this stage is to consider different renormalization schemes and definitions of the coupling and gather numerical evidence and formal arguments supporting power corrections to $`\alpha _s(p)`$. In this way, scheme-independent features may eventually be identified. For example, on the basis of our results, we note the following:
* Theoretical arguments suggest $`1/p^2`$ corrections both for the coupling as defined from the static potential and for the one obtained from the three-gluon vertex. The arguments for the former case were outlined in Sections 3.1 and 3.2. As far as the coupling from the three-gluon vertex is concerned, $`1/p^2`$ corrections appear in an OPE analysis if one keeps into account the fact that for such a gauge dependent coupling dimension 2 condensates are expected.
* In the static potential case, the theoretical arguments also provide an estimate for the order of magnitude of the coefficient of the $`1/p^2`$ correction while in the three-gluon vertex case the OPE arguments do not, suggesting instead that it may depend on the gauge. Nevertheless , the order of magnitude of our numerical result in the Landau gauge is roughly the same as the one from the static potential case. This calls for further investigation, which may be performed for example by attempting a similar calculation in a different gauge and particularly in the Coulomb gauge in which the static potential is naturally defined.
This issue of scheme dependence could be the focus of possible future work.
## 5 Acknowledgements
We thank Alain Le Yaouanc and Chris Michael for stimulating discussions. J. R-Q is indebted to Spanish Fundación Ramón Areces for his financial support. F. DR acknowledges both support from PPARC and from MURST (contract 9702213582) and INFN (i.s. PR11). C. Pi. warmly thanks the “Groupe de Phys. Nucl. Th. de l’Univ. de Liège” for kind hospitality and acknowledges the partial support of IISN. These calculations were performed on the QUADRICS QH1 located in the Centre de Ressources Informatiques (Paris-sud, Orsay) and purchased thanks to a funding from the Ministère de l’Education Nationale and the CNRS.
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# Towards a Stable Numerical Evolution of Strongly Gravitating Systems in General Relativity: The Conformal Treatments
## I Introduction
Three dimensional (3D) numerical relativity is an important technique for exploring the strong field dynamics in realistic astrophysical phenomena involving black holes and neutron stars. It is expected to play a role in analyzing gravitational waveforms to be observed soon, one expects, with the new generation of gravitational wave detectors going online worldwide in the next few years . However, progress in 3D numerical relativity, which has traditionally been based on the Arnowitt-Deser-Misner (ADM) system of evolution equations, has been slow. This is not only because of the immense computational difficulties that 3D simulations represent, but to a large extent it is due to severe instabilities often encountered during such simulations. Presently there is no complete understanding of the causes of these instabilities in numerical evolutions of the ADM equations. This has prompted much recent effort in developing alternative formulations of the 3+1 Einstein equations.
In this and a companion paper we focus on an alternative approach based on a conformal decomposition of the metric and the trace-free components of the extrinsic curvature. The conformal-tracefree (CT) approach was first devised by Nakamura in the 1980’s in 3D calculations , and then modified and applied to work on gravitational waves , and on neutron stars . This approach was not taken up by others in the community until a recent paper by Baumgarte and Shapiro , where a similar formulation was compared with the standard ADM approach and shown to be superior, in terms of both accuracy and stability, on tests involving weak gravitational waves, with geodesic and harmonic slicing. In a followup paper, Baumgarte, Hughes, and Shapiro applied the same formulation to systems with given (analytically prescribed) matter sources, and found similar stability properties. More recently fully hydrodynamical simulations employing the CT approach have been reported in in the context of collapse of rapidly-rotating (isolated) neutron stars and coalescence and merger of binary neutron stars. As we were preparing this manuscript we have also become aware of work by Lehner, Huq and Garrison where a comparison between the ADM and CT formulation in spherical symmetry has been carried out in the context of black hole excision.
In the companion paper we perform an analytic investigation of the stability properties of the ADM and the CT evolution equations. Using a linearized plane wave analysis, we identify features of the equations that we believe are responsible for the difference in their stability properties.
In this paper we report the results of simulations of weak and strong gravitational wave packets, black holes, boson stars and neutron stars in various slicing conditions, including maximal slicing and a family of algebraic slicings, and compare the results obtained by the ADM and CT equations in different implementations. We begin with a brief presentation of the relevant equations in Sec. II. We then discuss the results of our numerical simulations in section III. We consider vacuum spacetimes in section III A, and matter spacetimes in section III B. In section III A 1, we describe the various implementations of the CT equations using gravitational wave spacetimes as an example. We identify two particular implementations, which we call AFA and AF2, that give the best performance in long term evolutions. The essence of these implementations, is to “actively force” (AF, see below) the trace of the conformally rescaled extrinsic curvature (AFA), and for maximal slicing also the trace of the extrinsic curvature (AF2), to zero in each step of the numerical evolution. In the sections that follow, we focus on comparing the AFA and AF2 implementations to the results of the ADM equations for evolutions of strong field systems including black holes, boson stars and neutron stars. We demonstrate that for this wide range of systems, these two implementations of the CT equations always lead to more stable long term evolutions. However, we also find that for a given resolution, the ADM results are often more accurate than the CT results at early times, before the instabilities become apparent. We conclude with section IV. A study of the stability properties of the iterative Crank-Nicholson (ICN) scheme, used for the spacetime evolution of the simulations presented in this paper, can be found in the Appendix.
## II Formulation
We start reviewing briefly the formulations used for the comparisons.
The standard ADM equations are:
$`{\displaystyle \frac{d}{dt}}\gamma _{ij}`$ $`=`$ $`2\alpha K_{ij},`$ (1)
$`{\displaystyle \frac{d}{dt}}K_{ij}`$ $`=`$ $`D_iD_j\alpha +\alpha (\text{}R_{ij}+KK_{ij}`$ (3)
$`2K_{ik}K^k{}_{j}{}^{}{}_{}{}^{(4)}R_{ij}^{}),`$
with
$$\frac{d}{dt}=_t_\beta $$
(4)
and where $`_\beta `$ is the Lie derivative with respect to the shift vector $`\beta ^i`$. Here $`R_{ij}`$ is the Ricci tensor and $`D_i`$ the covariant derivative associated with the three-dimensional metric $`\gamma _{ij}`$. The 4-dimensional Ricci tensor $`{}_{}{}^{(4)}R_{ij}^{}`$ is usually written in terms of the energy density $`\rho `$ and stress tensor $`S_{ij}`$ of the matter as seen by the normal (Eulerian) observers:
$${}_{}{}^{(4)}R_{ij}^{}=8\pi \left[S_{ij}\frac{1}{2}\left(S\rho \right)\right].$$
(5)
The conformal, trace-free reformulations of these equations make use of a conformal decomposition of the three-metric, and the trace-free part of the extrinsic curvature. Here we follow closely the presentation of Ref. . The conformal three-metric $`\stackrel{~}{\gamma }_{ij}`$ is written as
$$\stackrel{~}{\gamma }_{ij}=e^{4\varphi }\gamma _{ij},$$
(6)
with the conformal factor chosen to be
$$e^{4\varphi }=\gamma ^{1/3}det(\gamma _{ij})^{1/3}.$$
(7)
In this way the determinant of $`\stackrel{~}{\gamma }_{ij}`$ is unity. The trace-free part of the extrinsic curvature $`K_{ij}`$, defined by
$$A_{ij}=K_{ij}\frac{1}{3}\gamma _{ij}K,$$
(8)
where $`K=\gamma ^{ij}K_{ij}`$ is the trace of the extrinsic curvature, is also conformally decomposed:
$$\stackrel{~}{A}_{ij}=e^{4\varphi }A_{ij}.$$
(9)
So far, these are just definitions of new variables, with no clear motivation for their introduction. Evolution equations for these new quantities are easy to find, and we summarize here the Baumgarte-Shapiro discussion on these equations, but with an emphasis on the possible numerical implications of various choices one can make.
The evolution equations for the conformal three–metric $`\stackrel{~}{\gamma }_{ij}`$, and its related conformal factor $`\varphi `$ are trivially written as
$`{\displaystyle \frac{d}{dt}}\stackrel{~}{\gamma }_{ij}`$ $`=`$ $`2\alpha \stackrel{~}{A}_{ij},`$ (10)
$`{\displaystyle \frac{d}{dt}}\varphi `$ $`=`$ $`{\displaystyle \frac{1}{6}}\alpha K.`$ (11)
The evolution equation for the trace of the extrinsic curvature $`K`$ can easily be found to be
$$\frac{d}{dt}K=\gamma ^{ij}D_iD_j\alpha +\alpha \left[\stackrel{~}{A}_{ij}\stackrel{~}{A}^{ij}+\frac{1}{3}K^2+\frac{1}{2}\left(\rho +S\right)\right],$$
(12)
where the Hamiltonian constraint was used to eliminate the Ricci scalar.
For the evolution equation of the trace-free extrinsic curvature $`\stackrel{~}{A}_{ij}`$ there are many possibilities. A trivial manipulation of Eq. (3) yields:
$`{\displaystyle \frac{d}{dt}}\stackrel{~}{A}_{ij}`$ $`=`$ $`e^{4\varphi }\left[D_iD_j\alpha +\alpha \left(R_{ij}S_{ij}\right)\right]^{TF}`$ (14)
$`+\alpha \left(K\stackrel{~}{A}_{ij}2\stackrel{~}{A}_{il}\stackrel{~}{A}_j^l\right),`$
but as shown previously there are many ways to write several of the terms, especially those involving the Ricci tensor. For example, one could eliminate the Ricci scalar $`R`$ again through the use of the Hamiltonian constraint.
With the conformal decomposition of the three–metric, the Ricci tensor now has two pieces, which we write as
$$R_{ij}=\stackrel{~}{R}_{ij}+R_{ij}^\varphi .$$
(15)
The “conformal-factor” part $`R_{ij}^\varphi `$ is given directly by straightforward computation of derivatives of $`\varphi `$:
$`R_{ij}^\varphi `$ $`=`$ $`2\stackrel{~}{D}_i\stackrel{~}{D}_j\varphi 2\stackrel{~}{\gamma }_{ij}\stackrel{~}{D}^l\stackrel{~}{D}_l\varphi `$ (17)
$`+4\stackrel{~}{D}_i\varphi \stackrel{~}{D}_j\varphi 4\stackrel{~}{\gamma }_{ij}\stackrel{~}{D}^l\varphi \stackrel{~}{D}_l\varphi ,`$
while the “conformal” part $`\stackrel{~}{R}_{ij}`$ can be computed in the standard way from the conformal three–metric $`\stackrel{~}{\gamma }_{ij}`$. To simplify notation, it is convenient to define what Ref. calls the “conformal connection functions”:
$$\stackrel{~}{\mathrm{\Gamma }}^i:=\stackrel{~}{\gamma }^{jk}\stackrel{~}{\mathrm{\Gamma }}_{jk}^i=\stackrel{~}{\gamma }_{,j}^{ij},$$
(18)
where the last equality holds if the determinant of the conformal three–metric $`\stackrel{~}{\gamma }`$ is actually unity (notice that this should be true analytically, but may not be numerically).
Using the conformal connection function, the Ricci tensor can be written:
$`\stackrel{~}{R}_{ij}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\stackrel{~}{\gamma }^{lm}\stackrel{~}{\gamma }_{ij,lm}+\stackrel{~}{\gamma }_{k(i}_{j)}\stackrel{~}{\mathrm{\Gamma }}^k+\stackrel{~}{\mathrm{\Gamma }}^k\stackrel{~}{\mathrm{\Gamma }}_{(ij)k}`$ (20)
$`+\stackrel{~}{\gamma }^{lm}\left(2\stackrel{~}{\mathrm{\Gamma }}_{l(i}^k\stackrel{~}{\mathrm{\Gamma }}_{j)km}+\stackrel{~}{\mathrm{\Gamma }}_{im}^k\stackrel{~}{\mathrm{\Gamma }}_{klj}\right).`$
Here again, one has choices in how the terms involving the conformal connection functions $`\stackrel{~}{\mathrm{\Gamma }}^i`$ are computed. A straightforward computation based on the Christoffel symbols could be used (and usually is in standard ADM formulations), but this approach leads to derivatives of the three–metric in no particular elliptic form. One would like to see an elliptic form as the principal part of this expression, as it brings the $`\stackrel{~}{\gamma }_{ij}\stackrel{~}{A}_{ij}`$ system a step closer to being hyperbolic. Thanks to the definition of the $`\stackrel{~}{\mathrm{\Gamma }}^i`$’s, an explicitly elliptic operator is singled out. However, if the terms involving the $`\stackrel{~}{\mathrm{\Gamma }}^i`$ are evaluated directly in terms of derivatives of the three–metric, this elliptic operator serves no special purpose, as other second derivatives appear through derivatives of the $`\stackrel{~}{\mathrm{\Gamma }}^i`$ which spoils the elliptic nature of the operator as a whole. If, on the other hand, the $`\stackrel{~}{\mathrm{\Gamma }}^i`$ are promoted to independent variables, for which evolution equations can be derived, then the expression for the Ricci tensor retains its elliptic character. The price to pay is that one must now evolve three additional quantities in the evolution system. Whether this has any numerical advantage will depend on details of the implementation, and will be discussed below.
Following this argument of promoting the $`\stackrel{~}{\mathrm{\Gamma }}^i`$ to independent variables, it is straightforward to derive their evolution equation:
$`{\displaystyle \frac{}{t}}\stackrel{~}{\mathrm{\Gamma }}^i`$ $`=`$ $`{\displaystyle \frac{}{x^j}}(2\alpha \stackrel{~}{A}^{ij}2\stackrel{~}{\gamma }^{m(j}\beta _{,m}^{i)}`$ (22)
$`+{\displaystyle \frac{2}{3}}\stackrel{~}{\gamma }^{ij}\beta _{,l}^l+\beta ^l\stackrel{~}{\gamma }_{,l}^{ij}).`$
However, again there is a choice one can make in writing this evolution equation; as pointed out in Ref. it turns out that the above choice leads to an unstable system. A choice which will be shown to be better can be obtained by eliminating the divergence of $`\stackrel{~}{A}^{ij}`$ with the help of the momentum constraint:
$`{\displaystyle \frac{}{t}}\stackrel{~}{\mathrm{\Gamma }}^i`$ $`=`$ $`2\stackrel{~}{A}^{ij}\alpha _{,j}+2\alpha (\stackrel{~}{\mathrm{\Gamma }}_{jk}^i\stackrel{~}{A}^{kj}`$ (25)
$`{\displaystyle \frac{2}{3}}\stackrel{~}{\gamma }^{ij}K_{,j}\stackrel{~}{\gamma }^{ij}S_j+6\stackrel{~}{A}^{ij}\varphi _{,j})`$
$`{\displaystyle \frac{}{x^j}}\left(\beta ^l\stackrel{~}{\gamma }_{,l}^{ij}2\stackrel{~}{\gamma }^{m(j}\beta _{,m}^{i)}+{\displaystyle \frac{2}{3}}\stackrel{~}{\gamma }^{ij}\beta _{,l}^l\right).`$
With this reformulation, in addition to the evolution equations for the conformal three–metric $`\stackrel{~}{\gamma }_{ij}`$ (10) and the conformal-traceless extrinsic curvature variables $`\stackrel{~}{A}_{ij}`$ (14), there are evolution equations for the conformal factor $`\varphi `$ (11), and the trace of the extrinsic curvature $`K`$ (12). If the $`\stackrel{~}{\mathrm{\Gamma }}^i`$ are promoted to the status of fundamental variables, as in Ref. , they can be evolved with (25). (Note that the mixed first and second order evolution system for $`\{\varphi ,K,\stackrel{~}{\gamma }_{ij},\stackrel{~}{A}_{ij},\stackrel{~}{\mathrm{\Gamma }}^i\}`$ is not in any immediate sense hyperbolic ). In the original formulation of Shibata and Nakamura , the auxiliary variables $`F_i=_j\stackrel{~}{\gamma }_{ij,j}`$ are used instead of the $`\stackrel{~}{\mathrm{\Gamma }}^i`$, and the final system of equations is somewhat more complicated.
Ref. shows that the CT system can also be interpreted as a “conformal second-order” version of the Bona-Massó system with $`2V_i=(\stackrel{~}{\mathrm{\Gamma }}_i+8_i\varphi )`$.
### A Gauge
Systems of the CT type have been investigated with various slicing conditions in the past. The paper of Baumgarte and Shapiro considered geodesic and harmonic slicing, while earlier work by Shibata and Nakamura, and the more recent paper by Baumgarte, Hughes, and Shapiro have also considered maximal slicing. Here we have studied maximal slicing and a number of algebraic slicings, and used them with different implementations of the CT equations, on numerical evolutions of many different spacetimes.
Maximal slicing has the property that $`K=0`$, leading to an elliptic equation for the lapse
$$^2\alpha =\alpha \left[K_{ij}K^{ij}+4\pi \left(\rho +S\right)\right].$$
(26)
Notice that in maximal slicing the evolution equations for $`\varphi `$ and $`K`$ become simply
$$d\varphi /dt=0,dK/dt=0.$$
(27)
The algebraic slicings that we will consider here correspond to the family originally introduced by Bona and Massó , building on earlier work of Bernstein
$$_t\alpha =f(\alpha )\alpha ^2K,$$
(28)
with $`f(\alpha )>0`$ but otherwise arbitrary. This family contains many well known slicing conditions. For example, taking $`f=1`$ we recover the “harmonic” slicing condition, which after a trivial integration becomes
$$\alpha =F(x^i)+\gamma ^{1/2},$$
(29)
with $`F`$ an arbitrary function of space. The name “harmonic” slicing comes from the fact that it corresponds to the choice of a harmonic time coordinate
$$\mathrm{}t=0.$$
(30)
Another useful slicing condition is obtained by taking $`f=N/\alpha `$. This corresponds to the generalized “1+log” slicing condition which after integration becomes
$$\alpha =F(x^i)+\mathrm{log}\gamma ^{N/2}.$$
(31)
(There is in fact some inconsistency in terminology as to whether the $`N=1`$ or the $`N=2`$ case corresponds to the standard “1+log” slicing; different choices being made by different authors.)
These type of algebraic slicings have an advantage over maximal slicing in terms of computational efficiency: It is much faster to integrate an evolution equation for the lapse than to solve an elliptic equation. On the other hand, such algebraic slicings are prone to the development of gauge pathologies . The possibility of the appearance of such pathologies when using algebraic slicings should always be kept in mind, as a gauge pathology can easily be confused with a numerical instability: one can lose a lot of sleep trying to cure an “instability” that is in fact a true solution of our system of differential equations.
To finish discussing our choice of gauge, we need to mention the fact that all the simulations described here have been carried out with the shift vector set to zero.
### B Boundary conditions
In standard 3+1 numerical simulations, the computational domain covers only a finite region of space. One must therefore apply some sort of artificial boundary condition at the edges of the numerical grid. Ideally, one would like to find a boundary condition that does not introduce numerical instabilities and allows gravitational waves to leave the grid cleanly, with no artificial reflections. This is in itself a very difficult problem, since in the first place, there is no local boundary condition that allows waves coming from any arbitrary direction to leave the grid with no reflections, and second, there does not even exist a clear way to define what a wave is in general relativity except at asymptotic infinity. In practice, what one looks for is a condition that remains stable and allows some “wave-like” solutions to leave the grid without introducing large reflections at the boundaries. The amount of artificial reflection that results typically depends on the specific form of the boundary condition, and on the direction of motion of the wave fronts as they hit the boundary .
Since in this paper we are interested in the question of the stability of the interior evolution, we will not worry too much about the boundary conditions, and we will limit ourselves to describing a few conditions that we have found to work well in practice. The conditions we have used are the following:
* Static boundary condition: The evolved variables are simply not updated at the boundary, and remain with their initial values there. This condition is very bad at handling waves since it reflects everything back in, but it can be very useful when studying situations that are supposed to remain static (as are some of the systems studied below), and where all the dynamics comes from numerical truncation errors.
* Zero-order extrapolation or “flat” boundary condition: After evolving the interior, the value of a given variable at the boundary is simply copied from its value one grid point in (along the normal direction to the boundary). This condition allows for some dynamics at the boundaries, and is somewhat better at absorbing waves than the static boundaries, but it still introduces a considerable amount of reflections.
* Sommerfeld or “radiative” boundary condition: In this case we assume that the dynamical variables behave like a constant plus an outgoing radial wave at the boundaries, that is:
$$f(x^i,t)=f_0+u(rt)/r,$$
(32)
where $`r=\sqrt{x^2+y^2+z^2}`$ and where the constant $`f_0`$ is taken to be one for diagonal metric components and zero for everything else. The radiative condition assumes that the boundaries are in the wave zone, where the speed of light is essentially one, and where the gravitational waves behave as spherical wavefronts. This boundary condition has been used before by other authors , and it has been found that in practice it is very good at absorbing waves.
It is in fact easier to implement a differential form of the radiative boundary condition than to use (32) directly. Consider a boundary that corresponds to a coordinate plane $`x_i`$=constant. Condition (32) then implies:
$$\frac{x_i}{r}_tf+_if+\frac{x_i}{r^2}(ff_0)=0.$$
(33)
One can now use simple finite differences to implement this last condition. In our code we have implemented condition (33) consistently to second order in both time and space.
* Robin boundary condition: This is a different type of “extrapolating” boundary condition, where one assumes that for large $`r`$ a given field behaves as:
$$f(x^i)=f_0+k/r,$$
(34)
with $`k`$ constant. This condition is clearly related to the radiative condition described above, but it contains no information about the time evolution. Just as we did with the radiative condition, we in fact implement the Robin condition in differential form:
$$_if+\frac{1}{r}\left(ff_0\right)=0.$$
(35)
The Robin boundary condition is usually better suited for solving elliptic problems than for use on dynamical variables.
Most of the simulations discussed below have been performed using the radiative boundary condition (33) for the dynamical variables, and the Robin boundary condition (35) both for constructing the initial data and for solving the maximal slicing condition. Whenever a different boundary condition is used, we say so explicitely.
## III Applications
In this section we will apply the previous system of conformal trace-free equations, exploring different implementations, in a series of numerical experiments with different spacetimes. The various implementations we consider are:
* Promoting the $`\mathrm{\Gamma }`$’s to independent variables.
* Use the momentum constraints on the evolution equation for the $`\mathrm{\Gamma }`$’s.
* Enforcing $`\mathrm{tr}\stackrel{~}{A}=0`$.
* For maximal slicing, enforcing $`\mathrm{tr}K=0`$.
We will study the effects of these different implementations using strong gravitational waves spacetimes.
All the numerical simulations presented here are carried out with the Cactus code for numerical relativity co-developed in our NCSA/Potsdam/Wash U collaboration and elsewhere.
### A Vacuum Spacetimes
We begin our discussion of the numerical simulations with vacuum spacetimes in this subsection, examining the evolution of both strong gravitational wave and black hole spacetimes. In particular, we use the gravitational wave simulations to illustrate the effects of the various implementations of the CT approach.
#### 1 Pure Gravitational Waves
We first turn to pure gravitational wave spacetimes. The low amplitude linear case has been studied, with a full 3D code, and published previously, (a) in both the standard ADM formulation and the Bona-Massó hyperbolic formulation by , where no fundamental differences were seen in performance at that time, and (b) by Shibata and Nakamura and Baumgarte and Shapiro in the CT approach as described above. The Baumgarte and Shapiro work particularly showed the strength of the CT formulation in the linearized case. Here we extend the study of these systems to include highly dynamic, strong field regimes. The study here is limited to tests that show the strengths and weaknesses of the different formulations. A study of the physics of collapsing waves in full 3D numerical relativity is presented elsewhere .
We consider here a three–metric of the form originally considered by Brill :
$$ds^2=\mathrm{\Psi }^4\left[e^{2q}\left(d\rho ^2+dz^2\right)+\rho ^2d\varphi ^2\right]=\mathrm{\Psi }^4\widehat{ds}^2,$$
(36)
where $`q`$ is a free function subject to certain regularity and fall-off conditions. Different forms of the function $`q`$ have been considered by different authors , but most work so far has concentrated only in constructing and analyzing the initial data.
As in Ref. , we use a generalized form for the function $`q`$, giving it a full 3D dependence, following :
$$q=a\rho ^2e^{r^2}\frac{(1+c\rho ^2\mathrm{cos}^2(\mathrm{m}\theta ))}{(1+\rho ^2)},$$
(37)
where $`a`$ and $`c`$ are constants, $`r^2=\rho ^2+z^2`$ and $`m`$ is an integer. In this paper we focus on the axisymmetric case, $`c=0`$, for simplicity, although using a non-zero value of $`c`$ does not affect the results we discuss below. All the runs discussed here where performed using an iterative Crank-Nicholson (ICN) scheme with 3 iterations (see appendix), and radiative boundary conditions.
The first case presented is an initial configuration with amplitude $`a`$=4, corresponding to a strong wave, but not quite strong enough to collapse to a black hole. In the evolution of this data set the wave implodes through the origin, oscillates a few times, and finally disperses back to infinity leaving flat space behind, but in a non-trivial spatial coordinate system . The evolution of this spacetime is highly non-linear, and the final configuration has metric components with a large dynamical range.
In Fig. 1a we show the evolution of the minimum value of the lapse over the grid for a simulation done with the standard ADM formulation, using maximal slicing, no shift and a radiative boundary condition. For this particular simulation we used a resolution of $`\mathrm{\Delta }x`$=0.08 and $`67^3`$ grid points. Also, we used the fact that our data is symmetric across coordinate planes to evolve only one octant. The simulation crashes at $`t8`$ when the lapse collapses catastrophically in response to a blow up of the extrinsic curvature. Fig. 1b shows the evolution of the maximum value of the trace of the extrinsic curvature $`K`$. Notice that even though we are using maximal slicing, $`K`$ does not remain zero, and blows up towards the end of the simulation. The fact that $`K`$ does not remain zero is not surprising, since the maximal slicing condition is solved numerically, and thus a residual time derivative of $`K`$ is to be expected. The catastrophic blow-up, however, is a different matter and points towards the existence of an unstable solution of our system of equations.
Fig. 2 shows the same simulation, but now using the so-called “K-driving” technique . The idea here is to add counter terms to the elliptic equation for the lapse to drive the numerically produced non-zero $`K`$ (the trace of the extrinsic curvature) back towards zero. With K-driving, $`K`$ remains much smaller until close to the point of crashing at $`t`$=9, with a catastrophic blow-up of the lapse at the end. This shows that a better control of the time slicing is not enough to cure the instability in the evolution: There exist unstable modes that are not controlled by keeping the value of $`K`$ small. (For an analysis of possible unstable modes of the ADM equations, see .)
Next, we show the evolution of the same system using again maximal slicing, and different implementations of the CT formulation. In Fig. 3 we show again the central value of the lapse for the same initial data. The different runs correspond to the following cases:
| | use of | use momentum | force | remove |
| --- | --- | --- | --- | --- |
| | $`\mathrm{\Gamma }^i`$ | constraints | $`K`$=0 | $`\mathrm{tr}\stackrel{~}{A}`$ |
| Res | no | - | no | no |
| Gam | yes | no | no | no |
| Mom | yes | yes | no | no |
| AFK | yes | yes | yes | no |
| AFA | yes | yes | no | yes |
| AF2 | yes | yes | yes | yes |
The first run uses the implementation denoted “Res” (for rescale). It differs from the standard ADM equations only in the conformal rescaling and the fact that $`\varphi `$ and $`K`$ (which enter into the evolution equation for $`\stackrel{~}{A}_{ij}`$) are now evolved separately. The second run, with the implementation denoted “Gam” (for gamma), introduces the $`\mathrm{\Gamma }^i`$, but does not use the momentum constraints to rewrite their evolution equations. The third run uses the implementation “Mom” (for momentum constraints) and represents a straightforward coding of the the full set of CT equations , where the momentum constraints are used to modify the evolution equations for the $`\mathrm{\Gamma }^i`$, but without adding anything else. In the fourth run, which uses the implementation “AFK” (for “actively enforcing K”), we have forced $`K`$ to remain zero by simply not evolving it, and we have also kept $`\varphi `$ time independent (see Eq. (11)). In the fifth run we use the implementation “AFA”, where we allow $`K`$ to evolve freely, but actively force $`\stackrel{~}{A}`$ (the trace of $`\stackrel{~}{A}_{ij}`$) to remain zero by subtracting it from $`\stackrel{~}{A}_{ij}`$ after each time step:
$$\stackrel{~}{A}_{ij}\stackrel{~}{A}_{ij}\frac{1}{3}\stackrel{~}{\gamma }_{ij}\mathrm{tr}\stackrel{~}{A}.$$
(38)
And finally, in the sixth run we use the implementation “AF2” that combines implementations AFK and AFA above by actively enforcing both $`K`$=0 and $`\stackrel{~}{A}`$=0. Notice that both $`K`$ and $`\stackrel{~}{A}`$ should be zero in principle in an exact evolution using the CT equations with maximal slicing, but they do not remain so in actual numerical evolutions unless actively enforced.
As can be seen from the figure, runs Res, Gam, Mom, AFK and AFA eventually crash, but run AF2 with double active enforcement does not, at least for the time scale under study. The lapse returns to unity, and the final static spacetime can be followed for a long time with no sign of an instability (we have in fact followed run AF2 past $`t`$=100 and it still remains stable). From the figure we also see that by enforcing only $`K`$=0 or $`A`$=0 separately, as is done in runs AFK and AFA, one still obtains improved stability, with the simulations crashing at late times after the lapse has already returned to 1. This shows that by enforcing only one of the two constraints, and keeping the other options turned on, we still get a rather robust system when compared to standard ADM. Moreover, enforcing $`\stackrel{~}{A}`$=0 appears to be more important than enforcing $`K`$=0, as can be seen from the fact that run AFA crashes much later than run AFK.
Finally, notice that run Gam crashes even sooner than run Res, which shows that it is in fact better not to use the $`\mathrm{\Gamma }^i`$ than to use them without modifying their evolution equation. For understanding the need to use the momentum constraints in the CT approach, see the companion paper .
We note that the results found above for the different implementation are generic for strong gravitational wave spacetimes, quite independent of the precise parameter choices. However, for weak gravitational waves in the linear regime, the straightforward coding of the CT equations (implementation “Mom”) leads also to stable evolutions as do the AFK, AFA and AF2 cases. In Fig. 4 we show again the minimum value of the lapse for the evolution of a wave with an amplitude of $`a=0.01`$, using the ADM formulation and also the Mom, AFK and AFA implementations of the CT system (since the lapse remains very close to 1, we are in fact plotting $`(\alpha 1)\times 10^5`$). We see that while the ADM run crashes at an early time ($`t15`$) with a catastrophic collapse of the lapse, all three implementations of the CT equations give stable evolutions and yield basically the same results for a weak wave. We have followed these three runs past $`t`$=100 with no instabilities developing (the AF2 implementation is in fact just as stable, but we don’t include it in the figure).
From these studies (and many others with different parameters that we have done) we can conclude that, for maximal slicing, the CT formulation has better stability properties for the evolution of strong field systems, as long as:
* The $`\mathrm{\Gamma }^i`$ are promoted to independent variables.
* The momentum constraints are used to transform the evolution equation for the $`\mathrm{\Gamma }^i`$. Evolving the $`\mathrm{\Gamma }^i`$ without modifying their evolution equation is worse than not using them at all.
* The trace of the extrinsic curvature $`K`$ is actively forced to be zero (the definition of maximal slicing).
* The trace of the $`\stackrel{~}{A}_{ij}`$ is also actively forced to be zero.
So far we have focused on the issue of long term stability. Now we want to compare accuracy of the CT and ADM formulations. We concentrate on the best implementation of the CT equations, the one we labelled AF2. In Fig. 5 we show the L2-norms of the hamiltonian constraint for the $`a`$=0.01 and $`a`$=4 cases discussed above, using the ADM (solid line) and the AF2 systems (dashed line). In both cases we see that for the ADM system, the L2-norm of the hamiltonian constraint grows more or less linearly for some time (this is more evident in the $`a`$=0.01 case) until just before the crash when it blows up catastrophically. In contrast, in the AF2 runs the L2-norm of the hamiltonian constraint initially grows faster, but it later settles on a constant value. The fact that the ADM runs are more accurate than the AF2 runs at early times appears to be quite generic: we have found essentially the same behavior for all the different parameters that we have studied.
We have also performed convergence tests by running the same initial data with different resolutions, and we have found that both the ADM and AF2 evolutions are second order accurate. As an example of this, Fig. 6 shows the L2-norms of the hamiltonian constraint for both the ADM and the AF2 systems for two different resolutions: The dashed lines show the L2 norm for a resolution of $`dx`$=0.16 ($`35^3`$ grid points), while the solid lines show the L2 norm for a resolution of $`dx`$=0.08 ($`67^3`$ grid points) multiplied by a factor of four. For second order convergence the solid and dashed lines should fall on top of each other. From the figure we see that this is indeed true for most of the run in both cases. For the ADM run, second order convergence starts to fail shortly before the crash. On the other hand, for the AF2 run we obtain slightly degraded convergence (but still better than first order) for times between $`t`$=5 and $`t`$=15 when the spacetime is very dynamic, indicating that we haven’t quite reached the second order convergence regime for the resolutions considered here.
Though in this section we have concentrated in the case of maximal slicing, we should mention that we have also performed many simulations using the generalized “1+log” slicings. The results are in fact very similar to those reported here, except for the fact that implementations AFK and AF2 can no longer be used (since $`K`$ in non-zero for these slicing conditions). We find that for these algebraic slicings, implementation AFA is by far the best performer.
In the following subsections, we show that the above results on the stability and accuracy of the ADM and CT systems are basically the same for systems ranging from black holes to spacetimes coupled to dynamical source fields.
#### 2 Black Holes
Black holes have been the target of an intense research effort in recent years in numerical relativity, and have proved particularly difficult to handle in 3D evolutions. In the “standard” numerical evolution of black holes using the ADM equations together with singularity avoiding slicings, 3D simulations generally develop instabilities and crash before $`t=50M`$, where $`M`$ is the mass of the system . This falls far short of the time required to model the complete inspiral of two black holes, or even the head-on collision. Still, singularity avoiding slicings combined with the ADM equations make it possible to evolve through a brief part of the merger phase of two black holes with momenta and spins, and from this point of view give the most generally applicable method available. Future cures for grid stretching are expected to be based on black hole excision or characteristic slicings .
In the following we carry out a preliminary study of the CT formulation in black hole evolutions with grid stretching. It is inevitable that the sharp peaks that develop in the metric function due to grid stretching will cause the code to crash at some point in the evolution. We consider the evolution of the Misner data as a concrete example. The 3D numerical evolution of the Misner data in the standard ADM setting with singularity avoiding slicing has previously been studied using the so-called “G” code and its derivatives , developed by the NCSA/WashU group. Comparable results for a single black hole can be found in .
In Fig. 7 and Fig. 8 we compare the results of evolutions of Misner data with the separation parameter $`\mu =2.2`$, corresponding to two initially well separated black holes, on a grid of size $`130^3`$ with grid spacing $`0.08`$. The only difference in the simulations is the system of equations used to carry out the evolution (ADM vs. AF2); all computational parameters, such as parameters in the ICN finite differencing scheme, grid parameters, radiative boundary conditions, and maximal slicing condition are the same. In Fig. 7, first panel, we show the radial-radial metric component along a line on the equatorial plane at various times for the ADM case. We can clearly see the familiar ever-growing peak caused by the grid stretching associated with singularity avoiding slicings. In the first panel of Fig. 8 we show the lapse function along a line on the equatorial plane at various times for the ADM case, and here an instability becomes apparent at around $`t=14M`$ which is not yet reflected in the metric. This short wave length instability grows rapidly and causes the code to crash at $`t=14M`$. In the second panel we show the AF2 case. No metric instability is seen until towards the end of the simulation at $`t=24M`$, although the peak appears to be deformed. At this time the radial metric function peak has grown to about two times higher than that attained in the ADM case. The lapse for the AF2 case in Fig. 8 does not show an instability.
However, note that a smooth and stable evolution of the lapse does not mean that the computed data is still useful. To emphasize this point, Fig. 9 shows the same run as above with AF2 on a smaller grid with only $`66^3`$ points, but with the same grid spacing as before (so the boundaries are much closer in). While ADM crashes when the gradients in the metric become too severe, the AF2 run is able to continue with a smooth lapse even after the metric becomes deformed (cmp. where the evolution of the metric is not discussed). The lapse eventually collapses in the whole grid, freezing the evolution (so one could keep running “forever”, but the evolution becomes meaningless).
Next, we compare the accuracy of both simulations. In Fig. 10 we show the L2 norm of the Hamiltonian constraint for a grid size of $`130^3`$. The dashed line represents the ADM run, and the solid line the AF2 run. We see that the ADM results are more accurate than the AF2 results until just after time $`t=14M`$, when the instability in the ADM evolution begins to dominate and the code crashes (with higher resolution this crash time can be delayed somewhat). Starting at around $`t=20M`$ for AF2, there is a spurious growth in the Hamiltonian constraint that corresponds to the deformation in the metric. For maximal slicing one expects continuous growth of a smooth metric peak, but with AF2 the shoulder in the lapse seems to overtake the outward movement of the metric peak, freezing its growth in an irregular manner.
These results for black holes with grid stretching cannot be compared directly to the wave runs in the previous section because in the case of the black hole runs we do not approach a static final state. However, the CT formulation still offers some advantages over ADM in achievable run time. We find stability far beyond were the runs are meaningful, and it remains to be explored how far one can push the CT runs while maintaining convergence.
### B Matter Spacetimes
In the previous sections we studied the stability properties of the vacuum Einstein equations. What will happen if these equations are coupled to dynamical matter sources that are themselves governed by evolution equations coupled to the spacetime geometry? The complete set of equations can now have more complicated types of unstable modes. What would be the effects of switching from the ADM formulation to the CT formulation?
To respond to this question we consider next the following systems: (i) the evolution of boson stars governed by the scalar field Klein-Gordon equation and (ii) the evolution of neutron stars governed by the hydrodynamical equations (general relativistic Euler equations). The numerical evolution of the Klein-Gordon equation is straightforward with many well-known stable schemes. However, the numerical evolution of the hydrodynamical equations is considerably more challenging, especially in the presence of shocks or highly relativistic flows. For this purpose we use a recently developed hydrodynamical code which employs a conservative formulation of the equations together with high-resolution shock-capturing (HRSC) schemes based on approximate Riemann solvers. In we demonstrated that this code is capable of handling hydrodynamical evolutions in a stable and accurate fashion for a range of scenarios.
We focus here on analyzing the stability and accuracy of evolutions of both static boson stars and static neutron stars using the ADM formulation and the AFA implementation of the CT equations discussed above. We use the AFA implementation rather than AF2 because the simulations discussed here have all been performed using algebraic slicings and implementation AF2 applies only to maximal slicing. The main motivation for this has been the fact that, as we will show below, implementation AFA with algebraic slicings already gives excellent results when compared with standard ADM and is far less computationally expensive than runs that use maximal slicing.
#### 1 Boson Stars
We begin with a simple kind of matter source: self-gravitating scalar fields. This system has served as a useful testbed for developing numerical techniques for dealing with relativistic matter coupled to the Einstein equations , and also has a distinguished history in the field, having provided the first example of critical phenomena in relativity .
The dynamics of a massive scalar field are described by the minimally coupled Klein-Gordon (KG) equation
$$\mathrm{}_g\varphi =m^2\varphi ,$$
(39)
(see, e.g. ). The KG equation can be obtained from the Lagrangian
$$=\frac{1}{2}g^{\mu \nu }\varphi ,_\mu \varphi ^{},_\nu +\frac{1}{2}m|\varphi |^2,$$
(40)
which leads to the stress-energy tensor
$$T_{\mu \nu }=\frac{2}{\sqrt{g}}\frac{\delta }{\delta g^{\mu \nu }},$$
(41)
which is used as the matter source for the Einstein equations.
Self-gravitating massive scalar fields have bound, star-like solutions called boson stars with stability properties very much like those of neutron stars. These objects have been studied numerically, extensively in 1D and also in 3D . Apart from the fact that their evolution equation is much simpler than the hydrodynamical equations, boson stars are also easier to handle numerically when compared to neutron stars because they have no sharp changes in the density distribution near the surface layer of the star. For more details on the properties of boson stars and their behavior under perturbations see and references cited therein.
We perform our numerical evolutions of boson stars by writing the KG equation as a flux-conservative system of the form
$$\dot{u}_a=_bF_a^b+S_a^bu_b$$
(42)
where $`\stackrel{}{u}`$ contains the scalar field and its time and space derivatives. The method used to integrate this equation is a symmetrized MacCormack with both directional and Strang splitting. Symmetrized here means that the order of left-hand and right-hand differencing changes every time step (this improves the stability of the scalar field evolution). The code for solving the KG equation converges to second order in time and space. See Ref. for details of the numerical methods.
We have carried out evolutions of equilibrium boson star configurations with the metric background held fixed artificially (not updating the metric functions), and evolutions of the metric of such configurations with the scalar field held fixed artificially (not updating the scalar field), for a range of compactness of the boson stars, using both the ADM and AFA schemes. For all these cases, we have seen that the simulations are stable and second order convergent. The case of coupled spacetime-scalar field evolution is much more challenging, and we focus on that below.
We begin by showing an equilibrium boson star with a central density near the maximum stable value (field strength at center $`\varphi _0=0.26`$, total mass $`M=0.6322m_{p}^{}{}_{}{}^{2}/m`$, with $`m_p`$ the Planck mass, $`m`$ the mass of the scalar field). In Fig. 11, we show the evolution of radial metric component $`g_{rr}`$ in a fully coupled simulation, using a three step ICN scheme, 1+log slicing with $`N=2`$, no shift, a radiative boundary condition on the metric, and a flat boundary condition on the scalar field. A $`32^3`$ grid is used to cover only one octant. In the first panel we show the results of the ADM evolution. We see that for a short time, the spacetime remains nearly static (as it should). However, a short wavelength instability becomes significant by time $`t`$=7, and quickly grows causing the code to crash. The time $`t`$ here is expressed in terms of the intrinsic oscillation time scale of the scalar field (the exact equilibrium boson star field has the form $`\psi (r)e^{it}`$). In the second panel we show the evolution with exactly the same setup but using now implementation AFA instead of ADM. We see that the static configuration is maintained for a much longer time. Towards the end of the evolution, near $`t=150`$, we see that numerical error starts to build-up near the boundary of the computational domain.
In Fig. 12 below, we compare the L2-norm of the hamiltonian constraint for the ADM and AFA runs. We see that at early times the ADM run gives a more accurate result, but instabilities cause the L2-norm to blow by $`t8`$. For the AFA run the constraint violation is larger at first, but the evolution remains stable or a much longer time. The oscillation of the hamiltonian constraint we see here can be understood as a reaction of the scalar field to the numerical truncation error, which can be interpreted as a kind of perturbation. The frequency of these oscillations coincides with the ones obtained in 1D studies of perturbed boson stars. Notice that with the ADM run the code crashes so early that one can not even see the first oscillation.
#### 2 Static Neutron Stars
We turn now to the study of hydrodynamical evolutions of neutron stars. In we developed a three-dimensional, fully relativistic code to integrate the hydrodynamical equations coupled to the ADM equations. Convergence studies using polytropic neutron stars showed that the code is second order accurate in both space and time. For the integration of the hydrodynamical equations we used HRSC schemes of the total-variation-diminishing (TVD) class, with a piecewise-linear reconstruction of a sufficient set of hydrodynamical variables (rest-mass density, three-velocity and internal energy density). For more details on the schemes available in the code, see . In the studies reported in this paper we use the ICN scheme for the integration of the spacetime equations (either ADM or AFA) and Roe’s approximate Riemann solver for the hydrodynamical equations. We use “1+log” slicing with $`N=2`$.
As in the boson star studies we have first considered evolutions which test separately the individual components of the code. In these, we either solve the hydrodynamical equations in a prescribed (static) spacetime or the gravitational field equations for a prescribed matter source. In particular, we have evolved static neutron star configurations with a zero-temperature polytropic equation of state, of the form $`P=K\rho ^\mathrm{\Gamma }`$ (where $`P`$ is pressure and $`\rho `$ is rest-mass density). This included stars with a large polytropic index $`\mathrm{\Gamma }`$ (very stiff) having density profiles with a discontinuous first derivative at the surface. In the case of prescribed matter sources, we have confirmed that the comparison of the AFA and AF2 systems to the ADM system, in terms of stability and accuracy, remains the same as in the vacuum cases studied above. Static neutron stars with polytropic index $`\mathrm{\Gamma }=2`$ have also been studied in using the CT equations with prescribed hydrodynamical sources.
We focus next on the coupled spacetime and hydrodynamical evolution of static Tolman-Oppenheimer-Volkoff (TOV) neutron stars (in isotropic coordinates). Again, we compare the results obtained using the AFA implementation of the CT equations to those of the ADM equations. In principle both the matter distribution inside the star and the spacetime should remain static. In practice they evolve due to truncation errors of the finite-difference scheme, with the hydrodynamics and the spacetime responding to one another. The static TOV solution provides a reference to monitor the accuracy of the coupled numerical evolution. Note that in these evolutions, static outer boundary conditions were used.
In Fig. 13, we show the evolution of the L2-norm of the Hamiltonian constraint for a polytropic, $`N=1`$, TOV star of gravitational mass $`1.4M_{}`$ and compactness ratio $`M/R=0.146`$. A $`64^3`$ grid is used to cover the first octant, with $`dx=dy=dz=0.34`$km. The dashed line corresponds to the ADM system and the solid line to the AFA system. Again, as in the vacuum studies, we see that the ADM evolution suddenly becomes unstable at roughly 2.7ms, while the AFA evolution remains stable after more than 6ms (we followed the evolution for more than twice that).
In Fig. 14 we show the evolution of the radial component of the metric (constructed from the evolved Cartesian metric components). The first panel of Fig. 14 corresponds to the evolution obtained with ADM. We see that the star basically maintains its initial equilibrium, until the high-frequency instability crashes the code. In the second panel, we show $`g_{rr}`$ at various times, obtained with the AFA implementation. All other parameters are the same as in the ADM evolution. The ADM run is more accurate, before it becomes unstable, while the AFA run is stable but less accurate (there is a secular drift away from the initial configuration).
The truncation errors of the coupled evolution code initiate a pulsation of the star in, mainly, its radial modes of pulsation. These pulsations are damped in time due to the viscosity of the numerical scheme (see ). The TVD schemes we are using describe well the physical pulsations of the fluid, except in a small region around the center of the star, where short wavelength noise appears in the radial velocity. Our trials with other HRSC schemes show that this behavior seems to be generic for higher order HRSC schemes <sup>*</sup><sup>*</sup>*We have extensively experimented with other hydrodynamical evolution schemes. If one uses a first-order (Godunov) scheme, using piecewise constant reconstructed data for the Riemann problem, instead of piecewise linear, the radial velocity oscillates around zero near the center of the star, without any short wave length noise. But, a low-order scheme is not capable of accurately describing the evolution of the stellar surface where the density distribution is changing rapidly (unless prohibitively large grids are used) and large errors from the surface layers soon propagate inside the star. We have also experimented with a mixed system: first-order near the center and second-order near the surface. In this case the error grows at the interface of the two regime, yielding a even less accurate evolution overall.. In all such schemes, the radial momentum near the center has a small residual value of constant sign. This momentum appears in the r.h.s. of the evolution equation for $`\stackrel{~}{\mathrm{\Gamma }}^i`$ (Eq. (25)). This, in turn, leads to an error in the spacetime evolution. It is noteworthy that this does not cause an instability in the coupled evolution, except at very late times, when the violation of the Hamiltonian constraint has already become extremely large.
We note that as the TVD schemes are only first-order accurate at local extrema, such as the maximum of the density at the center of the star, so the increase in the Hamiltonian constraint at the center converges to roughly first order with increasing resolution. Away from the center, the scheme is second order convergent. The convergence of the L2-norm of the Hamiltonian constraint with the AFA system, for different grid-sizes (and for the same initial configuration as above), is shown in Fig. 15.
## IV Discussion and Conclusions
In this paper we have studied the stability of three-dimensional numerical evolutions of the Einstein equations in a formulation that separates out the conformal and traceless parts of the system. In our study we have considered different spacetimes including gravitational waves, black holes, boson stars and neutron stars.
We investigated several implementations of the conformal-traceless (CT) evolution equations. We identified two of them which give the best long term stability behavior: the AF2 implementation for maximal slicing, and the AFA for algebraic slicings. The AFA implementation actively enforces the trace of the conformally rescaled extrinsic curvature ($`\stackrel{~}{A}`$) to zero at each step of the time evolution, while the AF2 implementation enforces as well the fact that the trace of the extrinsic curvature ($`K`$) should vanish in maximal slicing. On the analytic level, the CT evolution equations imply that $`\stackrel{~}{A}=0`$ throughout the evolution, but this is inevitably violated in numerically evolution due to truncation error, unless actively enforced. Similarly, for maximal slicing, $`K`$ will not remain zero numerically unless actively forced to do so. We find that these two implementations of the CT equations lead to a more stable evolution compared to what one can obtain using the standard ADM evolution equations, under the same resolution, boundary condition and grid parameter choices, for all systems investigated. In comparison, a straightforward implementation of the CT equations (“Mom”) is capable of giving a stable evolution for weak but not strong field systems. We should also mention that we have recently become aware of the work of Lehner, Huq and Garrison where a comparison of the ADM and CT formulations has been carried out and where it is also found that freezing the evolution of $`K`$ (what these authors call “locked evolution”) improves considerably the stability of simulations that use the CT formulation.
Beyond stability, we have also compared the accuracy of the evolutions obtained by the ADM equations and CT equations. For all spacetimes considered we have found that the ADM system is consistently more accurate than the CT system in short term evolutions, before the instabilities set in. Although at present we can offer no explanation of this difference in accuracy between the different formulations, we believe that it is not a consequence of our numerical implementation, but is rather a property of the system of differential equations. It therefore points in the direction for a possible improvement of the CT approach. We note that formulations combining the CT approach and the hyperbolic approach have been proposed . A similar investigation of the stability and accuracy properties of such formulations will be presented elsewhere.
In this paper we have focused on the implementations and the numerical properties of their evolutions. Some understanding of the different stability of properties on the analytic level is discussed in a companion paper .
###### Acknowledgements.
We would like to thank many colleagues for discussions that have aided the development of this work. We are especially grateful to Mark Miller, Malcolm Tobias and Wai-Mo Suen of the Washington University gravity group, to Toni Arbona, Carles Bona, and Joan Massó of the Universitat de les Illes Balears, to Gabrielle Allen, Gerd Lanfermann and Daniel Holz at the AEI, and to Vince Moncrief. The research was supported by AEI, NCSA, the NSF grant Phy 96-00507, NSF NRAC MCS93S025, and NASA HPCC Grand Challenge Grant NCCS5-153. J.A.F. acknowledges financial support from a TMR grant from the European Union (contract nr. ERBFMBICT971902).
## A Stability Analysis of the Iterative Crank-Nicholson Scheme
The numerical scheme used for the simulations described in this paper is the so-called iterative Crank-Nicholson (ICN) scheme, which is an iterative, explicit version of the standard implicit Crank-Nicholson (CN) scheme . The idea behind this method is to solve the implicit equations by an iterative procedure, where each iteration is an explicit operation depending only on previously computed data. Normally, this process is stopped after a certain number of iterations, or until some tolerance is achieved. For a linear equation (and in particular in one dimension), the iterative procedure can easily be much more computationally expensive than the matrix inversion required to solve the original implicit scheme. For a non-linear system, however, solving the implicit scheme directly can prove to be extremely difficult.
In this appendix we study the stability properties of the ICN scheme in the particular case of the simple wave equation, and derive two very important results:
* In order to obtain a stable scheme one must do at least three iterations, and not just the two one would normally expect (two iterations are enough to achieve second order accuracy, but they are unstable!).
* The iterative scheme itself is only convergent if the standard Courant-Friedrichs-Lewy (CFL) stability condition is satisfied, otherwise the iterations diverge.
These two results taken together imply that there is no reason (at least from the point of view of stability) to ever do more that three ICN iterations. Three iterations are already second order accurate, and provide us with a (conditionally) stable scheme. Increasing the number of iterations will not improve the stability properties of the scheme any further. In particular, we will never achieve the unconditional stability properties of the full implicit CN scheme, since if we violate the CFL condition the iterations will diverge. As we were finishing this manuscript we became aware of a paper by S. Teukolsky were he does essentially the same analysis and obtains the same results . His analysis and ours complement each other, since he considers any finite number of iterations, while we consider only 1, 2 and 3 iterations. On the other hand, here we also consider the question of the convergence properties of an infinite number of iterations.
For our stability analysis we will consider the simple wave equation in N-dimensions. Numerical experiments have shown that the full Einstein equations have essentially the same stability properties.
Consider then the N-dimensional wave equation written in “3+1 like” form:
$$_t\varphi =A,_tA=\underset{i=1}{\overset{N}{}}_i^2\varphi .$$
(A1)
For the finite difference approximation to these equations we employ the usual notation
$$f_𝐦^n:=f(x_i=m_i\mathrm{\Delta }x,t=n\mathrm{\Delta }t),$$
(A2)
with $`n`$ and $`𝐦=(m_1,\mathrm{},m_N)`$ integers. The implicit CN scheme is then given by
$`\varphi _𝐦^{n+1}`$ $`=`$ $`\varphi _𝐦^n+{\displaystyle \frac{\mathrm{\Delta }t}{2}}\left(A_𝐦^{n+1}+A_𝐦^n\right),`$ (A3)
$`A_𝐦^{n+1}`$ $`=`$ $`A_𝐦^n+{\displaystyle \frac{\mathrm{\Delta }t}{2(\mathrm{\Delta }x)^2}}{\displaystyle \underset{i=1}{\overset{N}{}}}\delta _i^2\left(\varphi _𝐦^{n+1}+\varphi _𝐦^n\right),`$ (A4)
where the finite difference operators $`\delta _i^2`$ are defined as
$$\delta _i^2f_{m_i}^n:=f_{m_i+1}^n2f_{m_i}^n+f_{m_i1}^n.$$
(A5)
The implicit CN scheme is well known to be unconditionally stable for the wave equation (i.e. stable for any value of $`\mathrm{\Delta }t`$).
The ICN scheme is defined in the following way
$`\varphi _𝐦^{(1)}`$ $`=`$ $`\varphi _𝐦^n+\mathrm{\Delta }tA_𝐦^n,`$ (A6)
$`A_𝐦^{(1)}`$ $`=`$ $`A_𝐦^n+\mathrm{\Delta }t{\displaystyle \underset{i=1}{\overset{N}{}}}\varphi _𝐦^n,`$ (A7)
$`\varphi _𝐦^{(i)}`$ $`=`$ $`\varphi _𝐦^n+{\displaystyle \frac{\mathrm{\Delta }t}{2}}\left(A_𝐦^{(i1)}+A_𝐦^n\right),`$ (A9)
$`A_𝐦^{(i)}`$ $`=`$ $`A_𝐦^n+{\displaystyle \frac{\mathrm{\Delta }t}{2(\mathrm{\Delta }x)^2}}{\displaystyle \underset{i=1}{\overset{N}{}}}\delta _i^2\left(\varphi _𝐦^{(i1)}+\varphi _𝐦^n\right),`$ (A10)
and finally,
$`\varphi _𝐦^{n+1}`$ $`=`$ $`\varphi _𝐦^{(i_{\mathrm{max}})},`$ (A11)
$`A_𝐦^{n+1}`$ $`=`$ $`A_𝐦^{(i_{\mathrm{max}})},`$ (A12)
From these expressions it is clear that if the iterations converge, we will recover the implicit CN scheme.
For the stability analysis of the ICN scheme we use the standard von Neumann ansatz
$`\varphi _𝐦^n`$ $`=`$ $`\xi _1\lambda ^ne^{i(𝐤𝐦)\mathrm{\Delta }x},`$ (A13)
$`A_𝐦^n`$ $`=`$ $`\xi _2\lambda ^ne^{i(𝐤𝐦)\mathrm{\Delta }x},`$ (A14)
with $`𝐤`$ the “wave vector”. Notice that the highest wave number that can be represented on the finite difference grid corresponds to $`k_i\mathrm{\Delta }x=\pi `$. The stability condition for our numerical scheme will then be
$$|\lambda |1.$$
(A15)
Let us consider first the “1-step” ICN scheme, that is, the so-called forward-time centered-space (FTCS) scheme. This scheme is well known to be only first order accurate, and unconditionally unstable. The fact that is only first order accurate can be easily seen from a simple Taylor expansion in time. For the stability analysis we substitute the von Neumann ansatz (A14) into the ICN scheme defined above with $`i_{\mathrm{max}}=1`$. Doing this we obtain
$$\lambda ^22\lambda +1+2\rho ^2u^2=0,$$
(A16)
where $`\rho :=\mathrm{\Delta }t/\mathrm{\Delta }x`$ is the Courant parameter and
$`u^2`$ $`:=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}u_i^2,`$ (A17)
$`u_i^2`$ $`:=`$ $`1\mathrm{cos}(k_i\mathrm{\Delta }x).`$ (A18)
Solving for $`\lambda `$ we find
$$\lambda =1\pm i\sqrt{2}\rho u,$$
(A19)
which implies
$$|\lambda |=1+2\rho ^2u^2>1.$$
(A20)
Comparing with (A15) we conclude that the 1-step scheme is unstable for any value of $`\mathrm{\Delta }t`$.
Let us now consider the 2-step scheme. If we take the ICN scheme above with $`i_{\mathrm{max}}=2`$, and do the appropriate substitutions we find
$`\varphi _𝐦^{n+1}`$ $`=`$ $`\varphi _𝐦^n+\mathrm{\Delta }tA_𝐦^n+{\displaystyle \frac{\rho ^2}{2}}{\displaystyle \underset{i=1}{\overset{N}{}}}\delta _i^2\varphi _𝐦^n,`$ (A21)
$`A_𝐦^{n+1}`$ $`=`$ $`A_𝐦^n+{\displaystyle \frac{\rho ^2}{2}}{\displaystyle \underset{i=1}{\overset{N}{}}}\delta _i^2\left(2\varphi _𝐦^n+\mathrm{\Delta }tA_𝐦^n\right).`$ (A22)
As before, a simple Taylor expansion shows that this approximation is now second order both in time and space.
Using again the ansatz (A14) we find now that
$$\lambda ^2+2\lambda \left(\rho ^2u^21\right)+1+\rho ^4u^4=0.$$
(A23)
Solving again for $`\lambda `$ we obtain
$$\lambda =1\rho ^2u^2\pm i\sqrt{2}\rho u,$$
(A24)
which implies
$$|\lambda |=1+\rho ^4u^4>1.$$
(A25)
Comparing again with (A15) we conclude that the 2-step ICN scheme is also unstable for any value of $`\mathrm{\Delta }t`$. This result is surprising, since a priori one might expect that the 2-step scheme should behave like a predictor-corrector scheme, and should therefore be stable.
Finally, let us consider the 3-step scheme. By taking the ICN scheme above with $`i_{\mathrm{max}}=3`$, and doing the appropriate substitutions we now find
$`\varphi _𝐦^{n+1}`$ $`=`$ $`\varphi _𝐦^n+\mathrm{\Delta }tA_𝐦^n+{\displaystyle \frac{\rho ^2}{4}}{\displaystyle \underset{i=1}{\overset{N}{}}}\delta _i^2\left(2\varphi _𝐦^n+\mathrm{\Delta }tA_𝐦^n\right),`$ (A26)
$`A_𝐦^{n+1}`$ $`=`$ $`A_𝐦^n+{\displaystyle \frac{\rho ^2}{2}}{\displaystyle \underset{i=1}{\overset{N}{}}}\delta _i^2\left(2\varphi _𝐦^n+\mathrm{\Delta }tA_𝐦^n\right)`$ (A28)
$`+{\displaystyle \frac{\rho ^3}{4\mathrm{\Delta }x}}\left({\displaystyle \underset{i=1}{\overset{N}{}}}\delta _i^2\right)^2\varphi _𝐦^n.`$
A Taylor expansion now shows that this 3-step scheme is still only second order accurate in both time and space. Using the ansatz (A14) on this scheme we now find
$$\lambda ^2+2\lambda \left(\rho ^2u^21\right)+1\rho ^4u^4+\frac{1}{2}\rho ^6u^6=0.$$
(A29)
And solving for $`\lambda `$ we obtain
$$\lambda =1\rho ^2u^2\pm i\sqrt{2}\rho u\left|1\rho ^2u^2/2\right|,$$
(A30)
which now implies
$$|\lambda |=1\rho ^4u^4+\frac{1}{2}\rho ^6u^6.$$
(A31)
Comparing now with (A15) we obtain the following stability condition
$$\rho ^2u^22.$$
(A32)
And finally, from the fact that the maximum value of $`u^2`$ is $`2\sqrt{N}`$ we find
$$\rho 1/\sqrt{N}.$$
(A33)
Notice that this is just the standard CFL condition in $`N`$ dimensions. We then conclude that in order to obtain a (conditionally) stable scheme we need to do at least three iterations.
Next, we address the question of the stability of the iterations themselves, that is, if we iterate an infinite number of times do we converge to something (that is, to the implicit CN scheme)? For this we consider two consecutive iteration steps $`(i1,i)`$, and subtract them to get
$`\varphi _𝐦^{(i)}\varphi _𝐦^{(i1)}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }t}{2}}\left(A_𝐦^{(i1)}A_𝐦^{(i2)}\right),`$ (A34)
$`A_𝐦^{(i)}A_𝐦^{(i1)}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }t}{2(\mathrm{\Delta }x)^2}}{\displaystyle \underset{i=1}{\overset{N}{}}}\left(\varphi _𝐦^{(i1)}\varphi _𝐦^{(i2)}\right).`$ (A35)
Let us now define $`F_{1}^{}{}_{𝐦}{}^{(i)}:=\varphi _𝐦^{(i)}\varphi _𝐦^{(i1)}`$ and $`F_{2}^{}{}_{𝐦}{}^{(i)}:=A_𝐦^{(i)}A_𝐦^{(i1)}`$. The above equations become
$`F_{1}^{}{}_{𝐦}{}^{(i)}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }t}{2}}F_{2}^{}{}_{𝐦}{}^{(i1)},`$ (A36)
$`F_{2}^{}{}_{𝐦}{}^{(i)}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }t}{2(\mathrm{\Delta }x)^2}}{\displaystyle \underset{i=1}{\overset{N}{}}}F_{1}^{}{}_{𝐦}{}^{(i1)}.`$ (A37)
We now use the von Neumann ansatz again
$`F_{1}^{}{}_{𝐦}{}^{(i)}`$ $`=`$ $`f_1\lambda ^ie^{i(𝐤𝐦)\mathrm{\Delta }x},`$ (A38)
$`F_{2}^{}{}_{𝐦}{}^{(i)}`$ $`=`$ $`f_2\lambda ^ie^{i(𝐤𝐦)\mathrm{\Delta }x},`$ (A39)
Substituting this ansatz back into the equations above we find
$$\lambda ^2+\frac{1}{2}\rho ^2u^2=0,$$
(A40)
from which we obtain
$$\lambda =\pm i\frac{\rho u}{\sqrt{2}}.$$
(A41)
In this case, the condition for the iterations to converge implies that the norm of the successive differences should go to zero, which in turn implies $`|\lambda |<1`$. Using again the fact that the maximum value of $`u^2`$ is $`2\sqrt{N}`$ we see that the convergence condition reduces to
$$\rho <1/\sqrt{N}.$$
(A42)
This is again the standard CFL stability condition. So we have just shown that if this condition is violated, the iterations will fail to converge. This means that there is no reason to try to iterate to convergence in the hope of improving stability. If $`\mathrm{\Delta }t`$ was too big in the first place the iterations will never converge.
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# Gravitational potential energy of simple bodies: the homogeneous bispherical concavo-convex lens
##
### Introduction
At present, the potential energy is known analytically for three types of homogeneous self-gravitating bodies: a) ellipsoids , b) concave bispherical lenses , and 3) rectangular parallelepipeds . In accompaning paper , we present a method of negative density which allows to obtain the analytical solutions for potential energy of new kinds of ”simple” self-gravitating bodies. Here we present the potential energy for the homogeneous bispherical concavo-convex lens with radii of curvature $`R_2=R_1`$. Such figure is obtained by cutting out a symmetric bispheric lens from the sphere (Fig. 1). If we introduce $`H`$, a central thickness of lens, then potential energy of the ”quasi-symmetric” bispheric lens is:
$$\begin{array}{c}W=\frac{\pi }{540H}G\rho ^2R^5[40aH(1216H^2+H^4)+\hfill \\ 3\pi (160240H^2+30H^4H^6)+960(2+3H^2)\mathrm{arctan}\frac{2H}{2a}];\hfill \end{array}$$
(1)
here $`G`$ is constant of Newtonian gravitation, $`\rho =constant`$ is a matter density, $`R`$ is radius of spherical surfaces, $`H`$ is central thickness (in units of $`R`$), and $`a=\sqrt{1H^2/4}`$. (Note that in the caption to Fig. 1 all values are dimensional.) We call the figure in question ”quasi-symmetrical”, as in general case, two spherical surfaces may have the different absolute values of curvature radii. This more general case will be considered elsewhere.
### Dependence of potential energy of the lens on central thickness
The series expansions of $`W`$ at points $`H=0`$ and $`H=2`$ are:
$$\begin{array}{c}W_s(0)/(G\rho ^2R^5)=\frac{64H^2\pi }{27}+\frac{8H^4\pi }{45}+\frac{H^3\pi ^2}{6};\hfill \\ W_s(2)/(G\rho ^2R^5)=\frac{16\pi ^2}{15}+\frac{2\left(2+H\right)^2\pi ^2}{3}.\hfill \end{array}$$
(2)
In general, the function $`W(H)`$ is monotonic (Fig. 2), simply because the volume and mass (and so potential energy) of the lens are all increasing with increasing $`H`$.
### Compactness factor
It is of some interest to consider ”compactness” of the lens as function of $`H`$. We introduce compactness factor, as characteristic of the potential energy of the honogeneous bodies, as the dimensionless coefficient :
$$w=\frac{W}{G\rho ^2V^{5/3}};$$
(3)
in general case of inhomogeneous bodies, it is better to use mean density in the definition of $`w`$. As an example, for the honogeneous sphere we have:
$$w_{sphere}=\frac{3}{5}(\frac{4\pi }{3})^{1/3}=.967195.$$
(4)
This is the maximal possible value of compactness factor for homogeneous bodies.
The volume of the lens is:
$$V=\left(H\frac{H^3}{12}\right)\pi R^3.$$
(5)
The dependence of compactness factor $`w`$ for quasi-symmetric concavo-convex lens on central thickness $`H`$ is shown in Fig. 3. The serial expansion of $`w`$ at point $`H=0`$ is:
$$w_s(0)=\frac{H^{\frac{1}{3}}\left(1289H\pi \right)}{54\pi ^{\frac{2}{3}}},$$
(6)
which is also shown in Fig. 3.
### Conclusion
In conclusion, we present here the analytic formula for gravitational energy of the homogeneous bispheric concavo-convex lens, when the radii of curvature of both surfaces of the lens have the same absolute value. Both potential energy and compactness factor are shown to be monotonic increasing functions of relative central thickness of the lens. As to applicational aspect of the concavo-convex lens problem, we may mention the modelling of the solar system’s small bodies with large craters, such as Phobos with his large Stickney crater.
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# 1 Introduction
## 1 Introduction
The study of the process $`\mathrm{e}^+\mathrm{e}^{}\mathrm{Z}\mathrm{Z}`$ has recently become possible since LEP now operates at center-of-mass energies above the threshold for on-shell $`\mathrm{Z}`$ boson pair production. In the Standard Model, the process $`\mathrm{e}^+\mathrm{e}^{}\mathrm{Z}\mathrm{Z}`$ occurs via the NC2 diagrams shown in Figure 1. The $`\mathrm{Z}`$-pair cross section depends on properties of the $`\mathrm{Z}`$ boson ($`m_\mathrm{Z}`$, $`\mathrm{\Gamma }_\mathrm{Z}`$ and the vector and axial vector coupling of the $`\mathrm{Z}`$ to electrons, $`g_\mathrm{V}^\mathrm{e}`$ and $`g_\mathrm{A}^\mathrm{e}`$) that have been measured with great precision at the $`\mathrm{Z}`$ resonance . The expected $`\mathrm{Z}`$-pair cross section increases from about 0.25 pb at $`\sqrt{s}=183`$ GeV to about 1.0 pb at $`\sqrt{s}=200`$ GeV, but remains more than an order of magnitude smaller than $`\mathrm{W}`$-pair production. In contrast to $`\mathrm{W}`$-pair production, where tree level $`\mathrm{WW}\gamma `$ and $`\mathrm{WW}\mathrm{Z}`$ couplings are important, no $`\mathrm{Z}`$$`\mathrm{Z}`$$`\mathrm{Z}`$ and $`\mathrm{Z}\mathrm{Z}\gamma `$ couplings are expected in the Standard Model. However, physics beyond the Standard Model could lead to effective couplings which could then be observed as deviations in the measured $`\mathrm{Z}`$-pair cross section from the Standard Model prediction. Such deviations have been proposed in the context of Higgs doublet models and in low scale gravity theories . In this paper we report on measurements of the NC2 $`\mathrm{Z}`$-pair cross section, including the extrapolation to final states with one or both $`\mathrm{Z}`$ bosons off-shell. These measurements, along with the angular distribution of the observed events, are then used to extract limits on possible $`\mathrm{Z}`$$`\mathrm{Z}`$$`\mathrm{Z}`$ and $`\mathrm{Z}\mathrm{Z}\gamma `$ couplings.
In Section 2 we describe the data sets used and the Monte Carlo simulation of signal and background. In Section 3 we describe the selection of the processes $`\mathrm{Z}\mathrm{Z}\mathrm{}^+\mathrm{}^{}\mathrm{}^+\mathrm{}^{}`$, $`\mathrm{Z}\mathrm{Z}\mathrm{}^+\mathrm{}^{}\nu \overline{\nu }`$, $`\mathrm{Z}\mathrm{Z}\mathrm{q}\overline{\mathrm{q}}\mathrm{}^+\mathrm{}^{}`$, $`\mathrm{Z}\mathrm{Z}\mathrm{q}\overline{\mathrm{q}}\nu \overline{\nu }`$, and $`\mathrm{Z}\mathrm{Z}\mathrm{q}\overline{\mathrm{q}}\mathrm{q}\overline{\mathrm{q}}`$, where $`\mathrm{}^+\mathrm{}^{}`$ denotes a charged lepton pair of opposite charge and $`\mathrm{q}\overline{\mathrm{q}}`$ any of the five lightest quark-antiquark pairs. We also describe analyses of $`\mathrm{Z}\mathrm{Z}\mathrm{b}\overline{\mathrm{b}}\mathrm{}^+\mathrm{}^{}`$, $`\mathrm{Z}\mathrm{Z}\mathrm{b}\overline{\mathrm{b}}\nu \overline{\nu }`$ and $`\mathrm{Z}\mathrm{Z}\mathrm{q}\overline{\mathrm{q}}\mathrm{b}\overline{\mathrm{b}}`$ which use b-tagging methods similar to those used in the OPAL Higgs search . The use of b-tagging improves the separation of the $`\mathrm{Z}`$-pair signal from background and allows us to check the $`\mathrm{b}\overline{\mathrm{b}}`$ content of our $`\mathrm{Z}`$-pair sample for consistency with the Standard Model. The description of the individual selections is followed by a discussion of possible systematic errors (Section 3.6). In Section 4 the selected events are used to measure the $`\mathrm{Z}`$-pair cross section. Then the cross section and angular distribution are compared with the Standard Model predictions and limits on anomalous neutral current triple gauge couplings are derived.
## 2 Data analysis and Monte Carlo
The OPAL detector<sup>1</sup><sup>1</sup>1OPAL uses a right-handed coordinate system in which the $`z`$ axis is along the electron beam direction and the $`x`$ axis is horizontal. The polar angle, $`\theta `$, is measured with respect to the $`z`$ axis and the azimuthal angle, $`\varphi `$, with respect to the $`x`$ axis., trigger and data acquisition system are described fully elsewhere . Our analyses use approximately $`55\mathrm{pb}^1`$ of data collected at center-of-mass energies between 181–184 GeV and approximately $`178\mathrm{pb}^1`$ collected at center-of-mass energies near 189 GeV. The corresponding luminosity-weighted mean center-of-mass energies are $`182.62\pm 0.05`$ GeV and $`188.63\pm 0.04`$ GeV . The luminosity was measured using small-angle Bhabha scattering events recorded in the silicon-tungsten luminometer and the theoretical calculation of Reference . The overall error on the luminosity measurement amounts to less than 0.5% and contributes negligibly to our cross-section measurement error.
Selection efficiencies and backgrounds were calculated using Monte Carlo simulations. All events were passed through a simulation of the OPAL detector and processed as for data. We define the $`\mathrm{Z}`$$`\mathrm{Z}`$ cross section as the contribution to the total four-fermion cross section from the NC2 $`\mathrm{Z}`$-pair diagrams shown in Figure 1. All signal efficiencies given in this paper are with respect to these $`\mathrm{Z}`$-pair processes. Contributions from all other four-fermion final states, including interference with NC2 diagrams, are considered as background. For studies of the signal efficiency we have used grc4f , YFSZZ and PYTHIA .
Backgrounds are simulated using several different generators. PYTHIA is used to simulate two-fermion final states such as $`\mathrm{e}^+\mathrm{e}^{}\mathrm{Z}^{}(n\gamma )\mathrm{q}\overline{\mathrm{q}}(n\gamma )`$ and $`\mathrm{e}^+\mathrm{e}^{}\gamma ^{}(n\gamma )\mathrm{q}\overline{\mathrm{q}}(n\gamma )`$, where $`(n\gamma )`$ indicates the generation of one or more initial state photons. HERWIG and KK2f are used as checks for these final states. These two-fermion generators include gluon radiation from the quarks which produce $`\mathrm{q}\overline{\mathrm{q}}\mathrm{g}`$, $`\mathrm{q}\overline{\mathrm{q}}\mathrm{q}\overline{\mathrm{q}}`$ and $`\mathrm{q}\overline{\mathrm{q}}\mathrm{gg}`$ final states. The grc4f generator, with the contribution exclusively due to NC2 diagrams removed, is used to simulate other four-fermion background. KORALW and EXCALIBUR are used as checks of the four-fermion background. Multiperipheral (“two-photon”) processes with hadronic final states are simulated by combining events from PYTHIA, for events without electrons scattered into the detector, and HERWIG , for events with electrons scattered into the detector. For the $`\mathrm{q}\overline{\mathrm{q}}\mathrm{e}^+\mathrm{e}^{}`$ final state, TWOGEN is used to simulate two-photon events with both the electron and positron scattered into the detector. The Vermaseren generator is used to simulate multiperipheral production of the final states $`\mathrm{e}^+\mathrm{e}^{}\mathrm{}^+\mathrm{}^{}`$.
To avoid background from four-fermion final states mediated by $`\mathrm{Z}\gamma ^{}`$, our selections were optimized to select simulated events with masses, $`m_1`$ and $`m_2`$, that satisfy $`m_1+m_2>170\mathrm{GeV}`$ and $`|m_1m_2|<20\mathrm{GeV}`$. At 189 GeV (183 GeV) more than 90% (80%) of the events produced via the NC2 diagrams are contained in this mass region. Events from the NC2 diagrams dominate in this mass region except for final states containing electron pairs. Backgrounds in these samples from two-photon and electroweak Compton scattering ($`\mathrm{e}\gamma \mathrm{e}\mathrm{Z}`$) processes are reduced by using electrons detected in the electromagnetic calorimeters with $`|\mathrm{cos}\theta _\mathrm{e}|<0.985`$, where $`\theta _\mathrm{e}`$ is the polar angle of the electron.
## 3 Event selection
In the following subsections we describe event selections which exploit every decay mode of the $`\mathrm{Z}`$ boson. Our selections cover all $`\mathrm{Z}`$$`\mathrm{Z}`$ final states except $`\nu \overline{\nu }\nu \overline{\nu }`$ and $`\tau ^+\tau ^{}\nu \overline{\nu }`$. In hadronic final states, the energy and direction of the jets are determined using reconstructed tracks and calorimeter clusters using the correction for double counting described in Reference . In the $`\mathrm{q}\overline{\mathrm{q}}\mathrm{}^+\mathrm{}^{}`$, $`\mathrm{q}\overline{\mathrm{q}}\mathrm{q}\overline{\mathrm{q}}`$ and $`\mathrm{q}\overline{\mathrm{q}}\mathrm{b}\overline{\mathrm{b}}`$ analyses four-constraint (4C) and five-constraint (5C) kinematic fits are used. The 4C fit imposes energy and momentum conservation. In the 5C fit the added constraint requires the masses of the two candidate $`\mathrm{Z}`$ bosons to be equal to one another. For final states with one or more $`\mathrm{Z}`$ bosons decaying to tau pairs, the energy and total momentum of the tau leptons are obtained by leaving the reconstructed direction of the four fermions fixed and scaling the energy and momentum of each of the fermions to obtain energy and momentum conservation. The scaled values of the tau momentum and energy are then used in the subsequent steps of the analysis. In the $`\mathrm{q}\overline{\mathrm{q}}\tau ^+\tau ^{}`$ and $`\mathrm{b}\overline{\mathrm{b}}\tau ^+\tau ^{}`$ final states, subsequent kinematic fits are effectively 2C and 3C fits.
### 3.1 Selection of $`𝐙𝐙\mathbf{}\mathbf{}^\mathbf{+}\mathbf{}^{\mathbf{}}\mathbf{}^\mathbf{+}\mathbf{}^{\mathbf{}}`$ events
$`\mathrm{Z}`$-pair events decaying to final states with four charged leptons ($`\mathrm{}^+\mathrm{}^{}\mathrm{}^+\mathrm{}^{}`$) produce low multiplicity events with a clear topological signature that is exploited to maximize the selection efficiency. The $`\mathrm{}^+\mathrm{}^{}\mathrm{}^+\mathrm{}^{}`$ analysis begins by selecting low multiplicity events (less than 13 tracks or clusters) with visible energy of at least $`0.2\sqrt{s}`$ and at least one track with momentum of 5 GeV or more. Using a cone algorithm, the events are required to have exactly four cones of $`15^{}`$ half angle containing between 1 and 3 tracks. Cones of opposite charge are paired<sup>2</sup><sup>2</sup>2Two-track cones are assigned the charge of the most energetic track if the momentum of one track exceeds that of the other by a factor of 4. Events with a cone which fails this requirement are rejected. to form $`\mathrm{Z}`$ boson candidates.
Lepton identification is only used to classify events as background or to reduce the number of cone combinations considered by preventing the matching of identified electrons with identified muons. Electrons are identified on the basis of energy deposition in the electromagnetic calorimeter, track curvature and specific ionization in the tracking chambers. Muons are identified using the association between tracks and hits in the hadron calorimeter and muon chambers.
To reduce background from two-photon events with a single scattered electron detected, we eliminate events with forward going electrons (backward going positrons) with the cut $`\mathrm{cos}\theta _\mathrm{e}^{}<0.85`$ ($`\mathrm{cos}\theta _{\mathrm{e}^+}>0.85`$). Here $`\theta _\mathrm{e}^{}`$ ( $`\theta _{\mathrm{e}^+}`$) is the angle of the electron (positron) with respect to the incoming electron beam. Background from partially reconstructed $`\mathrm{q}\overline{\mathrm{q}}(n\gamma )`$ events and two-photon events is reduced by requiring that most of the energy is not concentrated in a single cone, $`E_{\mathrm{vis}}E_{\mathrm{cone}}^{\mathrm{max}}>0.2\sqrt{s}`$. Here $`E_{\mathrm{vis}}`$ is the total visible energy of the event and $`E_{\mathrm{cone}}^{\mathrm{max}}`$ is the energy contained in the most energetic cone.
The invariant masses of the lepton pairs are calculated in three different ways which are motivated by the possibility of having zero, one or two tau pairs in the event. The events are classified according to the number of tau pairs in the event. (i) Events with $`E_{\mathrm{vis}}>0.9\sqrt{s}`$ are treated as $`\mathrm{e}^+\mathrm{e}^{}\mathrm{e}^+\mathrm{e}^{}`$, $`\mathrm{e}^+\mathrm{e}^{}\mu ^+\mu ^{}`$ or $`\mu ^+\mu ^{}\mu ^+\mu ^{}`$ events. We also treat all events with $`|\mathrm{cos}\theta _{\mathrm{miss}}|>0.98`$ ($`\theta _{\mathrm{miss}}`$ is the polar angle associated with the missing momentum in the event) as $`\mathrm{e}^+\mathrm{e}^{}\mathrm{e}^+\mathrm{e}^{}`$, $`\mathrm{e}^+\mathrm{e}^{}\mu ^+\mu ^{}`$ or $`\mu ^+\mu ^{}\mu ^+\mu ^{}`$ events to maintain efficiency for $`\mathrm{Z}`$-pairs with initial state radiation. In these events there are no missing neutrinos and the mass of each pair of lepton cones is evaluated. (ii) Events failing (i) with a cone-pair combination that has energy exceeding $`0.9m_\mathrm{Z}`$ are tried as an $`\mathrm{e}^+\mathrm{e}^{}\tau ^+\tau ^{}`$ or $`\mu ^+\mu ^{}\tau ^+\tau ^{}`$ final state. The mass of the tau-pair system is calculated from the recoil mass of the presumed electron or muon pair. (iii) Any remaining combinations are treated as $`\tau ^+\tau ^{}\tau ^+\tau ^{}`$ final states. The momenta of the tau leptons are determined with the scaling procedure described in the introduction to Section 3 and the invariant masses of the cone pairs are evaluated using the scaled momenta. In any event with more than one valid combination, each combination is tested using the invariant mass cuts listed below.
To reduce the combinatorial background, combinations with pair masses closest to $`m_\mathrm{Z}`$ are selected. In events with one or more combination satisfying $`|m_\mathrm{Z}m_{\mathrm{}\mathrm{}}|>0.1m_\mathrm{Z}`$ the cone-pair combination with the smallest value of $`(m_\mathrm{Z}m_{\mathrm{}\mathrm{}})^2+(m_\mathrm{Z}m_{\mathrm{}^{}\mathrm{}^{}})^2`$ is selected for further analysis. In the other combinations, the combination with the smallest value of $`|m_\mathrm{Z}m_{\mathrm{}\mathrm{}}|`$ or $`|m_\mathrm{Z}m_{\mathrm{}^{}\mathrm{}^{}}|`$ is selected. The final event sample is then chosen with the requirement $`m_{\mathrm{}\mathrm{}}+m_{\mathrm{}^{}\mathrm{}^{}}>160`$ GeV and $`|m_{\mathrm{}\mathrm{}}m_{\mathrm{}^{}\mathrm{}^{}}|<40`$ GeV. The signal detection efficiency, averaged over all $`\mathrm{}^+\mathrm{}^{}\mathrm{}^+\mathrm{}^{}`$ final states is given in Table 1 (line a). The efficiency for individual final states range from 30% for $`\tau ^+\tau ^{}\tau ^+\tau ^{}`$ to more than 70% for $`\mu ^+\mu ^{}\mu ^+\mu ^{}`$. The invariant masses of all cone pairs passing one of the selections are shown in Figure 2a. One candidate is found in the 183 GeV data and one candidate is found in the 189 GeV data.
### 3.2 Selection of $`𝐙𝐙\mathbf{}\mathbf{}^\mathbf{+}\mathbf{}^{\mathbf{}}𝝂\overline{𝝂}`$ events
The selection of the $`\mathrm{e}^+\mathrm{e}^{}\nu \overline{\nu }`$ and $`\mu ^+\mu ^{}\nu \overline{\nu }`$ final states is based on the OPAL selection of $`\mathrm{W}`$ pairs decaying to leptons . The mass and momentum of the $`\mathrm{Z}`$ boson decaying to $`\nu \overline{\nu }`$ are calculated using the beam energy constraint and the visible decay of the other $`\mathrm{Z}`$ boson to a charged lepton pair. A likelihood selection based on the visible and recoil masses as well as the polar angle of the leptons, is then used to separate signal from background.
The $`\mathrm{e}^+\mathrm{e}^{}\nu \overline{\nu }`$ selection starts with OPAL $`\mathrm{W}`$-pair candidates where both charged leptons are classified as electrons. Each event is then divided into two hemispheres using the thrust axis. The highest momentum charged (leading) track is selected from each hemisphere. The sum of the charges of these two tracks is required to be zero. The determination of the visible mass, $`m_{\mathrm{vis}}`$, and the recoil mass, $`m_{\mathrm{recoil}}`$, is based on the energy as measured in the electromagnetic calorimeter and the direction of the leading tracks.
Three variables were chosen for the likelihood selection: $`Q\mathrm{cos}\theta `$, where $`\theta `$ is the angle of the highest momentum charged track and $`Q`$ is its charge, the normalized sum of visible and recoil masses $`(m_{\mathrm{vis}}+m_{\mathrm{recoil}})/\sqrt{s}`$ and the difference of visible and recoil masses, $`m_{\mathrm{vis}}m_{\mathrm{recoil}}`$. The performance of the likelihood is improved with the following preselection: $`25\mathrm{GeV}<\mathrm{m}_{\mathrm{vis}}\mathrm{m}_{\mathrm{recoil}}<15\mathrm{GeV}`$ and $`(\mathrm{m}_{\mathrm{vis}}+\mathrm{m}_{\mathrm{recoil}})/\sqrt{\mathrm{s}}>0.90`$. One event with $`_{\mathrm{e}^+\mathrm{e}^{}\nu \overline{\nu }}>0.60`$ is selected (see Table 1 (line b) and Figure 3a).
The $`\mu ^+\mu ^{}\nu \overline{\nu }`$ selection starts with the OPAL $`\mathrm{W}`$-pair candidates where both charged leptons are classified as muons. The selection procedure is the same as for the $`\mathrm{e}^+\mathrm{e}^{}\nu \overline{\nu }`$ final states except that $`m_{\mathrm{vis}}`$, $`m_{\mathrm{recoil}}`$, and $`E_{\mathrm{vis}}`$ are calculated from the the momentum of the reconstructed tracks of the $`\mathrm{Z}`$ boson decaying to muon pairs. The likelihood preselections $`25\mathrm{GeV}<\mathrm{m}_{\mathrm{vis}}\mathrm{m}_{\mathrm{recoil}}<25\mathrm{GeV}`$ and $`(\mathrm{m}_{\mathrm{vis}}+\mathrm{m}_{\mathrm{recoil}})/\sqrt{\mathrm{s}}>0.90`$ are applied. Two events with $`_{\mu ^+\mu ^{}\nu \overline{\nu }}>0.60`$ are selected (see Table 1 (line c) and Figure 3b).
### 3.3 Selection of $`𝐙𝐙\mathbf{}𝐪\overline{𝐪}\mathbf{}^\mathbf{+}\mathbf{}^{\mathbf{}}`$ events
The lepton pairs in the $`\mathrm{q}\overline{\mathrm{q}}\mathrm{e}^+\mathrm{e}^{}`$ and $`\mathrm{q}\overline{\mathrm{q}}\mu ^+\mu ^{}`$ final states have a distinctive signature making possible selections with high efficiencies and a low background contamination. In the $`\mathrm{q}\overline{\mathrm{q}}\tau ^+\tau ^{}`$ final state, the decay of the tau leptons produces events which are more difficult to identify. The identification of this final state exploits the missing momentum and missing energy carried away by the neutrinos produced in the decay of the tau lepton.
#### 3.3.1 Selection of $`𝐙𝐙\mathbf{}𝐪\overline{𝐪}𝐞^\mathbf{+}𝐞^{\mathbf{}}`$ and $`𝐙𝐙\mathbf{}𝐪\overline{𝐪}𝝁^\mathbf{+}𝝁^{\mathbf{}}`$ events
The selection of $`\mathrm{q}\overline{\mathrm{q}}\mathrm{e}^+\mathrm{e}^{}`$ and $`\mathrm{q}\overline{\mathrm{q}}\mu ^+\mu ^{}`$ final states requires the visible energy of the events to be greater than 90 GeV and at least six reconstructed tracks. Among all tracks with momenta greater than 2 GeV, the highest momentum track is taken as the first lepton candidate and the second-highest momentum track with a charge opposite to the first candidate is taken as the second lepton candidate. Using the Durham jet algorithm, the event, including the lepton candidates, is forced into four jets and the jet resolution variable that separates the three-jet topology from the four-jet topology, $`y_{34}`$, is required to be greater than $`10^3`$. Excluding the electron or muon candidates and their associated calorimeter clusters, the rest of the event is forced into two jets. The 4C and 5C fits to the two lepton candidates and the two jets are required to converge.<sup>3</sup><sup>3</sup>3 In the context of this paper, convergence is defined as a fit probability greater than $`10^{10}`$.
In the $`\mathrm{q}\overline{\mathrm{q}}\mathrm{e}^+\mathrm{e}^{}`$ selection no explicit electron identification is used. Electron candidates are selected by requiring the sum of the electromagnetic cluster energies $`E_1+E_2`$ associated to the electrons to be greater than 70 GeV and the momentum of the most energetic electron track to exceed 20 GeV. We also reject the event if the angle between either electron candidate and any other track is less than 5.
In the $`\mathrm{q}\overline{\mathrm{q}}\mu ^+\mu ^{}`$ selection the muons are identified using (i) tracks which match a reconstructed segment in the muon chambers, (ii) tracks which are associated to hits in the hadron calorimeter or muon chambers , or (iii) isolated tracks associated to electromagnetic clusters with reconstructed energy less than 2 GeV. No isolation requirement is imposed on events with both muon tracks passing (i) or (ii). Events with at least one muon identified with criterion (iii) are accepted if both muon candidates in the event have an angle of at least 10 to the nearest track. We require the sum of the momenta of the two leptons to be greater than 70 GeV.
$`\mathrm{Z}`$-pair events are separated from $`\mathrm{Z}\gamma ^{}`$ background by requiring the fitted mass of the 5C fit to be larger than 85 GeV and the invariant masses $`m_{\mathrm{}\mathrm{}}`$ and $`m_{\mathrm{qq}}`$ obtained from the 4C fit to satisfy $`(m_{\mathrm{}\mathrm{}}+m_{\mathrm{qq}})>170`$ GeV and $`|m_{\mathrm{}\mathrm{}}m_{\mathrm{qq}}|<30`$ GeV. Figure 2b (2c) shows the distribution of $`m_{\mathrm{ee}}`$ ($`m_{\mu \mu }`$) and $`m_{\mathrm{qq}}`$ before the cuts on the masses from the 4C and 5C fits.
After all cuts the selection efficiency<sup>4</sup><sup>4</sup>4 Small amounts of feedthrough from other $`\mathrm{Z}`$$`\mathrm{Z}`$ final states, in this case $`\mathrm{q}\overline{\mathrm{q}}\tau ^+\tau ^{}`$, are counted as signal. for $`\mathrm{q}\overline{\mathrm{q}}\mathrm{e}^+\mathrm{e}^{}`$ signal events is $`(55.2\pm 2.7)`$% at 183 GeV and $`(65.1\pm 2.9)`$% at 189 GeV. The errors on these efficiencies include the systematic errors (see Section 3.6). No candidate is observed at 183 GeV. Six candidate events are found after all cuts in the data taken at 189 GeV. In Table 1 (lines d and e) we give the efficiency, background and observed number of events. The largest source of background after all cuts is from $`\mathrm{Z}\gamma ^{}`$ mediated $`\mathrm{q}\overline{\mathrm{q}}\mathrm{e}^+\mathrm{e}^{}`$ events and from two-photon events.
For the $`\mathrm{q}\overline{\mathrm{q}}\mu ^+\mu ^{}`$ final state the selection efficiency is $`(66.9\pm 2.6)`$% at 183 GeV and $`(72.3\pm 2.8)`$% at 189 GeV. Three events are observed in the 189 GeV data sample and none in the 183 GeV data sample. The errors on these efficiencies include the systematic errors given below. In Table 1 (lines f and g) we give the efficiency, backgrounds and observed number of events. The background after all cuts is expected to come from $`\mathrm{e}^+\mathrm{e}^{}\mathrm{Z}\gamma ^{}\mathrm{q}\overline{\mathrm{q}}\mu ^+\mu ^{}`$ events.
#### 3.3.2 Selection of $`𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}𝐙𝐙\mathbf{}𝐪\overline{𝐪}𝝉^\mathbf{+}𝝉^{\mathbf{}}`$ events
The $`\mathrm{q}\overline{\mathrm{q}}\tau ^+\tau ^{}`$ final state is selected from a sample of events with track multiplicity greater or equal to six. Events which have been selected as $`\mathrm{e}^+\mathrm{e}^{}\mathrm{Z}\mathrm{Z}\mathrm{q}\overline{\mathrm{q}}\mathrm{e}^+\mathrm{e}^{}`$ or $`\mathrm{e}^+\mathrm{e}^{}\mathrm{Z}\mathrm{Z}\mathrm{q}\overline{\mathrm{q}}\mu ^+\mu ^{}`$ are excluded from this selection. The tau-lepton candidates are selected using an artificial neural network algorithm which is described in detail in Reference . The tau candidate with the highest neural network output value is taken as the first candidate. The second best candidate is required to have its charge opposite to the first candidate and the highest output value among all remaining candidates. If a second candidate cannot be found the event is rejected.
Motivated by the presence of neutrinos in the final state, the visible energy of the event, $`E_{\mathrm{vis}}`$, is required to exceed 90 GeV and the missing energy $`\sqrt{s}E_{\mathrm{vis}}`$ is required to exceed 15 GeV. In addition, the sum of the momenta of the leading tracks from the tau-lepton decays is required to be less than 70 GeV. Since the direction of the missing momentum in signal events will tend to be along the direction of one of the decaying tau leptons, the angle $`\alpha _{\tau ,\mathrm{miss}}`$ between the missing momentum and a tau-lepton candidate is required to satisfy $`\alpha _{\tau ,\mathrm{miss}}<90^{}`$ for at least one of the two tau candidates.
The two hadronic jets are selected in the same way as in the $`\mathrm{e}^+\mathrm{e}^{}\mathrm{Z}\mathrm{Z}\mathrm{q}\overline{\mathrm{q}}\mathrm{e}^+\mathrm{e}^{}`$ selection. The initial estimate of the energy and the momenta of the tau candidates is found from the sum of the tracks associated to the tau by the neural network algorithm and all unassociated electromagnetic clusters in a cone with a half angle of 10 around the leading track from the tau decay. A 2C kinematic fit that imposes energy and momentum conservation (see the introduction to Section 3) is required to converge. A 3C kinematic fit, with the additional constraint of the equality of the fermion pair masses is also required to converge.
Using the network output for each tau lepton, a probability is calculated taking into account the different branching ratios, sensitivities, efficiencies and background levels for 1-prong and 3-prong tau-lepton decays. In the following, we combine the probabilities $`𝒫_1`$ and $`𝒫_2`$ to form a likelihood using
$$=\frac{𝒫_1𝒫_2}{𝒫_1𝒫_2+(1𝒫_1)(1𝒫_2)}.$$
(1)
The likelihood associated with probabilities of the two tau candidates is required to satisfy $`_{\tau \tau }>0.977`$. In addition, the common mass of the 3C fit is required to exceed 85 GeV. Using the 2C fit masses of the tau pair, $`m_{\tau \tau }`$, and the quark pair, $`m_{\mathrm{qq}}`$, as obtained from the kinematic fit, we also require $`m_{\mathrm{qq}}+m_{\tau \tau }>170`$ GeV and $`|m_{\mathrm{qq}}m_{\tau \tau }|<60`$ GeV.
After all cuts the selection efficiency for signal events is found to be $`(22.2\pm 1.7)`$% at 183 GeV and $`(26.9\pm 2.0)`$% at 189 GeV. One candidate event is found in the data at 189 GeV while no candidate is selected at 183 GeV. Figure 2d shows the masses of the candidate events before the invariant mass cuts. In Table 1 (lines h and i) we give the efficiencies, backgrounds and observed number of events.
#### 3.3.3 Selection of $`𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}𝐙𝐙\mathbf{}𝐛\overline{𝐛}\mathbf{}^\mathbf{+}\mathbf{}^{\mathbf{}}`$
Events with $`\mathrm{b}\overline{\mathrm{b}}`$ final states are selected using the algorithm described in Reference . The b probabilities of the two hadronic jets are combined to form a likelihood, $`_{\mathrm{bb}}`$, according to Equation 1. Because the $`\mathrm{q}\overline{\mathrm{q}}\mathrm{e}^+\mathrm{e}^{}`$ and $`\mathrm{q}\overline{\mathrm{q}}\mu ^+\mu ^{}`$ selections are pure, a relatively loose cut of $`_{\mathrm{bb}}>0.2`$ is used to select the $`\mathrm{b}\overline{\mathrm{b}}\mathrm{e}^+\mathrm{e}^{}`$ and $`\mathrm{b}\overline{\mathrm{b}}\mu ^+\mu ^{}`$ samples. For the selections with electron and muon pairs there are two classes of events since the selected $`\mathrm{b}\overline{\mathrm{b}}\mathrm{}^+\mathrm{}^{}`$ events are a subset of the $`\mathrm{q}\overline{\mathrm{q}}\mathrm{}^+\mathrm{}^{}`$ events. In Table 1 (lines e and g) we give the efficiencies of the b-tagged samples with respect to the expected fraction of $`\mathrm{b}\overline{\mathrm{b}}`$ events. The efficiencies for samples without b-tags are given with respect to the hadronic decays without $`\mathrm{b}\overline{\mathrm{b}}`$ final states.
In the $`\mathrm{b}\overline{\mathrm{b}}\mathrm{e}^+\mathrm{e}^{}`$ selection one candidate is found in the data. The selection efficiency is found to be (47$`\pm `$3)% at 183 GeV and (46$`\pm `$3)% at 189 GeV. <sup>5</sup><sup>5</sup>5 The efficiencies given in this section do not include feedthrough from other $`\mathrm{q}\overline{\mathrm{q}}\mathrm{}^+\mathrm{}^{}`$ final states. The efficiencies given in Table 1 include this feedthrough. In the $`\mathrm{b}\overline{\mathrm{b}}\mu ^+\mu ^{}`$ selection we find no candidate at 183 GeV and one candidate at 189 GeV. The selection efficiency is found to be (50$`\pm `$3)% at 183 GeV and (54$`\pm `$3)% at 189 GeV.
For the $`\mathrm{b}\overline{\mathrm{b}}\tau ^+\tau ^{}`$ selection the $`_{\tau \tau }`$ cut of the $`\mathrm{q}\overline{\mathrm{q}}\tau ^+\tau ^{}`$ selection is loosened and combined with $`_{\mathrm{bb}}`$ as follows. $`_{\tau \tau }`$ and $`_{\mathrm{bb}}`$ are both required to be greater than 0.1. The $`\mathrm{b}\overline{\mathrm{b}}\tau ^+\tau ^{}`$ probability for the event, $`_{\mathrm{bb}\tau \tau }`$, is calculated from Equation 1 with $`_{\tau \tau }`$ and $`_{\mathrm{bb}}`$ as inputs and required to exceed 0.95. After the cut on $`_{\mathrm{bb}\tau \tau }`$, the remaining cuts of the $`\mathrm{q}\overline{\mathrm{q}}\tau ^+\tau ^{}`$ selection are applied, giving a selection efficiency of $`(21\pm 3)`$% at 183 GeV and $`(24\pm 3)`$% at 189 GeV. No candidate event is found in the 183 GeV or 189 GeV data.
We also use an alternative jet-based $`\mathrm{b}\overline{\mathrm{b}}\tau ^+\tau ^{}`$ selection and accept any event which passes either $`\mathrm{b}\overline{\mathrm{b}}\tau ^+\tau ^{}`$ analysis, but events previously selected by another $`\mathrm{q}\overline{\mathrm{q}}\mathrm{}^+\mathrm{}^{}`$ selection are rejected. The alternative analysis uses a different approach to reconstruct the tau leptons. This event selection consists of a set of preselection cuts and a subsequent multivariate likelihood selection. Events are reconstructed as four jets using the Durham algorithm. Tau-lepton candidates are sought in the four jets using a likelihood technique to separate real tau leptons and fakes in quark jets. The tau and b-tag likelihood values are combined in a $`\mathrm{b}\overline{\mathrm{b}}\tau ^+\tau ^{}`$ likelihood which is maximized to choose the b jets and tau leptons of the event. This $`\mathrm{b}\overline{\mathrm{b}}\tau ^+\tau ^{}`$ likelihood uses tau and b-tag likelihoods and some topological variables as input. Events are accepted if their likelihood exceeds 0.6. In addition, the fitted 2C masses are required to satisfy $`m_{\mathrm{qq}}+m_{\tau \tau }>170`$ GeV and $`|m_{\mathrm{qq}}m_{\tau \tau }|<60`$ GeV.
At $`189\mathrm{GeV}`$ data the alternative selection has an efficiency of $`(30\pm 3)`$% for $`\mathrm{b}\overline{\mathrm{b}}\tau ^+\tau ^{}`$ events. After combining the two selections, the efficiency for $`\mathrm{b}\overline{\mathrm{b}}\tau ^+\tau ^{}`$ events for all cuts except those rejecting events found by the other $`\mathrm{q}\overline{\mathrm{q}}\mathrm{}^+\mathrm{}^{}`$ selections is 40%. A similar improvement is realized for the $`183\mathrm{GeV}`$ data. The efficiency and backgrounds after rejecting events found by the other selections are given in Table 1 (line j). One exclusive event is selected by the combined $`\mathrm{b}\overline{\mathrm{b}}\tau ^+\tau ^{}`$ analysis at $`189\mathrm{GeV}`$.
### 3.4 Selection of $`𝐙𝐙\mathbf{}𝐪\overline{𝐪}𝝂\overline{𝝂}`$ events
The $`\mathrm{q}\overline{\mathrm{q}}\nu \overline{\nu }`$ selection is based on the reconstruction of the $`\mathrm{Z}`$ boson decaying to $`\mathrm{q}\overline{\mathrm{q}}`$ which produces somewhat back-to-back jets. The selection uses contained events with a two-jet topology. The beam energy constraint is then used to determine the mass of the $`\mathrm{Z}`$ boson decaying to $`\nu \overline{\nu }`$. The properties of the $`\mathrm{q}\overline{\mathrm{q}}`$ decay and the inferred mass of the $`\nu \overline{\nu }`$ decay are then used in a likelihood analysis to separate signal from background.
Two-jet events are selected by dividing each event into two hemispheres using the plane perpendicular to the thrust axis. The number of charged tracks in each hemisphere is required to be four or more. The polar angles of the energy-momentum vector associated with each hemisphere, $`\theta _{\mathrm{hemi1}}`$ and $`\theta _{\mathrm{hemi2}}`$, are used to calculate the quantity $`\mathrm{cos}\theta _\mathrm{h}=\frac{1}{2}(\mathrm{cos}\theta _{\mathrm{hemi1}}\mathrm{cos}\theta _{\mathrm{hemi2}})`$. Contained events are selected by requiring $`|\mathrm{cos}\theta _\mathrm{h}|<0.80`$. The total energy in the forward detectors and in the forward region of the electromagnetic calorimeter ($`|\mathrm{cos}\theta |>0.95`$) is required to be less than 3 GeV. $`\mathrm{W}`$ decays identified by the OPAL $`\mathrm{W}`$-pair selection are rejected; the likelihood for $`\mathrm{e}^+\mathrm{e}^{}\mathrm{q}\overline{\mathrm{q}}\mathrm{}\nu `$ from Reference , $`_{\mathrm{WW}}`$, is required to be 0.5 or less.
An important background to our selection is $`\mathrm{q}\overline{\mathrm{q}}(n\gamma )`$ events with photons which escape detection. We discriminate against these events by looking for a significant amount of missing transverse momentum, $`p_\mathrm{t}`$. In each event, $`p_\mathrm{t}`$ can be resolved into two components, $`p_{\mathrm{t}i}`$, perpendicular to both the thrust axis and the beam axis and $`p_{\mathrm{t}j}`$, along the thrust axis and perpendicular to the beam axis. Since $`p_{\mathrm{t}i}`$ is based primarily on angular measurements, it is better measured than $`p_{\mathrm{t}j}`$. We approximate $`p_{\mathrm{t}i}`$ as $`p_{\mathrm{t}i}=\frac{1}{2}E_\mathrm{b}\mathrm{sin}\varphi \mathrm{sin}\theta _\mathrm{h}`$. Here $`E_\mathrm{b}=\sqrt{s}/2`$ is the beam energy, $`\varphi `$ is the acoplanarity of the momentum vectors of the two hemispheres and $`\mathrm{sin}\theta _\mathrm{h}=\sqrt{1\mathrm{cos}^2\theta _\mathrm{h}}`$. The resolution on $`p_{\mathrm{t}i}`$, $`\sigma _{p_{\mathrm{t}i}}`$, was parameterized as a function of thrust and $`\mathrm{cos}\theta _\mathrm{h}`$ using data taken at the $`\mathrm{Z}`$ resonance. The variable $`R_{p_{\mathrm{t}i}}=(p_{\mathrm{t}i}p_{\mathrm{t}i}^0)/\sigma _{p_{\mathrm{t}i}}`$ is used as an input to the likelihood. Here $`p_{\mathrm{t}i}^0`$ corresponds to the transverse momentum carried by a photon with half the beam energy which just misses the inner edge of our acceptance ( $`p_{\mathrm{t}i}^0=E_\mathrm{b}\mathrm{sin}(32\text{mrad})/2`$). We also use the variable $`\mathrm{cos}\theta _{\mathrm{miss}}`$, the direction of the missing momentum in the event, to discriminate against the $`\mathrm{q}\overline{\mathrm{q}}(n\gamma )`$ events.
In the final selection of events, we use a likelihood based on the following five variables: (i) the normalized sum of visible and recoil masses $`(m_{\mathrm{vis}}+m_{\mathrm{recoil}})/\sqrt{s}`$, (ii) the difference of visible and recoil masses ($`m_{\mathrm{vis}}m_{\mathrm{recoil}}`$), (iii) $`\mathrm{log}(y_{23})`$, where $`y_{23}`$ is the jet resolution parameter that separates the two-jet topology from the three-jet topology as calculated from the Durham jet algorithm, (iv) $`\mathrm{cos}\theta _{\mathrm{miss}}`$ and (v) $`R_{p_{\mathrm{t}i}}`$. The mass variables are useful for reducing background from $`\mathrm{W}`$-pair production and $`\mathrm{We}\nu `$ final states. The jet resolution parameter is useful in reducing the remaining $`\mathrm{q}\overline{\mathrm{q}}\mathrm{}\nu `$ final states. To improve the performance of the likelihood analysis we use only events with : $`|m_{\mathrm{vis}}m_{\mathrm{recoil}}|<50`$ GeV, $`(m_{\mathrm{vis}}+m_{\mathrm{recoil}})/\sqrt{s}>0.89`$ and $`R_{p_{\mathrm{t}i}}>1.2`$. Events are then selected using $`_{\mathrm{q}\overline{\mathrm{q}}\nu \overline{\nu }}>0.5`$, where $`_{\mathrm{q}\overline{\mathrm{q}}\nu \overline{\nu }}`$ is the likelihood for the $`\mathrm{q}\overline{\mathrm{q}}\nu \overline{\nu }`$ selection. The likelihood distribution of data and Monte Carlo is shown in Figure 3c. For the $`\mathrm{b}\overline{\mathrm{b}}\nu \overline{\nu }`$ selection we require, in addition, the b-tag variable of Reference to be greater than 0.65.
The efficiencies for the $`\mathrm{q}\overline{\mathrm{q}}\nu \overline{\nu }`$ selection alone are $`(30.5\pm 2.0)`$% at 183 GeV and $`(33.9\pm 2.3)`$% at 189 GeV. The errors on these efficiencies include the systematic errors discussed below in Section 3.6. The efficiencies after considering the results of the b-tagging, as well as the number of events selected at the two energies are given in Table 1 (lines k and l).
### 3.5 Selection of $`𝐙𝐙\mathbf{}𝐪\overline{𝐪}𝐪\overline{𝐪}`$ events
The fully hadronic channel of the $`\mathrm{Z}`$-pair decay has the largest branching fraction of all channels (about 50%), but suffers from large background from hadronic $`\mathrm{W}`$-pair decays. We apply two different event selections, one of which is based mainly on reconstructed mass information in order to accept all hadronic $`\mathrm{Z}`$-pair decays without flavor requirement, while a second analysis applies a flavor tag in order to select final states involving b quarks, allowing for looser requirements on the reconstructed boson mass.
For both subsamples, hadronic $`\mathrm{Z}^{}/\gamma ^{}\mathrm{q}\overline{\mathrm{q}}`$ events are an important background. We therefore start with a common preselection based on event shape variables which is mainly aimed at reducing this background. We use a likelihood method, described below in Section 3.5.1, in order to choose the most likely jet pairing for $`\mathrm{Z}`$-pair decays.
#### 3.5.1 Preselection and jet pairing
The event selection starts from the inclusive multihadron selection described in Reference . The radiative process $`\mathrm{e}^+\mathrm{e}^{}\mathrm{Z}\gamma \mathrm{q}\overline{\mathrm{q}}\gamma `$ is suppressed by requiring the effective center-of-mass energy after initial state radiation, $`\sqrt{s^{}}`$, to be larger than 150 GeV. $`\sqrt{s^{}}`$ is obtained from a kinematic fit that allows for one or more radiative photons in the detector or along the beam pipe. The final state particles are then grouped into jets using the Durham algorithm . A four-jet sample is formed by requiring the jet resolution parameter $`y_{34}`$ to be at least $`0.003`$ and each jet to contain at least two charged tracks. In order to suppress $`\mathrm{Z}^{}/\gamma ^{}\mathrm{q}\overline{\mathrm{q}}`$ background, the event shape parameter $`C_{\mathrm{par}}`$ , which is large for spherical events, is required to be greater than 0.25. A 4C kinematic fit using energy and momentum conservation is required to converge. A 5C kinematic fit which forces the two jet pairs to have the same mass is applied in turn to all three possible combinations of the four jets. This fit is required to converge for at least one combination. The efficiencies of these preselection cuts are (86.4 $`\pm `$ 0.5) % and (88.9 $`\pm `$ 0.5) % for signal events at 183 GeV and 189 GeV, respectively.
In order to determine which pair of jets comes from each $`\mathrm{Z}`$, we calculate a likelihood function using the mass obtained from the 5C fit, the corresponding fit probability, and the difference between the two di-jet masses obtained from the 4C fit. In the YFSZZ simulation of $`\mathrm{Z}`$-pair decays the correct jet pairing has the highest likelihood output in (86.8 $`\pm `$ 0.5) % of the events. This fraction rises to (93.8 $`\pm `$ 0.5) % for the events after the final selection.
#### 3.5.2 Likelihood for the inclusive $`𝐙𝐙\mathbf{}𝐪\overline{𝐪}𝐪\overline{𝐪}`$ event selection
We use a likelihood selection with eight input variables for the selection of $`\mathrm{Z}\mathrm{Z}\mathrm{q}\overline{\mathrm{q}}\mathrm{q}\overline{\mathrm{q}}`$ events. The first variable is the jet pairing likelihood described above. Excluding the jet pairing with the largest difference between the two di-jet masses as obtained from the 4C fit, we identify among the remaining two possible pairings the one for which the 5C-fit mass is closer to the $`\mathrm{W}`$ mass. We use the difference between this 5C-fit mass and the $`\mathrm{W}`$ mass in order to discriminate against hadronic $`\mathrm{W}`$-pair events. Two variables that are sensitive to unobserved particles along the beam direction are the fitted center-of-mass energy and the sum of the cosines of the polar angles of the four jets. In order to discriminate against $`\mathrm{Z}^{}/\gamma ^{}`$ events, we use the difference between the largest and smallest jet energies after the 4C fit, and the angular variable $`j_{\mathrm{ang}}=E_4(1\mathrm{cos}\theta _{12}\mathrm{cos}\theta _{13}\mathrm{cos}\theta _{23})/\sqrt{s}`$, where $`E_4`$ is the smallest of the four jet energies, and the $`\theta _{ij}`$ are the opening angles between jets $`i`$ and $`j`$, with the jets ordered by energy. Finally we calculate from the momenta configuration of the four jets the effective matrix element for the QCD processes $`\mathrm{Z}^{}/\gamma ^{}\mathrm{q}\overline{\mathrm{q}}\mathrm{gg}`$ and $`\mathrm{Z}^{}/\gamma ^{}\mathrm{q}\overline{\mathrm{q}}\mathrm{q}\overline{\mathrm{q}}`$ as defined in Reference , and the matrix element for the process $`\mathrm{WW}\mathrm{q}\overline{\mathrm{q}}\mathrm{q}\overline{\mathrm{q}}`$ from Reference .
The distribution of the likelihood function calculated from these eight variables is shown in Figure 3d. In order to maximize the significance of the measured cross section, assuming a 10% relative systematic error on the background, we place a cut on the likelihood at 0.65. This cut leads to an efficiency of $`(33.0\pm 1.8)`$% (189 GeV) relative to all fully hadronic NC2 $`\mathrm{Z}`$-pair final states. In the data, 52 events are selected. We expect a total of 57.0 events from Standard Model processes, of which 17.6 originate from the $`\mathrm{Z}`$-pair signal, while 27.0 events are expected from hadronic $`\mathrm{W}`$-pair decays and 12.4 events from hadronic two-fermion processes. At 183 GeV the selection efficiency for fully hadronic $`\mathrm{Z}`$-pair decays is $`(25.3\pm 1.3)`$%. In the data 8 events are selected. The Standard Model expectation is 1.7 signal events and 7.2 background events.
#### 3.5.3 Likelihood for $`𝐙𝐙\mathbf{}𝐪\overline{𝐪}𝐛\overline{𝐛}`$ event selection
Jets originating from b-quarks are selected using the same b-tagging algorithm used in the $`\mathrm{b}\overline{\mathrm{b}}\mathrm{}^+\mathrm{}^{}`$ and $`\mathrm{b}\overline{\mathrm{b}}\nu \overline{\nu }`$ selections. We evaluate the probability for each of the four jets to originate from a primary b quark, and use the two highest probabilities as input variables for a likelihood to select $`\mathrm{Z}\mathrm{Z}\mathrm{q}\overline{\mathrm{q}}\mathrm{b}\overline{\mathrm{b}}`$ events. In addition, we use the parameters $`y_{34}`$, $`C_{\mathrm{par}}`$, the difference between the largest and smallest jet energies and the output of the jet pairing likelihood. We also use the fit probabilities of a 5C kinematic fit which constrains one boson mass to the $`\mathrm{Z}`$ mass, and the probability of a 6C fit which forces both masses to be equal to the $`\mathrm{W}`$ mass.
Figure 3e shows the distribution of the likelihood function calculated from these eight variables for the preselected events. The signal likelihood is required to be larger than 0.80 for both 183 GeV and 189 GeV data. After the likelihood selection, we perform an additional cut on the mass obtained from the 5C-mass fit in the most likely jet pairing, which is required to be larger than 86 GeV. This choice of the cuts on the likelihood and on the 5C-fit mass was made by maximizing the expected statistical significance of signal over background. The final efficiency is $`(35.1\pm 2.4)`$% for the 183 GeV data and $`(38.6\pm 2.7)`$% for the 189 GeV data. The observed number of events, expected signal and background are given in Table 1 (line o).
#### 3.5.4 Combination
To account for overlap between the $`\mathrm{q}\overline{\mathrm{q}}\mathrm{q}\overline{\mathrm{q}}`$ and $`\mathrm{q}\overline{\mathrm{q}}\mathrm{b}\overline{\mathrm{b}}`$ selections we divide the data into three logical classes, exclusive $`\mathrm{q}\overline{\mathrm{q}}\mathrm{q}\overline{\mathrm{q}}`$ events (Table 1 (line m)), exclusive $`\mathrm{q}\overline{\mathrm{q}}\mathrm{b}\overline{\mathrm{b}}`$ (line n) events and the overlapping region (Table 1 (line o)). In Table 1 (line m) we give the efficiency relative to all fully hadronic final states without b-quarks. The other efficiencies are given relative to fully hadronic final states with b-quarks.
### 3.6 Selection systematic errors
Systematic errors have only a modest effect on our final result because of the large statistical error associated with the small $`\mathrm{Z}`$-pair cross section.
Detector effects can best be studied by comparing calibration data taken at the $`\mathrm{Z}`$ resonance with a simulation of the same process. These comparisons are important for final states such as $`\mathrm{}^+\mathrm{}^{}\nu \overline{\nu }`$, $`\mathrm{q}\overline{\mathrm{q}}\nu \overline{\nu }`$, and $`\mathrm{q}\overline{\mathrm{q}}\mathrm{q}\overline{\mathrm{q}}`$, where the tight cuts are needed to separate signal and background. In these cases, we add additional smearing to the energy and momentum of the simulated events to match data and simulation. We then apply the same smearing to the signal and background Monte Carlos and then correct our efficiencies and background accordingly. The full difference is used as the systematic error in these cases. At $`\sqrt{s}=189`$ GeV, these differences give relative systematic errors on the efficiency of 2.5% for the $`\mathrm{e}^+\mathrm{e}^{}\nu \overline{\nu }`$ final state, 5.2% for the $`\mu ^+\mu ^{}\nu \overline{\nu }`$ final state and 3.8% for $`\mathrm{q}\overline{\mathrm{q}}\nu \overline{\nu }`$ final state. In the $`\mathrm{q}\overline{\mathrm{q}}\mathrm{q}\overline{\mathrm{q}}`$ inclusive analysis, where uncertainties on the reconstructed angles are also important, similar studies lead to a relative detector systematic error of 6%.
Detector systematic errors for the $`\mathrm{q}\overline{\mathrm{q}}\mathrm{}^+\mathrm{}^{}`$ and $`\mathrm{}^+\mathrm{}^{}\mathrm{}^+\mathrm{}^{}`$ selections without $`\tau `$-pairs in the final state are small because of the good separation of signal and background. In the $`\mathrm{}^+\mathrm{}^{}\mathrm{}^+\mathrm{}^{}`$ final state the largest effect (3%) is from modeling of the multiplicity requirement which is important for final states containing $`\tau `$-pairs. In the $`\mathrm{q}\overline{\mathrm{q}}\tau ^+\tau ^{}`$ selection, the systematic uncertainties in the efficiencies were determined by overlaying hadronic and tau decays taken from $`\mathrm{Z}`$ resonance data giving a contribution to the systematic error of 6.0%. When the final states are combined to determine the cross section and limits on the anomalous triple gauge couplings, we assume a common relative systematic error of 3% on the efficiencies.
Another important detector effect comes from the simulation of the variables used by the OPAL b-tag which is discussed in Reference . We allow for a common 5% error in the efficiency of the b-tag, consistent with our studies on $`\mathrm{Z}`$ resonance data and Monte Carlo.
These detector effects were propagated through to our background errors.
In each channel the signal and background Monte Carlo generators have been compared against alternative generators. In almost all cases the observed differences are consistent within the finite Monte Carlo statistics and the systematic error has been assigned accordingly. One notable exception is in the $`\mathrm{q}\overline{\mathrm{q}}\mathrm{b}\overline{\mathrm{b}}`$ background where differences in the PYTHIA, grc4f and EXCALIBUR simulation of the $`\mathrm{W}`$-pair background are as large as 20%. In this case the full difference has been assigned as the systematic error.
In the $`\mathrm{}^+\mathrm{}^{}\mathrm{}^+\mathrm{}^{}`$ channels, which has a large background from two-photon events, we have compared the number of selected events at an early stage of the analysis with the Monte Carlo prediction and based our background systematic error on the level of agreement. This results in 20% systematic error.
## 4 Results
We combine the information from all of the analyses reported above using a maximum likelihood fit to determine the production cross section for $`\mathrm{e}^+\mathrm{e}^{}\mathrm{Z}\mathrm{Z}`$. The information which was used in the fit, as well as the Standard Model prediction for $`\mathrm{Z}`$-pair production, is summarized in Table 1. For each channel the table gives the number of events observed, $`n_{\mathrm{obs}}`$, the Standard Model prediction for all events, $`n_{\mathrm{SM}}`$, the expected signal, $`n_{\mathrm{ZZ}}`$, the expected background, $`n_{\mathrm{back}}`$, the efficiency $`ϵ_{\mathrm{chan}}`$, and the integrated luminosity, $`L_{\mathrm{int}}`$. $`B_{\mathrm{ZZ}}`$ is the branching ratio of $`\mathrm{Z}`$-pairs to the given final state, calculated from $`\mathrm{Z}`$ resonance data . In the table we give the overlap between the b-tag and non-b-tag analyses. Possible overlap between $`\mathrm{q}\overline{\mathrm{q}}\mathrm{q}\overline{\mathrm{q}}`$ and $`\mathrm{q}\overline{\mathrm{q}}\mathrm{}^+\mathrm{}^{}`$ has been studied, and found to be an order of magnitude smaller than the overlap of $`\mathrm{q}\overline{\mathrm{q}}\mathrm{q}\overline{\mathrm{q}}`$ and $`\mathrm{q}\overline{\mathrm{q}}\mathrm{b}\overline{\mathrm{b}}`$ and has therefore been ignored.
The cross section at each energy is determined with a maximum likelihood fit using a Poisson probability density convolved with Gaussians to describe the uncertainties on efficiencies and backgrounds. The expected number of events in each channel, $`\mu _\mathrm{e}`$, as function of the $`\mathrm{Z}`$-pair cross section, $`\sigma _{\mathrm{ZZ}}`$, is given by
$$\mu _\mathrm{e}=\sigma _{\mathrm{ZZ}}L_{\mathrm{int}}ϵ_{\mathrm{chan}}B_{\mathrm{ZZ}}+n_{\mathrm{back}}.$$
(2)
The efficiencies, $`ϵ_{\mathrm{chan}}`$, include the effects of off-shell $`\mathrm{Z}`$ bosons that are produced outside of our kinematic acceptance. Our main result, the NC2 $`\mathrm{Z}`$-pair cross sections obtained from the fits, is
$$\begin{array}{cccc}\sigma _{\mathrm{ZZ}}(183\mathrm{GeV})\hfill & =& 0.12_{0.18}^{+0.20}(\text{stat.})_{0.02}^{+0.03}(\text{syst.})\hfill & \mathrm{pb}\hfill \\ & & & \\ \sigma _{\mathrm{ZZ}}(189\mathrm{GeV})\hfill & =& 0.80_{0.13}^{+0.14}(\text{stat.})_{0.05}^{+0.06}(\text{syst.})\hfill & \mathrm{pb}.\hfill \end{array}$$
The 183 GeV result corresponds to a 95% C.L. upper limit (normalized to the region $`\sigma _{\mathrm{ZZ}}>0`$) of $`0.55\mathrm{pb}`$. The comparison of these measurements with the YFSZZ prediction, using the coupling of $`\mathrm{Z}`$ bosons to electrons measured at the $`\mathrm{Z}`$ resonance, is shown in Figure 4. The results are consistent with the YFSZZ prediction and with the measurements presented in Reference .
Our results assume that, apart from the backgrounds discussed above, only $`\mathrm{Z}`$$`\mathrm{Z}`$ production contributes inside our kinematic region. Possible effects of Higgs boson production have been ignored.
In order to check the results for consistency with the expected fraction of $`\mathrm{b}\overline{\mathrm{b}}`$ final states, we perform a second fit where the branching ratio of the $`\mathrm{Z}`$ boson to $`\mathrm{b}\overline{\mathrm{b}}`$ is a free parameter. The relative branching ratios of the $`\mathrm{Z}`$ to other fermion pairs are fixed to their measured values. The resulting branching ratio and cross section for the 189 GeV data are $`\text{Br(}\mathrm{Z}\mathrm{b}\overline{\mathrm{b}})=0.21_{0.06}^{+0.07}(\mathrm{stat}.)\pm 0.01(\mathrm{syst}.)`$ and $`\sigma _{\mathrm{ZZ}}=0.75_{0.14}^{+0.15}(\text{stat.})_{0.05}^{+0.06}(\text{syst.})\text{pb}`$. The measured branching ratio can be compared with the world average as measured at the $`\mathrm{Z}`$ resonance of $`\text{Br(}\mathrm{Z}\mathrm{b}\overline{\mathrm{b}})=0.1516\pm 0.0009`$. The cross section measured without constraining the branching ratio is also consistent with our main result and the YFSZZ prediction. The smaller error on the main result cross section illustrates the advantage of classifying the hadronic systems as $`\mathrm{b}\overline{\mathrm{b}}`$ or non-$`\mathrm{b}\overline{\mathrm{b}}`$. The 183 GeV data sample is too small to extract a meaningful value of the $`\mathrm{Z}\mathrm{b}\overline{\mathrm{b}}`$ branching ratio.
Limits on anomalous triple gauge couplings were set using the total cross section and the $`|\mathrm{cos}\theta _\mathrm{Z}|`$ distribution of our data. Here $`\theta _\mathrm{Z}`$ is the polar angle of the $`\mathrm{Z}`$ bosons produced. In this study we varied the real and imaginary parts of the $`\mathrm{Z}`$$`\mathrm{Z}`$$`\mathrm{Z}`$ and $`\mathrm{Z}\mathrm{Z}\gamma `$ anomalous couplings parameterized by the form factors $`f_4^{\mathrm{ZZZ}}`$, $`f_5^{\mathrm{ZZZ}}`$, $`f_4^{\mathrm{ZZ}\gamma }`$ and $`f_5^{\mathrm{ZZ}\gamma }`$ as defined in Reference and implemented in the YFSZZ Monte Carlo. The real and imaginary parts of each coupling were varied separately with all others fixed to zero.
For this study we consider the effect of the anomalous couplings on the total cross section at $`\sqrt{s}=183`$ GeV and on the cross section in four bins of $`|\mathrm{cos}\theta _\mathrm{Z}|`$ at $`\sqrt{s}=189`$ GeV. The selection efficiencies for all final states are parameterized as function of anomalous couplings. At $`\sqrt{s}=189`$ GeV the parameterization is done separately for each bin in $`|\mathrm{cos}\theta _\mathrm{Z}|`$. For values of the anomalous couplings larger than unity, much of the production of the final state fermions occurs at $`|\mathrm{cos}\theta |1`$ where the efficiency for most channels is reduced by a factor of $`0.5`$. An uncertainty of 10%, dominated by Monte Carlo statistical errors, is assigned to the correction we apply to these efficiencies.
The 95% C.L. limits on the anomalous couplings obtained from the maximum likelihood fit are given in Table 2. With the exception of the $`Re\{f_5^{\mathrm{ZZZ}}\}`$ coupling, the limits are insensitive to the sign and complex phase of the couplings.
## 5 Conclusion
The production cross section of $`\mathrm{e}^+\mathrm{e}^{}\mathrm{Z}\mathrm{Z}`$ has been measured using the final states $`\mathrm{}^+\mathrm{}^{}\mathrm{}^+\mathrm{}^{}`$, $`\mathrm{}^+\mathrm{}^{}\nu \overline{\nu }`$, $`\mathrm{q}\overline{\mathrm{q}}\mathrm{}^+\mathrm{}^{}`$, $`\mathrm{q}\overline{\mathrm{q}}\nu \overline{\nu }`$, and $`\mathrm{q}\overline{\mathrm{q}}\mathrm{q}\overline{\mathrm{q}}`$. The number of observed events, the background expectation from Monte Carlo and the calculated efficiencies have been combined to measure the production cross section of the process $`\mathrm{e}^+\mathrm{e}^{}\mathrm{Z}\mathrm{Z}`$. Our measured cross sections include the effects of background and efficiency uncertainties.
We have determined the cross section for $`\mathrm{e}^+\mathrm{e}^{}\mathrm{Z}\mathrm{Z}`$ separately at average center-of-mass energies of $`182.62\pm 0.05`$ GeV and $`188.63\pm 0.04`$ GeV. The NC2 $`\mathrm{Z}`$-pair cross sections were determined to be
$$\begin{array}{cccc}\sigma _{\mathrm{ZZ}}(183\mathrm{GeV})\hfill & =& 0.12_{0.18}^{+0.20}(\text{stat.})_{0.02}^{+0.03}(\text{syst.})\hfill & \mathrm{pb}\hfill \\ & & & \\ \sigma _{\mathrm{ZZ}}(189\mathrm{GeV})\hfill & =& 0.80_{0.13}^{+0.14}(\text{stat.})_{0.05}^{+0.06}(\text{syst.})\hfill & \mathrm{pb}.\hfill \end{array}$$
At the lower center-of-mass energy, the 95% C.L. upper limit on the cross section is 0.55 pb. The measurements at both energies are consistent with the Standard Model expectations. No evidence is found for anomalous neutral current triple gauge couplings. The 95% confidence level limits are listed in Table 2.
## Acknowledgements
We particularly wish to thank the SL Division for the efficient operation of the LEP accelerator at all energies and for their continuing close cooperation with our experimental group. We thank our colleagues from CEA, DAPNIA/SPP, CE-Saclay for their efforts over the years on the time-of-flight and trigger systems which we continue to use. In addition to the support staff at our own institutions we are pleased to acknowledge the
Department of Energy, USA,
National Science Foundation, USA,
Particle Physics and Astronomy Research Council, UK,
Natural Sciences and Engineering Research Council, Canada,
Israel Science Foundation, administered by the Israel Academy of Science and Humanities,
Minerva Gesellschaft,
Benoziyo Center for High Energy Physics,
Japanese Ministry of Education, Science and Culture (the Monbusho) and a grant under the Monbusho International Science Research Program,
Japanese Society for the Promotion of Science (JSPS),
German Israeli Bi-national Science Foundation (GIF),
Bundesministerium für Bildung, Wissenschaft, Forschung und Technologie, Germany,
National Research Council of Canada,
Research Corporation, USA,
Hungarian Foundation for Scientific Research, OTKA T-029328, T023793 and OTKA F-023259.
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# Reversing type II migration: resonance trapping of a lighter giant protoplanet
## 1 Introduction
In the past few years a number of extrasolar giant planets have been discovered around nearby solar–type stars. These objects masses range from $`0.17`$ $`M_J`$ to $`11`$ $`M_J`$ (where $`M_J`$ is Jupiter’s mass) and their orbital semi-major axis range from $`0.038`$ AU to $`3.3`$ AU (Marcy, Cochran & Mayor, 1999). Although many uncertainties remain about planet formation, it is now commonly accepted that planets have formed in and from protoplanetary discs. Necessarily, there must be some time interval over which a giant planet and the surrounding gaseous disc material coexist. The planet and the disc exchange angular momentum through tidal interactions which generally make the planet lose angular momentum This mechanism is called migration. It can roughly be divided in two regimes:
* If the planet mass is small enough, the disc response is linear. The migration rate, in that regime, is proportional to the planet and disc masses, is independent of the viscosity and weakly dependent of the disc surface density and temperature profiles. This is the so-called type I migration (Ward, 1997).
* When the protoplanet mass is above a certain threshold, the torques acting locally on the surrounding disc material open a gap (Papaloizou & Lin, 1984), whose width and depth are controlled by the balance between the tidal torques, which tend to open the gap, and the viscous torques which tend to close it. The disc response is significantly non-linear, and most of the protoplanet Lindblad resonances fall in the gap and therefore cannot contribute to the planet-disc angular momentum exchange. The migration rate slows down dramatically compared to type I migration. Furthermore, the tidal truncation splits the disc into two parts and the planet is locked to the disc viscous evolution (Nelson et al. 2000). This is the type II migration, which describes the orbital evolution of giant protoplanets.
In this letter we consider the coupled evolution of a system of giant protoplanets consisting of two non-accreting cores with masses $`1`$ $`M_J`$ and $`0.29`$ $`M_J`$, which we are going to call from now on respectively “Jupiter” and “Saturn”. Attempts have already been made to describe the behaviour of a system of planets embedded in a disc. Melita & Woolfson (1996) and Haghighipour (1999) considered an embedded Jupiter and Saturn system orbiting a solar mass star, and showed how resonance trapping would affect their evolution. However the dissipative force in these works was due to the dynamical friction with a uniform density interplanetary medium, hence type II migration effects were not taken into account. Resonance trapping of planetesimals by a fixed orbit Jupiter sized protoplanet has also been investigated by Beaugé et al. (1994), and shown to be able to build up a single planetary core with orbital characteristics close to Saturn’s ones. Kley (2000) studied the orbital evolution of two maximally accreting giant cores embedded in a minimal mass protosolar disc, and showed how the migration of the inner core could be halted by the presence of the outer one, and how the eccentricity of the inner core is pumped up by the outer one.
## 2 Results
### 2.1 Numerical codes description
In order to investigate the long-term behaviour of the embedded Jupiter and Saturn system, we have used two independent hydrocodes, which have been described elsewhere in full detail (Nelson et al. 2000). These two codes are fixed Eulerian grid based codes, one of them is NIRVANA (Ziegler & Yorke, 1997) and the other one has been written by one of us (FM). Both have been endowed with the fast advection FARGO algorithm (Masset, 2000), and can run either with this algorithm or with a standard advection algorithm. They gave very similar results. They consist of a pure N-body kernel based on either a fourth (NIRVANA) or fifth order adaptative timestep Runge-Kutta solver (sufficient for the short time-scales involved in this dissipative problem) embedded in a hydrocode which provides a tidal interaction with a 2D non self-gravitating gaseous disc. The simulations are performed in the non-inertial non-rotating frame centered on the primary. The grid outer boundary does not allow inflow nor outflow and is chosen sufficiently far from the planets in order for the spiral density waves that they launch to be damped before they reach it, while the grid inner boundary only allows outflow (inwards), so that the disc material can be accreted on to the primary. Failing to do so may lead to overestimate the inner disc density and artificially favours an outwards migration. In the following our length unit is $`5.2`$ AU, the mass unit is one solar mass, and the time unit is the initial orbital period of Jupiter (the actual period may vary as Jupiter migrates). The disc aspect ratio $`H/R`$ is uniform and constant. In the run presented here the grid resolution adopted is $`N_r=122`$ and $`N_\theta =300`$ with a geometric spacing of the interzone radii — such that all the zones are “as square as possible”, i.e. $`N_r\mathrm{log}(1+2\pi /N_\theta )=\mathrm{log}(R_{max}/R_{min})`$ —. The grid outer boundary is at $`R_{max}=5`$ and its inner boundary is at $`R_{min}=0.4`$. The geometric spacing is the most natural one since the disc thickness scales as $`r`$. On the other hand, a constant spacing leads to an oversampling of the outer disc and an undersampling of the inner one, and therefore is likely to favour an inwards migration.
### 2.2 Initial setup
The cores we consider are embedded in a gaseous minimal mass protosolar nebula around a unit mass central object, and we assume they start their evolution with semi-major axis $`a_j=1`$ for Jupiter and $`a_s=2`$ for Saturn. The disc surface density is uniform and corresponds to two Jupiter masses inside Jupiter’s orbit. The effective viscosity $`\nu `$, the nature of which remains unclear and is usually thought to be due to turbulence generated by MHD instabilities (Balbus & Hawley 1991), is assumed to be uniform through the disc and corresponds to a value of $`\alpha 610^3`$ in the vicinity of Jupiter’s orbit. The disc aspect ratio is $`H/r=0.04`$.
The mass of Jupiter is sufficient to open a deep gap and hence it settles in a type II migration (Nelson et al. 2000), whereas Saturn is unable to fully empty its coorbital region because: (i) its mass is smaller; (ii) The planet is in a regime known as the inertial limit (Ward & Hourigan, 1989) where the inwards migration speed is so high that it makes the planet pass through what would be the gap inner edge before it had time to actually open it.
Therefore Saturn does not clear a deep gap initially, and its migration rate is typical of type I migration, since all its Lindblad resonances can still contribute to the angular momentum exchange with the disc.
### 2.3 Run results
We present in fig. 1 the central star–planet distance curves as a function of time. We see how initially Jupiter migrates as if it was the only planet in the disc (see test run). In the meantime, Saturn starts a much faster migration (the obvious initial acceleration of its migration will be discussed elsewhere), and reaches the 1:2 resonance with Jupiter at time $`t110`$. The eccentricities at that time are small (see fig. 2), and in particular Saturn’s eccentricity is much smaller than the eccentricity threshold below which the capture into resonance is certain if the “adiabatic” condition on the migration rate is satisfied (Malhotra 1993): $`|\underset{s}{\overset{.}{a}}|/(a_s\mathrm{\Omega }_s)0.5j(j+1)\mu _Je_s`$ for the $`j`$:$`j+1`$ resonance, where $`\mu _J`$ is the mass ratio of Jupiter to the central object, and where $`e_s`$ is Saturn’s eccentricity. This condition is not satisfied when Saturn reaches the 1:2 resonance, and it passes through.
The planets then obtain higher eccentricities, and Saturn’s migration rate is reduced. Saturn’s eccentricity increases again rapidly as it passes through the 3:5 resonance with Jupiter at $`t220`$. Eventually the adiabatic condition on the migration rate is satisfied for the 2:3 resonance and Saturn’s eccentricity is still below the corresponding critical threshold, so it gets trapped into the 2:3 resonance with Jupiter (both $`e`$ and $`e^{}`$ resonances, since the two critical angles $`\varphi =3\lambda _s2\lambda _j\stackrel{~}{\omega }_s`$ and $`\varphi ^{}=3\lambda _s2\lambda _j\stackrel{~}{\omega }_j`$ librate, where $`\lambda `$ is the mean longitude and $`\stackrel{~}{\omega }`$ the longitude of perihelion). At that time both planets steadily migrate outwards.
### 2.4 Interpretation
We define the system of interest as the system composed of the two planets. This resonance locked system interacts with the inner disc through torques proportional to $`M_J^2`$, at Jupiter’s inner Lindblad resonances (ILR), whereas it interacts with the outer disc through torques proportional to $`M_S^2`$ at Saturn’s outer Lindblad resonances (OLR), as indicated on fig. 3.
It can be seen that Saturn’s ILR fall in Jupiter’s gap and Jupiter’s OLR fall in Saturn’s gap so their effect is weakened compared to the situation where Jupiter is alone. As $`M_J^2/M_S^210`$, the torque imbalance does not favour an inwards migration as strongly as in a one planet case, and may even lead to a positive differential Lindblad torque on the two planet system. Actually one can estimate what the maximum mass ratio of the outer planet to the inner one should be to get a migration reversal, if one neglects the Inner Lindblad torque on the outer planet and the Outer Lindblad torque on the inner planet. The Inner Lindblad torque on the inner planet reads as:
$$T_{ILR}=C_{ILR}\mu _J^2\mathrm{\Sigma }_0a_J^2(a_J\mathrm{\Omega }_J)^2h^3$$
(1)
where $`C_{ILR}`$ is a dimensionless coefficient which is a sizable fraction of unity (Ward, 1997), and where $`h^{}`$ is the disc aspect ratio. There is a similar formula for the Outer Lindblad torque on the outer planet (obtained by substituting the $`ILR`$ and $`J`$ indices in Eq. (1) respectively with $`OLR`$ and $`S`$). The resulting torque imbalance will be positive if: $`T_{ILR}>T_{OLR}`$, which reads here as:
$$\frac{\mu _S}{\mu _J}<\left(\frac{C_{ILR}}{C_{OLR}}\right)^{1/2}\left(\frac{2}{3}\right)^{1/3}$$
(2)
If we assume that $`C_{ILR}=C_{OLR}`$ then we get: $`\mu _S/\mu _J<0.87`$, whereas if we make the conservative assumption that $`C_{ILR}=\frac{1}{2}C_{OLR}`$, we have: $`\mu _S/\mu _J<0.62`$. This threshold is much bigger than the actual ratio, therefore if the common gap is deep enough to shut off Jupiter’s OLR torques (and Saturn’s ILR torques) then the net Lindblad torque on the two planet system is positive. As the two planet system proceeds outwards in the disc, it does not act on the gas as a snow-plough, but rather it allows the material from the outer disc to travel across the common gap and eventually feed the inner disc. We can find the gap “permeability” condition by requiring that the rate of angular momentum change of the ring of material lying immediately outside Saturn’s gap that is required to expand accordingly to Saturn’s orbit (snow-plough effect) is greater than the torque available from Saturn (at most the sum of its outer Lindblad torques, in which case we need to assume that the waves excited at their OLR are damped locally). In our case, this turns out not to be the case and most of the outer disc material flows through the common gap to the inner disc. We find that in all our runs it is possible to check that the rate of mass flow through the common gap (see fig. 4) can be expressed as :
$$\stackrel{.}{M}3\pi \nu \mathrm{\Sigma }_0+2\pi r_s\underset{s}{\overset{.}{r}}\mathrm{\Sigma }_0$$
(3)
with a reasonable precision ($`1020`$ %). Furthermore we have performed many “restart runs” which consist in restarting a run once Jupiter and Saturn are locked into resonance, and then by varying one parameter at one time, e.g. the viscosity or the aspect ratio (which changes the Lindblad torques and therefore the migration rate). The mass flux through the gap rapidly switches (in a few tens of orbits) to a new value after the restart, so that Eq. (3) remains fulfilled.
From the considerations above we can conclude that the presence of Saturn unlocks Jupiter from the disc evolution : the two planet system evolution (outwards) and the disc viscous evolution (inwards) are basically decoupled. This decoupling and the corresponding mass flow through the common gap has two consequences :
* A refilling of the inner disc, which is too depleted for the torques at Jupiter’s ILR to have any sizable effect in the one planet case (the inner disc is accreted on the primary on its short viscous time-scale and maintaining its surface density at not too low a value implies a permanent flow of material from the outer disc to the inner).
* The angular momentum lost by the material which flows from the outer disc to the inner one is gained by the planets. The exchange of angular momentum between a planet and a gas fluid element occurs during a “close encounter” between these two, the one planet version of which corresponds to the angular momentum exchange at each end of a horseshoe orbit of the fluid element. The resulting torque is the so-called coorbital corotation torque (Goldreich & Tremaine 1979, Ward 1991 and 1992). To the best of our knowledge, an analytical evaluation of the corotation torque in the case of a non-vanishing net mass flow through the orbit (either due to a viscous accretion on to the primary or to a radial migration or both) has not been performed yet. Obviously even the one planet case deserves a large amount of work on this specific topic, therefore the estimate of the corotation torque in this two planet problem is far beyond the scope of this paper. We will just comment that the corotation torque in our case might not be negligible compared to the differential Lindblad torque at some stage.
## 3 Discussions
We have performed a series of restart runs (see section 3) in order to check for a variety of behaviours.
### 3.1 Differential Lindblad torque sign
The one sided Lindblad torque has been shown to be proportional to $`h^3`$ (Ward 1997). We have performed two restart runs ($`h^{}=0.040.03`$ and $`h^{}=0.040.05`$) in order to check that the migration rate variation is consistent with this dependence. This is indeed the case. We note in passing that the migration rate varies as $`h^3`$, and not as $`h^2`$ as it would be the case in a one planet problem, since the Outer/Inner Lindblad torque asymmetry does not vanish as the disk thickness tends to zero (the OLRs would pile-up at Saturn’s orbit, whereas the ILRs would pile-up at Jupiter’s orbit). These results confirm that the behaviour we observe occurs mainly due to the differential Lindblad torque and shows as well that this latter quantity is positive, as expected from Eq. (2).
### 3.2 $`\alpha `$-viscosity vs. uniform viscosity
So far we have only considered a uniform viscosity. Switching to a uniform-$`\alpha `$ viscosity of the form $`\nu =\alpha c_sH`$ makes $`\nu `$ scale here as $`r^{1/2}`$, so the viscosity at the outer edge of the common gap is higher, whereas it is smaller in the inner disc. This has the following effect, which plays in favour of enhancing the migration reversal mechanism: the viscous time-scale of the inner disc is higher and therefore its surface density increases accordingly, since the material brought through the gap piles-up in the inner disc for a longer time before being accreted on the primary. This has been checked with a restart run.
### 3.3 Accretion on to the planets
The cores considered above do not accrete gas from the disc. One can wonder what would be the effects of accretion. We have performed a number of restart runs in order to investigate the effect of accretion on the mechanism presented here. We have only considered accretion on to Jupiter, as it is likely that the accretion rate on Saturn can be regarded as being negligible (i.e. its mass doubling time is much longer than the timescale of the outwards migration, see e.g. Pollack et al. 1996). The prescription we used to model accretion on to Jupiter consists in removing a proportion of the material which lies in the inner Roche lobe (i.e. a sphere with a radius of half the Hill radius). The amount which is removed in one timestep is calculated from the half emptying time of the inner Roche lobe $`\tau _{1/2}`$. We have performed four different restart runs, corresponding to the following values of $`\tau _{1/2}`$: $`\tau _{1/2}=T_0`$ (maximally accreting core, see Kley 1999), $`\tau _{1/2}=3T_0`$, $`\tau _{1/2}=10T_0`$ and $`\tau _{1/2}=30T_0`$, where $`T_0=2\pi /\mathrm{\Omega }_J`$ is Jupiter’s orbital time. In each of these cases, turning on accretion had no impact on the system migration rate, at least in the early stages: in the first case, the mass doubling time for Jupiter is relatively short, and when Jupiter’s mass is significantly larger than its initial mass some additional effects, which will be presented in much greater detail elsewhere, affect the migration rate which then differ from the non-accreting case.
### 3.4 Smoothing
The smoothing parameter of the potential can have a dramatic impact on Saturn’s initial migration rate. This rate is controlled by a subtle balance between outer disc and inner disc torques. In the case of Saturn, all the Lindblad resonances play a role, since there is no gap. Many prescriptions for the smoothing are unable to give trustworthy results for the balance between the outer and inner torques since, depending on the prescription, these two quantities are affected in a different way. On the other hand Jupiter’s migration rate is much more robust, since the presence of the gap prevents high-$`m`$ Lindblad resonances playing a role in the migration, which is therefore controlled only by remote, low $`m`$ resonances and thus almost insensitive to the smoothing parameter. For this reason we have adopted an approach which involves choosing a smoothing prescription which endows Saturn with a migration velocity of the order of magnitude of the linear analytical predictions (type I migration), which is needed to give correct results for the capture into resonance. Once Saturn is trapped into resonance with Jupiter, it is dynamically slaved by the latter and the system evolution is only very weakly affected by the exact value of the outer disc torque exerted on Saturn. We have found that using either of the two prescriptions below satisfactorily preserves the analytical torque imbalance on Saturn and therefore gives it a type I migration rate:
* The potential of a planet acting on the disc is smoothed over the length $`\epsilon =0.4R_H`$ where $`R_H`$ is the Hill radius of the planet under consideration, whereas the potential of the disc acting back on the planet is smoothed over $`\epsilon ^{}=\sqrt{H^2+d^2}`$ where $`H`$ and $`d`$ are respectively the local disc thickness and zone diagonal. Since $`\epsilon ^{}\epsilon `$ the action-reaction law is not fulfilled and the numerical biases which arise favour an inwards migration, as can be easily checked.
* The potential of a planet acting on the disc and the potential of the disc acting on the same planet are smoothed over $`\epsilon =0.4R_H`$. This prescription does fulfill the action-reaction law.
In both these two cases, as in any other which gives Saturn a type I migration rate, including runs performed with a uniform radial spacing, the migration gets reversed. The run presented here corresponds to the first prescription.
### 3.5 Impact of mass ratio and Long-term behaviour
One can wonder about the size of the interval of “Saturn”’s mass which causes the migration to be significantly slowed down or reversed. If “Saturn” is not massive enough it will not significantly affect Jupiter’s evolution (the common “gap” will be too full on Saturn’s side, and therefore Jupiter’s OLR torques will not be shut off), whereas if it is too massive, the torque imbalance will be negative again. Work is in progress to accurately determine which range of parameters leads to a migration reversal. It should be noted that the results presented here depend on the artificial initial conditions. We have performed other runs in which Saturn is initially very close either to the $`1:2`$ or $`3:5`$ resonance, and it turns out that neither of these resonances is able to struggle against the strong Lindblad torques on Saturn: no resonance angle can be found which provides a resonant torque on Saturn which counteracts the tide. Therefore a trapping into the $`2:3`$ resonance is the most likely outcome when the system is still embedded in a massive disc, whatever the initial conditions: catching-up of “Saturn” or in-situ assembling from smaller, type I migrating bodies.
The long-term behaviour of the system is twofold:
* The system is locked into resonance as long as :
+ The two-planet system can adjust its resonance angle in order to prevent the planets being “pushed” towards each other by the Lindblad torques exerted by the disk on each of them. In all our runs we have never observed this behaviour. Now, given the small eccentricities involved here, and given the fact that the adiabatic criterion threshold increases as $`j(j+1)`$, the most probable outcome is that Saturn would then be captured in the next order resonance, that is to say $`3:4`$, and all the physics exposed in this paper would still be valid (presence of a common gap, sharing of the coorbital material by the two planets, mass-weighted torque imbalance, etc.)
+ The planets are not pulled apart by any other torques. Now we have mentioned the possibly important role of the coorbital corotation torque in this problem, which may be sufficient to move the planets apart at some stage, in which case we may ultimately get a low eccentricity double giant planet system when the disc disappears. This will be presented in greater detail elsewhere.
* If the planets happen to be locked into resonance at the time that the gas effects become negligible, then the system is likely to be unstable (we mentioned already that at least two angles librate simultaneously, which strongly suggests a possible chaotic behaviour; see also Kley 2000), and the most likely outcome is that one planet will be ejected whereas the other planet will end up on an eccentric orbit. This could account for the observed eccentricities of the extrasolar planets which are not orbiting close to their host star, i.e. which have not migrated all the way to the star.
## 4 Acknowledgements
We wish to thank J.C.B. Papaloizou, R.P. Nelson, C. Terquem, J.D. Larwood, A.A. Christou and an anonymous referee for useful comments and criticism. This work was partially supported (for F.M.) by the research network “Accretion onto black holes, compact stars and protostars” funded by the European Commission under contract number ERBFMRX-CT98-0195, and additionally supported (for M.S.) by funding from a PPARC research studentship. Computational resources of the Grand HPC consortium were available and are gratefully acknowledged. We thank Udo Ziegler for making a FORTRAN version of his code NIRVANA publicly available.
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# On the Theory of Quantum Oscillations of the Elastic Moduli in Layered Conductors
## I I. Introduction
Many superconducting materials with large critical parameters created in the last two decades are layered structures with metallic-type conductivity (e.g organic metals of the $`\alpha `$-(BEDT-TTF)<sub>2</sub>MHg(SCH)<sub>4</sub> group). A characteristic feature of these materials is a strong anisotropy of the conductivity in the nonsuperconducting state: the conductivity in the layer plane is much higher than in the direction normal to the layers. It is common to assume that an anisotropy in electrical conductivity is a manifestation of the quasi-two-dimensional nature of the energy spectrum of the charge carriers in layered conductors. It follows from experiments that the Fermi surface (FS) of such conductors is a system of weakly rippled cylinders (isolated or connected by links) whose axes are perpendicular to the layers. The experimental data (see, e.g. Refs. 1 ; 2 ; 3 ; 4 ; 5 ; 6 ; 7 ; 8 ; 9 ) support this assumption.
The FS of a conductor with a quasi-two-dimensional energy spectrum can be described by the following equation
$$E(𝐩)=\underset{k=0}{\overset{\mathrm{}}{}}E_k(p_x,p_y)\mathrm{cos}\frac{akp_z}{\mathrm{}}$$
(1)
where $`𝐩`$ is the electron quasimomentum; $`E_k(p_x,p_y)`$ are coefficients with dimensions of energy, $`p_z`$ is the projection of the quasimomentum on the direction normal to the layers, and $`a`$ is the distance between the layers. If we ignore the anisotropy of the energy spectrum in the layer plane, instead of (1) we can write the simpler equation:
$$E_p=\frac{p_{}^2}{2m_{}}+\underset{k=1}{\overset{\mathrm{}}{}}E_k(p_{})\mathrm{cos}\frac{akp_z}{\mathrm{}},$$
(2)
where $`𝐩_{}`$ is the projection of the quasimomentum on the layer plane, and $`m_{}`$ is the effective mass corresponding to the motion of the quasiparticles in that plane. Equation (2) describes an axially symmetric open FS with the axis of symmetry $`(z`$ axis of the chosen coordinate system) directed along a normal to the layers. The profile of the longitudinal section of the FS corresponding to Eq. (2) is a periodic function of period $`2p_02\pi \mathrm{}/a.`$
Experimental data concerning quantum oscillations in various characteristics of metals were widely used as the instruments of study of their electron energy spectra. Quantum oscillatory phenomena were repeatedly observed in organic layered conductors (Refs. 1 ; 3 ; 4 and 10 ). Quantum oscillations in thermodynamic observables in these materials can exhibit some characteristic properties which occur due to the quasi-2D character of the energy spectrum of charge carriers. For instance it was shown in Ref. 11 that the magnetic susceptibility of the noninteracting 2D electron system with open FS can exhibit sharp maxima due to the orbital effect. However, a systematic theory of quantum oscillatory phenomena in organic metals which takes into account the quasi-2D character of their electron spectra and a strong correlation among the electrons is not well developed at present. In this paper we give such theoretical analysis and we apply our results to study the quantum oscillations of the thermodynamic characteristics, especially of the elastic moduli.
At first, we derive the expressions for the electron contributions to the elastic constants corresponding to the compression in the direction perpendicular to layers $`(\lambda _0)`$ and to the compression applied along the layers $`(\lambda )`$ in the presence of external magnetic field $`𝐁.`$ In the geometry where the $`\mathrm{`}\mathrm{`}z\mathrm{"}`$ axis is directed along the normal to the layers, and assuming that the system is axially symmetric our elastic constants $`\lambda _0`$ and $`\lambda `$ correspond to the electron contribution to the elastic moduli $`c_{33}`$ and $`c_{11}=c_{22}`$ in Voigt notation. A possible instability of the lattice arises near the peaks of the oscillations of DOS in a quantizing magnetic field. Therefore we analyze these oscillations of DOS in detail, and we show that they can be strongly affected due to the interaction among the carriers. In the following section of the work we introduce a concrete model of the FS which permits us to analyze how the local geometry of the FS can strengthen the singularities in DOS which can give rise to the considered lattice instabilities.
## II 2. Main equations
To simplify further calculations we do not consider the deformation interaction between the electrons and the lattice. Using the generally accepted assumption the effect of the electrons on the lattice arises due to a self-consistent electric field, which occurs under deformation. Also, the deformation of the lattice causes the occurence of an additional non-uniform magnetic field $`𝐛(𝐫).`$ The presence of these fields leads to a redistribution of the electron density $`N.`$ The local change in the electronic density $`\delta N(𝐫)`$ equals:
$`\delta N(𝐫)`$ $`=`$ $`{\displaystyle \frac{N}{\zeta }}e\phi (𝐫)+{\displaystyle \frac{N}{𝐁}}𝐛(𝐫)`$ (3)
$``$ $`N_\zeta ^{}\left[e\phi (𝐫)+{\displaystyle \frac{\zeta }{𝐁}}𝐛(𝐫)\right].`$
Here, $`e`$ is the absolute value of the electron charge, $`𝐁`$ is the external magnetic field.
The magnetic field $`𝐛(𝐫)`$ satisfies the equation:
curl b(r) $`=`$ $`4\pi \text{curl }\delta \text{M(r)}`$ (4)
$`=`$ $`4\pi \text{curl}\left[{\displaystyle \frac{𝐌}{\zeta }}e\phi (𝐫)+{\displaystyle \frac{𝐌}{𝐁}}𝐛(𝐫)\right]\text{div }\text{b(r)}`$
$`=`$ $`0.`$
Here, $`𝐌`$ is the magnetization vector; $`\zeta `$ is the chemical potential of charge carriers; $`\phi (𝐫)`$ is the potential of the electric field, arising due to the deformation; $`N_\zeta `$ is the “bare” density of states (DOS) of quasiparticles on the Fermi surface of the layered conductor:
$$N_\zeta =\underset{\nu }{}\frac{f_\nu }{E_\nu }.$$
(5)
$`f_\nu `$ is the Fermi distribution function for the quasiparticles with energies $`E_\nu .`$
Within the framework of quantum phenomenological Fermi-liquid theory the quantity $`N_\zeta ^{}`$ is the DOS of the quasiparticles renormalized due to the Fermi-liquid interaction among them (see Ref. 13 ):
$$N_\zeta ^{}=\underset{\nu \nu ^{}}{}\frac{f_\nu f_\nu ^{}}{E_\nu E_\nu ^{}}n_{\nu \nu ^{}}^{}(\text{q})n_{\nu ^{}\nu }(\text{q})|_{q=0}.$$
(6)
and the renormalized operator of electron density $`n_{\nu \nu ^{}}^{}(𝐪)`$ is connected with the “bare”operator $`n_{\nu \nu ^{}}(𝐪)`$ by the relation:
$$n_{\nu \nu ^{}}^{}(𝐪)=n_{\nu \nu ^{}}(𝐪)+\underset{\nu _1\nu _2}{}\frac{f_{\nu _1}f_{\nu _2}}{E_{\nu _1}E_{\nu _2}}F_{\nu \nu ^{}}^{\nu _1\nu _2}n_{\nu _1\nu _2}^{}(𝐪)$$
(7)
where $`F_{\nu \nu ^{}}^{\nu _1\nu _2}`$ are the matrix elements of the Fermi-liquid kernel \[See Eq. (13) below\].
The relations (3), (4) have to be complemented by the condition of electrical neutrality of the system:
$$\delta N(𝐫)+eN\text{div }\text{u(r) }=0,$$
(8)
where $`𝐮(𝐫)`$ is the lattice displacement vector, $`N`$ is the electron density. It follows from equations (3), (4) and (8) that:
$$\text{curl}\{(14\pi \chi )𝐛(𝐫)\}=4\pi \text{curl}\left\{\frac{𝐌}{\zeta }\frac{N}{N_\zeta ^{}}\text{div }\text{u(r)}\right\}.$$
(9)
The set of simultaneous equations (3), (4) and (8) was first presented in Refs. 14 and 15 . It allows us to exclude $`𝐛(𝐫)`$ and to express the potential $`\phi (𝐫)`$ in terms of the lattice displacement vector. As a result we can derive the expression for the electron force $`𝐅(𝐫)`$ acting upon the lattice under its displacement by the vector $`𝐮(𝐫)`$. For the conductor with axially symmetric FS in the magnetic field $`\underset{¯}{\text{B}}`$ directed along the symmetry axis, the force $`𝐅(𝐫)`$ equals:
$$𝐅(𝐫)=\lambda _0𝐛_0(𝐛_0(𝐮(𝐫)))+\lambda [𝐛_0\times [(𝐮(𝐫))\times 𝐛_0]].$$
(10)
Here $`𝐛_0`$ is the unit vector directed along $`𝐁`$. We assume that the magnetic field $`𝐁`$ is directed along the normal to the layers. Then $`𝐛_0`$ coincides with the normal vector to the layers $`𝐧.`$ It follows from Eq. (10), that under longitudinal (parallel to $`𝐛_0`$ or perpendicular to layers) direction of gradients of the deformation tensor the electron contribution to the corresponding elastic modulus is equal to $`\lambda _0,`$ and under their direction across $`𝐛_0`$ (in the plane of layers) it equals $`\lambda .`$ The elastic constants $`\lambda _0`$ and $`\lambda `$ are described by the expressions:
$`\lambda _0`$ $`=`$ $`{\displaystyle \frac{N^2}{N_\zeta ^{}}},`$ (11)
$`\lambda `$ $`=`$ $`\lambda _0\left(1+{\displaystyle \frac{4\pi \chi _\zeta }{14\pi \chi _{||}}}\right).`$ (12)
Here, $`\chi _{||}`$ is the longitudinal part of the magnetic susceptibility; $`\chi _\zeta =(M_z/\zeta )(\zeta /B).`$ More careful treatment carried out before including the deformation interaction (see e.g. Ref. 16 ), do not give a qualitative change of results. If we will take into account such interaction, it will lead us to a replacement of the electron density $`N`$ in the expression (11) for $`\lambda _0`$ by the quantity $`N(1+a),`$ where the value of the dimensionless constant $`a`$ depends on the corresponding component of the tensor of deformation potential averaged over the FS. The constant $`N`$ in Eq. (12) also has to be replaced by $`N(1+a^{})`$ where $`a^{}`$ is a dimensionless constant depending on the averaged over the FS deformation potential.
As follows from Eq. (11), in our case, when we neglect the deformation interaction the quantity $`\lambda _0`$ coincides with the compression modulus of the electron liquid. The structure of the quantity $`\lambda `$ is more complicated. Besides the contribution, connected with the electron compression itself, $`\lambda `$ also contains the contribution arising due to the magnetostriction (the second term in the expression enclosed in the brackets). It is shown below that at low temperatures when $`\theta 1(\theta =2\pi ^2T/\mathrm{}\mathrm{\Omega };T`$ is the temperature expressed in energy units, $`\mathrm{\Omega }`$ is the cyclotron frequency of the charge carriers) there is a significant distinction between the oscillating corrections to the elastic moduli $`\lambda _0`$ and $`\lambda `$ arising due to the magnetostriction. In the considered case $`(𝐁||𝐧)`$ the velocity of the longitudinal sound, propagating along the normal to the layers is proportional to the square root of $`\lambda _0,`$ and when it propagates in the plane of layers – to the square root of $`\lambda .`$ Therefore, the distinction between $`\lambda _0`$ and $`\lambda `$ due to the magnetostriction will be displayed in the quantum oscillations of the velocity of sound. Above the low temperature range $`(\theta 1)`$ the contribution connected with the magnetostriction becomes negligible and the distinction between $`\lambda _0`$ and $`\lambda `$ is smoothed.
## III 3. low temperature quantum oscillations of the DOS
Here, we calculate the oscillating corrections to the renormalized DOS $`N_\zeta ^{}`$ which is defined by Eq. (6). The calculation is carried out for an axially symmetric FS of an arbitrary shape and this needs a special consideration because the phenomenological Fermi-liquid theory is well developed only for isotropic model of a metal. We start from the well known expressions for matrix elements of the Fermi-liquid kernel $`F_{\nu \nu ^{}}^{\nu _1\nu _2}`$ 17 :
$$F_{\nu \nu ^{}}^{\nu _1\nu _2}=\phi _{\alpha \alpha ^{}}^{\alpha _1\alpha _2}\delta _{\sigma \sigma ^{}}\delta _{\sigma _1\sigma _2}+4\psi _{\alpha \alpha ^{}}^{\alpha _1\alpha _2}(𝐬_{\sigma \sigma ^{}}𝐬_{\sigma _1\sigma _2}),$$
(13)
where $`\alpha `$ is the set of the orbital quantum numbers, $`\sigma `$ is the spin number, $`𝐬`$ is the spin operator.
The FS under consideration is axialy symmetric. Therefore Fermi-liquid functions $`\phi (𝐩,𝐩^{})`$ and $`\psi (𝐩,𝐩^{})`$ on the FS depend only on the cosine of the angle $`\theta `$ between the vectors $`𝐩_{}`$ and $`𝐩_{}^{}`$ and on the projections of the quasimomenta on the symmetry axis $`p_z,p_z\mathrm{´}.`$ We can represent the function $`\phi (𝐩,𝐩^{})`$ as a sum of even and odd contributions in $`\mathrm{cos}\theta :`$
$$\phi (𝐩,𝐩^{})=\phi _0(p_z,p_z^{},\mathrm{cos}\theta )+\mathrm{cos}\theta \phi _1(p_z,p_z^{},\mathrm{cos}\theta ),$$
(14)
where $`\phi _0(p_z,p_z^{},\mathrm{cos}\theta )`$ and $`\phi _1(p_z,p_z^{},\mathrm{cos}\theta )`$ are even function in $`\mathrm{cos}\theta .`$ By virtue of the FS invariance under replacements $`𝐩𝐩`$ and $`𝐩^{}𝐩^{},`$ the functions $`\phi _0(p_z,p_z^{},\mathrm{cos}\theta )`$ and $`\phi _1(p_z,p_z^{},\mathrm{cos}\theta )`$ should not vary under simultaneous change of sign of $`p_z`$ and $`p_z^{}.`$ It enables us to separate each of the mentioned functions into the parts which are even and odd in $`p_zp_z^{}`$ and to rewrite Eq. (14) in the form:
$$\phi (𝐩,𝐩^{})=\phi _{00}+p_zp_z^{}\phi _{01}+\mathrm{cos}\theta (\phi _{10}+p_zp_z^{}\phi _{11}).$$
(15)
Here, the functions $`\phi _{00},\phi _{01},\phi _{10}`$ and $`\phi _{11}`$ depend on $`p_z,p_z^{}`$ and $`\mathrm{cos}\theta `$ and are even functions in all their arguments. We can write a similar expression for $`\psi (𝐩,𝐩^{}).`$
To proceed we expand functions included into Eq. (15) in Fourier series in the angles $`\mathrm{\Phi }`$ and $`\mathrm{\Phi }^{}.`$ The functions are even in $`\mathrm{cos}\theta ,`$ therefore the expansions have the form:
$$\phi _{ik}(p_z,p_z^{},\mathrm{cos}\theta )=\underset{s=\mathrm{}}{\overset{\mathrm{}}{}}\phi _{ik}^{|2s|}(p_z,p_z^{})e^{2is(\mathrm{\Phi }\mathrm{\Phi }^{})}.$$
(16)
Here $`i,k=0,1.`$ Coefficients $`\phi _{ik}^{|2s|}(p_z,p_z^{}),`$ in the expansions (16) are even functions in both arguments. Besides, they should not change when the arguments are interchanged. We can simplify further consideration assuming that:
$$\phi _{ik}^{|2s|}(p_z,p_z^{})=\frac{1}{g_0}P_{ik}^{|2s|}(p_z)P_{ik}^{|2s|}(p_z^{}),$$
(17)
where both factors are even functions, and $`g_0`$ is the “bare” DOS of the quasiparticles on the FS in the absence of the magnetic field $`𝐁.`$
The expansions of the functions $`\psi _{ik}^{|2s|}(p_z,p_z^{},\mathrm{cos}\theta )`$ are similar to Eq. (16). The coefficients in these expansions also can be presented as:
$$\psi _{ik}^{|2s|}(p_z,p_z^{})=\frac{1}{g_0}Q_{ik}^{|2s|}(p_z)Q_{ik}^{|2s|}(p_z^{}),$$
(18)
To calculate the quantity $`N_\zeta ^{}`$ under the conditions of the semiclassical quantization $`(\gamma 1`$ where $`\gamma ^2=2\zeta /\mathrm{}\mathrm{\Omega })`$ we can replace the matrix elements $`\phi _{\alpha \alpha ^{}}^{\alpha _1\alpha _2}`$ and $`\psi _{\alpha \alpha ^{}}^{\alpha _1\alpha _2}`$ by their semiclassical analogs. The analog for $`\phi _{\alpha \alpha ^{}}^{\alpha _1\alpha _2}`$ is the Fourier component in the expansion of the function:
$`\phi (𝐪;𝐩,𝐩^{})`$ $`=`$ $`\mathrm{exp}\left[{\displaystyle \frac{icq_yp_{}}{eB}}\mathrm{sin}\mathrm{\Phi }\right]\phi (𝐩,𝐩^{})`$ (19)
$`\times `$ $`\mathrm{exp}\left[{\displaystyle \frac{icq_yp_{}^{}}{eB}}\mathrm{sin}\mathrm{\Phi }^{}\right],`$
in the Fourier series in the angles $`\mathrm{\Phi }`$ and $`\mathrm{\Phi }^{}.`$ Here, $`\phi (𝐩,𝐩^{})`$ is the semiclassical Fermi-liquid function. With the help of the known identity for the Bessel functions:
$$\mathrm{exp}(\pm iz\mathrm{sin}\mathrm{\Phi })=\underset{m=\mathrm{}}{\overset{\mathrm{}}{}}J_m(z)\mathrm{exp}(\pm im\mathrm{\Phi }),$$
(20)
we arrive at the result:
$`\phi _{ll^{}}(𝐪;𝐩,𝐩^{})`$ $`=`$ $`{\displaystyle \frac{1}{4\pi ^2}}{\displaystyle \underset{m=\mathrm{}}{\overset{\mathrm{}}{}}}J_m\left({\displaystyle \frac{cq_yp_{}}{eB}}\right)`$ (21)
$`\times `$ $`{\displaystyle \underset{m^{}=\mathrm{}}{\overset{\mathrm{}}{}}}J_m^{}({\displaystyle \frac{cq_yp_{}^{}}{eB}}){\displaystyle _0^{2\pi }}d\mathrm{\Phi }\mathrm{exp}[i(ml]\mathrm{\Phi }]`$
$`\times `$ $`{\displaystyle _0^{2\pi }}𝑑\mathrm{\Phi }^{}\mathrm{exp}[i(m^{}l^{})\mathrm{\Phi }^{}]\phi (𝐩,𝐩^{}).`$
Substituting the expression (14) for $`\phi (𝐩,𝐩^{})`$ into Eq. (21) and taking into account the expansions (16) we arrive at the final result:
$`\phi _{ll^{}}(𝐪;𝐩,𝐩^{})`$ (22)
$`=`$ $`{\displaystyle \frac{1}{g_0}}{\displaystyle \underset{s=\mathrm{}}{\overset{\mathrm{}}{}}}\{[P_{00}^{|2s|}(p_z)P_{00}^{|2s|}(p_z^{})`$
$`+`$ $`p_zp_z^{}P_{01}^{|2s|}(p_z)P_{01}^{|2s|}(p_z^{})]`$
$`\times `$ $`J_{l2s}\left({\displaystyle \frac{cq_yp_{}}{eB}}\right)J_{l^{}2s}\left({\displaystyle \frac{cq_yp_{}^{}}{eB}}\right)`$
$`+`$ $`\left[P_{10}^{|2s|}(p_z)P_{10}^{|2s|}(p_z^{})+p_zp_z^{}P_{11}^{|2s|}(p_z)P_{11}^{|2s|}(p_z^{})\right]`$
$`\times `$ $`[J_{l2s1}\left({\displaystyle \frac{cq_yp_{}}{eB}}\right)J_{l^{}2s1}\left({\displaystyle \frac{cq_yp_{}^{}}{eB}}\right)`$
$`+`$ $`J_{l2s+1}\left({\displaystyle \frac{cq_yp_{}}{eB}}\right)J_{l^{}2s+1}\left({\displaystyle \frac{cq_yp_{}^{}}{eB}}\right)]\}.`$
In the expansions of $`\psi _{ll^{}}(𝐪;𝐩,𝐩^{})`$ the functions $`P_{ik}^{|2s|}(p_z)`$ and $`P_{ik}^{|2s|}(p_z^{})`$ in (22) have to be replaced by the functions $`Q_{ik}^{|2s|}(p_z)`$ and $`Q_{ik}^{|2s|}(p_z^{}).`$ The derived expressions (22) enable us to study any Fermi-liquid effects for conductors whose FSs are axially symmetric.
Here, we consider the most interesting case when the gradients of the deformation of the lattice are directed across $`𝐛_0(𝐪𝐛_0).`$ Introducing the notations $`2N=n+n^{};2N^{}=n_1+n_2;l=nn^{};l^{}=n_1n_2`$ we can write the expansions for the matrix elements $`\phi (𝐪;Nlp_z;N^{}l^{}p_z^{})`$ and $`\psi (𝐪;Nlp_z;N^{}l^{}p_z^{}).`$ We arrive at these expansions replacing the quantities $`p_{},p_{}^{}`$ in the arguments of Bessel functions in Eq. (22) by their quantum analogs $`\mathrm{}\sqrt{2N+1}`$ and $`\mathrm{}\sqrt{2N^{}+1}`$.
The functions $`P_{ik}^{|2s|}(p_z)`$ and $`Q_{ik}^{|2s|}(p_z)`$ are even, and the matrix elements $`n_{\nu \nu ^{}}(𝐪)`$ are diagonal in spin number $`\sigma .`$ Owing to this, one can take into account only the first term in the first square brackets in Eq. (22) in calculation of the quantity $`N_\zeta ^{}.`$ The contributions from the other terms will be equal to zero. Thus in the following calculations we can assume that $`\phi (𝐪;Nlp_z;N^{}l^{}p_z^{})`$ and $`\psi (𝐪;Nlp_z;N^{}l^{}p_z^{})`$ have the form
$`\phi (𝐪;Nlp_z;N^{}l^{}p_z^{})`$ $`=`$ $`{\displaystyle \frac{1}{g_0}}{\displaystyle \underset{s=\mathrm{}}{\overset{\mathrm{}}{}}}P_{00}^{|2s|}(p_z)P_{00}^{|2s|}(p_z^{})`$
$`\times `$ $`J_{l2s}\left({\displaystyle \frac{\mathrm{}cq}{eB}}\sqrt{2N+1}\right)`$
$`\times `$ $`J_{l^{}2s}\left({\displaystyle \frac{\mathrm{}cq}{eB}}\sqrt{2N^{}+1}\right);`$
$`\psi (𝐪;Nlp_z;N^{}l^{}p_z^{})`$ $`=`$ $`{\displaystyle \frac{1}{g_0}}{\displaystyle \underset{s=\mathrm{}}{\overset{\mathrm{}}{}}}Q_{00}^{|2s|}(p_z)Q_{00}^{|2s|}(p_z^{})`$ (23)
$`\times `$ $`J_{l2s}\left({\displaystyle \frac{\mathrm{}cq}{eB}}\sqrt{2N+1}\right)`$
$`\times `$ $`J_{l^{}2s}\left({\displaystyle \frac{\mathrm{}cq}{eB}}\sqrt{2N^{}+1}\right).`$
Further we omit the lower indices of the functions $`P_{00}^{|2s|}(p_z)`$ and $`Q_{00}^{|2s|}(p_z)`$ for simplicity.
Under conditions of the semiclassical quantization, when $`\gamma 1,`$ we have for the arbitrary single particle operator $`\mathrm{\Phi }_{\nu \nu ^{}}`$ (see Ref. 13 ):
$`{\displaystyle \underset{\nu \nu ^{}}{}}{\displaystyle \frac{f_\nu f_\nu ^{}}{E_\nu E_\nu ^{}}}\mathrm{\Phi }_{\nu \nu ^{}}`$ (24)
$`=`$ $`{\displaystyle \frac{1}{4\pi ^2\mathrm{}^3}}{\displaystyle \underset{l}{}}{\displaystyle \underset{\sigma \sigma ^{}}{}}{\displaystyle }dp_z[m_{}(p_z)\mathrm{\Phi }_l^{\sigma \sigma ^{}}(p_z)`$
$`+`$ $`\delta _{l_0}\delta _{\sigma \sigma ^{}}{\displaystyle \underset{m}{}}(\mathrm{\Delta }_m+\sigma \mathrm{\Delta }_m^{})m_{}(p_m)\mathrm{\Phi }_l^{\sigma \sigma ^{}}(p_m)].`$
Here, $`m_{}(p_z)`$ is the cyclotron mass which is assumed to be a constant in further calculations; $`p_m`$ is the value of the longitudinal component of the quasimomentum, corresponding to the $`m`$th external cross section of the FS. The functions $`\mathrm{\Delta }`$ and $`\mathrm{\Delta }^{}`$ describe oscillating corrections which appear in quantizing magnetic fields. Their form depends on the particular character of the energy-momentum relation for the quasiparticles. We will analyze these quantum oscillating corrections below within the framework of our model for the electron energy spectrum in organic metals. Integration with respect to $`p_z`$ in Eq. (24) is carried out within the limits determined by the shape and size of the FS.
Using the asymptotic formula (24), and the expansions (23) we can obtain the expression for the renormalized DOS:
$$N_\zeta ^{}=g_0\left[1\alpha _0+\frac{(1\overline{\alpha }_0)^2[\mathrm{\Delta }+\beta (\mathrm{\Delta }^2\mathrm{\Delta }^2)]}{1+(\alpha +\beta )\mathrm{\Delta }+\alpha \beta (\mathrm{\Delta }^2\mathrm{\Delta }^2)}\right].$$
(25)
Here:
$`\alpha _0`$ $`=`$ $`{\displaystyle \frac{A_0}{1+\stackrel{~}{A}_0}};\overline{\alpha }_0={\displaystyle \frac{\overline{A}_0}{1+\stackrel{~}{A}_0}};`$
$`\alpha `$ $`=`$ $`{\displaystyle \frac{A_0^{}}{1+\stackrel{~}{A}_0}};\beta ={\displaystyle \frac{B_0^{}}{1+\stackrel{~}{B}_0}};`$
$`A_0`$ $`=`$ $`{\displaystyle \frac{m_{}}{2\pi ^2\mathrm{}^3g_0}}\left({\displaystyle P^0(p_z)𝑑p_z}\right)^2;`$
$`\overline{A}_0`$ $`=`$ $`{\displaystyle \frac{P^0(0)m_{}}{2\pi ^2\mathrm{}^3g_0}}{\displaystyle P^0(p_z)𝑑p_z},`$
$`\stackrel{~}{A}_0`$ $`=`$ $`{\displaystyle \frac{m_{}}{2\pi ^2\mathrm{}^3g_0}}{\displaystyle [P^0(p_z)]^2𝑑p_z}.`$
$`\stackrel{~}{B}_0`$ $`=`$ $`{\displaystyle \frac{m_{}}{2\pi ^2\mathrm{}^3g_0}}{\displaystyle [Q^0(p_z)]^2𝑑p_z};`$
$`A_0^{}`$ $`=`$ $`[P^0(0)]^2;B_0^{}=[Q^0(0)]^2.`$ (26)
Neglecting terms proportional to the differencies $`\mathrm{\Delta }^2\mathrm{\Delta }^2`$ in the numerator and the denominator of the Eq. (26) we can simplify the expression for the renormalized density of states:
$$N_\zeta ^{}=g_0\left[1\alpha _0+\frac{(1\overline{\alpha }_0)^2\mathrm{\Delta }}{1+(\alpha +\beta )\mathrm{\Delta }}\right].$$
(27)
The expression (27) is easily generalized to cover the case when the Fermi surface has several extremal cross sections. In this case we obtain
$$N_\zeta ^{}=g_0[1\alpha _0+\frac{\underset{m}{}(1\overline{\alpha }_m)^2\mathrm{\Delta }_m}{1+\underset{m}{}(\alpha _m+\beta _m)\mathrm{\Delta }_m}]g_0(1\alpha _0+K).$$
(28)
Here, summation over $`m`$ is carried out over the extremal cross sections:
$`\alpha _m={\displaystyle \frac{A_m^{}}{1+\stackrel{~}{A}_0}};\overline{\alpha }_m={\displaystyle \frac{\overline{A}_m}{1+\stackrel{~}{A}_0}};\beta _m={\displaystyle \frac{B_m^{}}{1+\stackrel{~}{B}_0}};`$
$`A_m^{}=[P^0(p_m)]^2;B_m^{}=[Q^0(p_m)]^2;`$
$`\overline{A}_m={\displaystyle \frac{m_{}}{2\pi ^2\mathrm{}^3g_0}}P^0(p_m){\displaystyle P^0(p_z)𝑑p_z}.`$ (29)
Neglecting all terms arising due to the interaction among quasiparticles, we arrive at the result for the “bare” DOS:
$$N_\zeta =g_0(1+\mathrm{\Delta }).$$
(30)
The oscillating term $`\mathrm{\Delta }`$ describes quantum oscillations of the “bare” DOS. This term is small compared to unity when the temperature is moderately low $`(\theta 1).`$ When $`\theta 1`$ the magnitude of the function $`\mathrm{\Delta }`$ can reach values of the order of unity. Comparison of Eqs. (28) and (30) shows that there are pronounced distinctions between the “bare” DOS $`N_\zeta `$ and the DOS $`N_\zeta ^{}`$ renormalized due to the interaction among quasiparticles within low temperature range. Under $`\theta 1`$ the amplitude and the form of the function $`K`$ which describes the quantum oscillations of $`N_\zeta ^{}`$ can differ significantly from the corresponding characteristics of the function $`\mathrm{\Delta }.`$ In particular, the denominator in Eq. (22) can turn zero at the peaks of the oscillations. Thus, the function $`K`$ has singularities caused by the Fermi-liquid interaction.
The specific features of the low temperature quantum oscillations in renormalized density of states have to be manifested in oscillations of the thermodynamic observables. We consider, for example, the oscillating contribution to the magnetic susceptibility for a uniform field $`𝐁.`$ To simplify the following calculations we assume that the unique extreme cross-section of the axisymmetric FS is located at $`p_z=0.`$ We remark here that possible singularities in the longitudinal magnetic susceptibility near the peaks of quantum oscillations of the renormalized DOS were first analyzed in Ref. 18 for an isotropic electron liquid (see also Ref. 19 ). Generalizing the corresponding result of Ref. 18 , we can show that in the case of the axisymmetric FS, $`\chi _{||}`$ is proportional to the oscillating function $`K^{}:`$
$`K^{}`$ $`=`$ $`{\displaystyle \frac{K}{1+[1\overline{\alpha }_0(0)]K}}`$ (31)
$`=`$ $`{\displaystyle \frac{\mathrm{\Delta }(0)}{1+[1+\alpha _0(0)\overline{\alpha }_0(0)+\beta _0(0)]\mathrm{\Delta }(0)}}.`$
Within the low temperature range the amplitude of the oscillations in $`\mathrm{\Delta }(0)`$ might increase to such an extend that it reaches the values of the order of unity. Under these conditions the form of the de Haas-van Alphen oscillations is determined by the denominator in Eq. (31) to a great extent, and it depends critically on the Fermi-liquid interaction. When $`[1+\alpha _0(0)\overline{\alpha }_0(0)+\beta _0(0)]<0`$ the denominator in Eq. (31) can go to zero at the peaks of the oscillations and the susceptibility $`\chi _{||}`$ would increase beyond all bounds. It would lead to the magnetic instability $`(\chi _{||}<1/4\pi )`$ at the peaks of the oscillations (see Fig. 1). Thus, the low temperature analog for the Shoenberg effect can be observed. It follows from the results of Ref. 18 as well as from our expression (31) that the magnetic susceptibility $`\chi _{||}`$ of noninteracting 3D electron gas can exhibit a series of maxima of a finite magnitude when temperature goes to zero. G. Montambaux et al. 11 came to a similar conclusion irrespective of the earlier results of Ref. 18 .
These singularities of the renormalized DOS $`N_\zeta ^{}`$ originating due to the Fermi-liquid interaction can give rise to specific features in the quantum oscillations of thermodynamic observables in layered conductors. However, their manifestations would depend on the local geometry of the FS near its extreme cross sections. To analyze these anomalies in more detail we introduce a concrete model of the FS for layered organic metals which allows us to consider a broad class of the FSs of various profiles.
## IV 4. model and results
The usual approach in theoretical studies of electron properties of layered conductors is to keep only the first few terms in the sum over $`k`$ in Eq. (2). As a rule (see Refs. 20 ; 21 ; 22 , only the first term is taken into account. This corresponds to results obtained in the tight-binding approximation. Here, we use a different approach to describe the electron energy spectrum of the charge carriers in layered conductors, whose FS is defined by the equation
$$E(𝐩)=\frac{p_{}^2}{2m_{}}\eta v_0p_0E\left(\frac{p_z}{p_0}\right),$$
(32)
where $`v_0=(2\zeta /m_{})^{1/2},E(p_z/p_0)`$ is an even function periodic in its argument $`p_z/p_0`$ with a period equal to 2, and $`\eta `$ is a dimensionless parameter characterizing the rate of rippling of the FS. The quantity $`\eta v_0p_0E(p_z/P_0)`$ is the sum of the trigonometric series in (2). By selecting the type of this function we can obtain FS’s shaped as pinched cylinders with different profiles. This approach provides broad possibilities in analyzing the effect of the shape of the Fermi surface on observed characteristics of layered conductors.
Let us assume that the function $`E(p_z/p_0)`$ in the interval $`p_0p_zp_0`$ is described by the expression
$$E\left(\frac{p_z}{p_0}\right)=\frac{1}{rl}\left[1\left|\frac{p_z}{p_0}\right|^l\right]^r,$$
(33)
where the parameters $`l`$ and $`r`$ take values more than unity. The model specified by (32) and (33) makes it possible to describe a broad class of FSs.
The Gaussian curvature of the FS is described by the expression
$$\kappa (p_z)=\frac{m_{}^2[v_z^2+(p_{}^2/m_{})v_z/p_z]}{(p_{}^2+m_{}^2v_z^2)^2}.$$
(34)
At $`l=r=2`$ the curvature of the FS at its sections by the planes $`p_z=0`$ and $`p_z=\pm p_0`$ equals:
$`\kappa (0)`$ $`=`$ $`{\displaystyle \frac{\delta S}{S_{\mathrm{max}}}}{\displaystyle \frac{1}{p_0^2}};`$ (35)
$`\kappa (\pm p_0)`$ $`=`$ $`{\displaystyle \frac{2\delta S}{S_{\mathrm{min}}}}{\displaystyle \frac{1}{p_0^2}},`$ (36)
where $`S_{\mathrm{max}}`$ and $`S_{\mathrm{min}}`$ are the maximum and minimum sectional areas of the FS: $`S_{\mathrm{max}}=S(0),S_{\mathrm{min}}=S(\pm p_0),`$ and $`\delta S=S_{\mathrm{max}}S_{\mathrm{min}}=(\pi /2)m_{}\eta v_0p_0.`$ Thus, if the FS remains a pinched cylinder $`(\eta 0),`$ its curvature at all points of the sections with extremal diameters is finite. Similar results are obtained if the tight-binding approximation is used to describe the electron energy spectrum.
For $`r2`$ and $`l=2`$ the curvature of the FS near $`p_z\pm p_0`$ remains finite and $`\kappa (0)`$ is still described by (35). However, the asymptotic behavior of the curvature of the FS near $`p_z=\pm p_0`$ is different:
$$\kappa (p_z)=2(r1)\frac{\delta S}{S_{\mathrm{min}}}\frac{1}{p_0^2}\left[1\left(\frac{p_z}{p_0}\right)^2\right]^{r2}.$$
(37)
Thus, for $`1<r<2`$ the curvature of the FS has singularities in these sections. For $`p_z=\pm p_0`$ the curvature $`\kappa (p_z)`$ vanishes at $`p_z=\pm p_0.`$ The corresponding sections of the Fermi surface are lines of parabolic points. The larger the value of the parameter $`r,`$ the more the FS near these sections resembles a cylinder.
The anomalies in the curvature of the FS near $`p_z=0`$ can be described by the model (32) and (33) with $`r=2`$ and $`l2.`$ Here the curvature of the FS near $`p_z=0`$ is described by the asymptotic expression
$$\kappa (p_z)=(l1)\frac{\delta S}{S_{\mathrm{max}}}\frac{1}{p_0^2}\left|\frac{p_z}{p_0}\right|^{l2}.$$
(38)
For $`1<l<2`$ the curvature of the FS has a singularity at $`p_z=0;`$ if $`l>2`$ the Fermi surface near $`p_z=0`$ resembles a cylinder, and the larger the value of $`l`$ the closer the resemblance. Finally, if $`r2`$ and $`l2,`$ we have a surface in the form of a pinched cylinder with curvature singularities in all the sections with extremal diameters. The profiles of the Fermi surface described by (32) and (33) are depicted schematically in Fig. 2.
We shall use our model to obtain the expressions for the oscillating function $`\mathrm{\Delta }.`$ The form of the function $`\mathrm{\Delta }`$ depends importanty on the characteristic properties of the electron energy spectrum. Assuming that FS curvature on the effective strips is nonzero $`(r=l=2)`$ we can obtain:
$`\mathrm{\Delta }`$ $`=`$ $`\mathrm{\Delta }(0)+\mathrm{\Delta }(p_0)={\displaystyle \frac{1}{\gamma }}\sqrt{{\displaystyle \frac{S_0}{\delta S}}}\{{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^n}{\sqrt{n}}}\psi _n(\theta )`$ (39)
$`\times `$ $`\mathrm{cos}\left({\displaystyle \frac{ncS_{\mathrm{max}}}{\mathrm{}eB}}{\displaystyle \frac{\pi }{4}}\right)\mathrm{cos}\left(\pi n{\displaystyle \frac{\mathrm{\Omega }_0}{\mathrm{\Omega }}}\right)+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^n}{\sqrt{n}}}\psi _n(\theta )`$
$`\times `$ $`\mathrm{cos}({\displaystyle \frac{ncS_{\mathrm{min}}}{\mathrm{}eB}}+{\displaystyle \frac{\pi }{4}})\mathrm{cos}\left(\pi n{\displaystyle \frac{\mathrm{\Omega }_0}{\mathrm{\Omega }}}\right)\}.`$
Here, $`S_0=2\pi m_{}\zeta `$ is the area of the FS cross-section in the limit $`\eta 0,`$ when the FS turns a pure cylinder; $`\mathrm{}\mathrm{\Omega }_0`$ is the spin splitting energy.
In the case when there is an anomaly in the FS curvature at $`p_z=\pm p_0(l=2;r2)`$ the contribution to the oscillating function $`\mathrm{\Delta }`$ from the vicinities of corresponding cross-sections assumes a form:
$`\mathrm{\Delta }(p_0)`$ $`=`$ $`\left({\displaystyle \frac{rS_0}{2\gamma ^2\delta S}}\right)^{1/r}{\displaystyle \frac{\mathrm{\Gamma }(1/r)}{r}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^n}{n^{1/r}}}\psi _n(\theta )`$ (40)
$`\times `$ $`\mathrm{cos}\left({\displaystyle \frac{ncS_{\mathrm{min}}}{\mathrm{}eB}}+{\displaystyle \frac{\pi }{2r}}\right)\mathrm{cos}\left(\pi n{\displaystyle \frac{\mathrm{\Omega }_0}{\mathrm{\Omega }}}\right),`$
where $`\mathrm{\Gamma }`$ is the gamma function. Under this condition the first term in Eq. (40) retains its form. Correspondingly when the local geometry of the FS is characterized with the anomalous curvature at $`p_z=0(r=2;l2)`$ the oscillating function $`\mathrm{\Delta }(0)`$ is described by the expression:
$`\mathrm{\Delta }(0)`$ $`=`$ $`\left({\displaystyle \frac{lS_0}{2\gamma ^2\delta S}}\right)^{1/l}{\displaystyle \frac{\mathrm{\Gamma }(1/l)}{l}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^n}{n^{1/l}}}\psi _n(\theta )`$ (41)
$`\times `$ $`\mathrm{cos}\left({\displaystyle \frac{ncS_{\mathrm{max}}}{\mathrm{}|e|B}}{\displaystyle \frac{\pi }{2l}}\right)\mathrm{cos}\left(\pi n{\displaystyle \frac{\mathrm{\Omega }_0}{\mathrm{\Omega }}}\right),`$
At the same time the expression (40) for $`\mathrm{\Delta }(p_0)`$ remains valid.
When the effective strip on the FS centered at $`p_z=0`$ or $`p_z=\pm p_0`$ is close to a cylinder in shape, its contribution to the function $`\mathrm{\Delta }`$ has the same form as for the conductor with two-dimensional energy spectrum. For instance, if $`l1`$ the expression (41) can be transformed to the form
$$\mathrm{\Delta }(0)=\underset{n=1}{\overset{\mathrm{}}{}}(1)^n\psi _n(\theta )\mathrm{cos}\left(\frac{ncS_{\mathrm{max}}}{\mathrm{}|e|B}\right)\mathrm{cos}\left(\pi n\frac{\mathrm{\Omega }_0}{\mathrm{\Omega }}\right).$$
(42)
In the case when both parameters $`(r`$ and $`l)`$ are large compared to unity the expression for $`\mathrm{\Delta }`$ contains two terms of the form (42). Both terms represent the contributions from the quasicylindrical strips on the FS. These oscillating terms have to differ in period. The difference in their periods arises due to the distinction in the extremal cross sectional areas. It may be noticeable if the magnetic field is not too strong and the inequality $`\gamma ^2\delta S/S_01`$ is satisfied.
In a very strong magnetic field the inequality $`\gamma ^2\delta S/S_01`$ is satisfied. Under this condition the difference in form between the FS and the cylinder does not affect the oscillating function $`\mathrm{\Delta }:`$
$$\mathrm{\Delta }=2\underset{n=1}{\overset{\mathrm{}}{}}(1)^n\psi _n(\theta )\mathrm{cos}(\pi n\gamma ^2)\mathrm{cos}\left(\pi n\frac{\mathrm{\Omega }_0}{\mathrm{\Omega }}\right).$$
(43)
When a temperature is moderately low $`(\theta 1)`$ we can keep only the first term in the sum over $`\mathrm{`}\mathrm{`}n\mathrm{"}`$ in Eqs. (39)–(43). Thus, under $`\theta 1,`$ the dependence of the function $`\mathrm{\Delta }`$ on the inverse magnetic field has to be of a harmonic type. Because of the factor $`\mathrm{exp}(\theta ),`$ the amplitude of the oscillations of the function $`\mathrm{\Delta }`$ under $`\theta 1`$ is small compared to unity even in strong magnetic fields.
When $`\theta 1,`$ the form of the oscillations becomes much more complicated and their amplitude increases. To evaluate the amplitude of the oscillations under this condition one can use the asymptotic formulae following from the Euler-Maclaurin summation formula. As a result, under this condition we obtain the following estimation for the amplitude of the oscillations:
$`\mathrm{\Delta }_{\mathrm{max}}(p_0)`$ $``$ $`\left({\displaystyle \frac{S_0\theta }{\gamma ^2\delta S}}\right)^{1/r}{\displaystyle \frac{1}{\theta }},`$ (44)
$`\mathrm{\Delta }_{\mathrm{max}}(0)`$ $``$ $`\left({\displaystyle \frac{S_0\theta }{\gamma ^2\delta S}}\right)^{1/l}{\displaystyle \frac{1}{\theta }}.`$ (45)
Here, $`r`$ and $`l`$ are the dimensionless parameters whose values determine the local geometry of the FS in our model (33). In the case when the FS curvature remains finite and nonzero at $`p_z=0`$ and $`p_z=\pm p_0`$ we have:
$`\mathrm{\Delta }_{\mathrm{max}}(p_0)=\mathrm{\Delta }_{\mathrm{max}}(0)`$ $`=`$ $`{\displaystyle \frac{1}{\gamma }}\sqrt{{\displaystyle \frac{S_0}{\delta S}}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{2\sqrt{2n}}}\psi _n(\theta )`$ (46)
$``$ $`{\displaystyle \frac{3}{2}}{\displaystyle \frac{1}{\gamma \sqrt{\theta }}}\sqrt{{\displaystyle \frac{S_0}{\delta S}}}.`$
For the same values of the parameters $`\gamma ,\delta S/S_0`$ and $`\theta `$ the ratio of the amplitudes of oscillations described with the formulae (41) and (39) is of the order of $`(\gamma ^2\delta S/S_0\theta )^{(l2)/2l}.`$ Thus, under typical experimental conditions when $`\gamma ^2\delta S/S_01,`$ the oscillations associated with a quasicylindrical extremal cross-section of the FS have an amplitude much larger than the amplitude of the ordinary oscillations described by formula (39). If $`1<l<2`$ one of the principal radii of curvature vanishes at all points of the extremal cross-section of the Fermi surface at $`p_z=0.`$ In this case the amplitude of oscillations described by formula (39) is much smaller than the amplitude of the oscillating parts of the contributions from the FS ordinary extremal cross-sections.
When $`\theta 1`$ the oscillation amplitude may be of the order of unity. Charge carriers concentration in layered organic metals is far below than in ordinary metals. This generates much more favorable conditions for observation of the salient features in low temperature quantum oscillations in thermodynamic quantities. If we take for $`m_{}`$ and $`S_0`$ the values obtained in experimental study of the Fermi surface of the organic metal $`\beta \text{(ET)}_2\text{IBr}_2`$ (Ref. 23 ) $`(m_{}4.5m_0;S_01.26\times 10^{39}`$ g<sup>2</sup>cm$`{}_{}{}^{2}/`$s$`{}_{}{}^{2};m_0`$ is free electron mass) the condition $`(3/2)1/\gamma \sqrt{\theta }\sqrt{S_0/\delta S}1`$ in magnetic fields $`B200`$ kG and for $`\delta S/S_00.05`$ will be satisfied at temperatures of the order of one Kelvin.
Now we can return back to our expressions (11) and (12) for the elastic moduli and to proceed the analysis of their low temperature anomalies in layered conductors. For definiteness, we assume that the FS of the considered organic metal is nearly cylindrical in shape near $`p_z=0`$ and its curvature at $`p_z=\pm p_0`$ is finite and nonzero. Under these assumptions the oscillating correction $`\mathrm{\Delta }(0)`$ predominates and we can write:
$$\lambda _0=\frac{N^2}{N_\zeta ^{}}=\frac{N^2[1(1\overline{\alpha }_0(0))K^{}]}{(1\alpha _0)g_0},$$
(47)
$$\lambda =\lambda _0\left(1+\frac{4\pi \chi _\zeta }{14\pi \chi _{||}}\right)=\frac{N^2[1(1\overline{\alpha }_0(0))K^{\prime \prime }]}{(1\alpha _0)g_0}.$$
(48)
Here, the function $`K^{\prime \prime }`$ includes the oscillating term $`\mathrm{\Delta }(0),`$ therefore it oscillates itself in a strong magnetic field:
$$K^{\prime \prime }=\frac{\mathrm{\Delta }(0)}{1+[1+\alpha _0(0)\overline{\alpha }_0(0)+\beta _0(0)4\pi \chi _0\gamma ^4]\mathrm{\Delta }(0)}.$$
(49)
To obtain the expression for $`K^{\prime \prime }`$ we used the result (31) for the longitudinal part of the magnetic susceptibility $`\chi _{||};`$ the quantity $`\chi _0`$ is related to the Landau diamagnetic susceptibility. The latter equals $`\chi _0/3.`$
The sharpest difference between $`\lambda `$ and $`\lambda _0`$ is displayed at those values of the magnetic field $`𝐁,`$ at which the quantity $`14\pi \chi _{||}`$ goes to zero and the denominator in Eq. (49) turns zero at the peaks of the quantum oscillations. It makes the elastic constant $`\lambda `$ (or the electron contribution to the elastic modulus $`c_{11})`$ vanish near the oscillations maxima. Thus, there can arise lattice instability. It is connected with the instability in the magnetization and arises because of the magnetostriction. Under these conditions the Fermi-liquid correlation plays the major part. When $`4\pi \chi _0\gamma ^4<1`$ the instability can arise only under $`\alpha _0(0)\overline{\alpha }_0(0)+\beta _0(0)<0.`$
The lattice instability can, in principle, be displayed in the case, when the magnetostriction does not play a part (the strain gradient is directed along $`𝐁),`$ as the function $`K`$ \[see Eq. (27)\] in the low temperature region can also turn to infinity. Such instability was predicted in time in Refs. 23 and 24 . However, actually it is probably unattainable. The reason is that when $`\mathrm{\Delta }`$ increases the quantity $`1+[1+\alpha _0(0)\overline{\alpha }_0(0)+\beta _0(0)4\pi \kappa \gamma ^4]\mathrm{\Delta }(0)`$ reaches zero before the quantity $`1+[1+\alpha _0(0)\overline{\alpha }_0(0)+\beta _0(0)]\mathrm{\Delta }(0).`$ Hence, the jump of the magnetic state will take place before the constant $`\lambda _0`$ runs into zero.
## V 5. Concluding Remarks
The point of the present work is that the Fermi-liquid correlation among the quasiparticles can strongly affect their renormalized DOS. This was analyzed in detail for isotropic electron liquid 13 ; 14 ; 15 ; 18 . Here, we generalize this analysis to cover a broad class of layered conductors whose FSs are supposed to be axially symmetric. It is shown as a result of lengthy but straightforward calculations that $`N_\zeta ^{}`$ significantly differs from the “bare”DOS $`N_\zeta `$ under the condition that oscillating corrections $`\mathrm{\Delta }_m`$ which arise in the presence of the quantizing magnetic field are not small in magnitude compared to unity. Under this condition the oscillating part of $`N_\zeta ^{}`$ can have singularities in the peaks of quantum oscillations. These singularities occur due to the Fermi-liquid interaction. The singularities of $`N_\zeta ^{}`$ can cause specific features manifested in the oscillations of observables. In particular, they can lead to the singularities of the longitudinal magnetic susceptibility $`\chi _{||}`$ at the magnetic fields corresponding to the peaks of the quantum oscillations of DOS, which leads to the violation of the stability of the considered substance $`(\chi _{||}<1/4\pi ).`$ A similar effect was briefly discussed in Ref. 11 . Besides we obtain a possible singularity in the electron contribution to the elastic constant $`c_{22}`$ which is connected with the instability in the magnetization and can give rise to the lattice instability of a special kind.
These interesting phenomena can be available for observation in organic metals (including substances of the $`\alpha \text{(BEDT-TTF)}_2\text{MHg(SCH)}_4`$ group) because of the specific character of the energy spectra of the charge carriers in these substances. It is too early to draw any conclusions about the local features of the geometry of the FS of the majority of layered organic metals, since there is a lot to study in the electron energy spectra of such materials. It can be assumed, however, that here, as in ordinary metals, the FS contains quasicylindrical bands or sections with an anomalously large curvature. These features of the local geometry of the FS can be created (if they are absent) or enhanced by applying an agent that changes the shape of the constant-energy surfaces, e.g. by applying external pressure along the normal to the conducting planes.
The above analysis shows that the special features in the profile of the rippled cylinder, which is the main part of the FS of layered organic metals, can substantially influence quantum oscillations of thermodynamic observables in these materials. The model developed in this paper makes it possible to study in detail the observable manifestations of the local geometry of the FS of layered conductors. It resolves some of difficulties that emerge when we use the model of tightly bound electrons. For instance, the characteristic features of the observable properties of layered conductors, arising due to strong anisotropy of the electrical conductivity can be described and analyzed without passing to the limit $`\eta 0,`$ which corresponds to a conductor with a two-dimensional spectrum of the charge carriers.
The model specified by (32) and (33) enables us to carry out a detailed analysis of quantum oscillations of the elastic constants in organic metals and this leads us to the conclusion that the local geometry of their Fermi surfaces can provide favorable conditions for observation of the above described magnetic and lattice instabilities.
## VI Acknowledgments
We thank G.M. Zimbovsky for help with the manuscript.
## VII Appendix
Here we present in more detail calculation of the renormalized DOS $`N_\zeta ^{}.`$ We start from the expansion (23) for the matrix elements of the Fermi-liquid kernel. To proceed we use the notations:
$`I_s`$ $`=`$ $`{\displaystyle \underset{\nu \nu ^{}}{}}{\displaystyle \frac{f_\nu f_\nu ^{}}{E_\nu E_\nu ^{}}}P^{|2s|}(p_z)J_{l2s}\left({\displaystyle \frac{\mathrm{}cq}{eB}}\sqrt{2N+1}\right)`$ (50)
$`\times `$ $`n_{\nu \nu ^{}}^{}(q)\delta _{\sigma \sigma ^{}}\delta _{x_0;x_0^{}\mathrm{}cq/eB},`$
$`I_s^{}`$ $`=`$ $`{\displaystyle \underset{\nu \nu ^{}}{}}{\displaystyle \frac{f_\nu f_\nu ^{}}{E_\nu E_\nu ^{}}}Q^{|2s|}(p_z)J_{l2s}\left({\displaystyle \frac{\mathrm{}cq}{eB}}\sqrt{2N+1}\right)`$ (51)
$`\times `$ $`n_{\nu \nu ^{}}^{}(q)\delta _{\sigma \sigma ^{}}\delta _{x_0;x_0^{}\mathrm{}cq/eB},`$
$`g_{ss^{}}`$ $`=`$ $`{\displaystyle \frac{1}{g_0}}{\displaystyle \underset{\nu \nu ^{}}{}}{\displaystyle \frac{f_\nu f_\nu ^{}}{E_\nu E_\nu ^{}}}P^{|2s|}(p_z)P^{|2s^{}|}(p_z)`$ (52)
$`\times `$ $`J_{l2s}\left({\displaystyle \frac{\mathrm{}cq}{eB}}\sqrt{2N+1}\right)J_{l2s^{}}\left({\displaystyle \frac{\mathrm{}cq}{eB}}\sqrt{2N+1}\right)`$
$`\times `$ $`\delta _{\sigma \sigma ^{}}\delta _{x_0;x_0^{}\mathrm{}cq/eB},`$
$`q_{ss^{}}`$ $`=`$ $`{\displaystyle \frac{1}{g_0}}{\displaystyle \underset{\nu \nu ^{}}{}}{\displaystyle \frac{f_\nu f_\nu ^{}}{E_\nu E_\nu ^{}}}Q^{|2s|}(p_z)Q^{|2s^{}|}(p_z)`$ (53)
$`\times `$ $`J_{l2s}\left({\displaystyle \frac{\mathrm{}cq}{eB}}\sqrt{2N+1}\right)J_{l2s^{}}\left({\displaystyle \frac{\mathrm{}cq}{eB}}\sqrt{2N+1}\right)`$
$`\times `$ $`\delta _{\sigma \sigma ^{}}\delta _{x_0;x_0^{}\mathrm{}cq/eB},`$
$`r_{ss^{}}`$ $`=`$ $`{\displaystyle \frac{1}{g_0}}{\displaystyle \underset{\nu \nu ^{}}{}}{\displaystyle \frac{f_\nu f_\nu ^{}}{E_\nu E_\nu ^{}}}P^{|2s|}(p_z)Q^{|2s^{}|}(p_z)\times `$ (54)
$`\times `$ $`J_{l2s}\left({\displaystyle \frac{\mathrm{}cq}{eB}}\sqrt{2N+1}\right)J_{l2s^{}}\left({\displaystyle \frac{\mathrm{}cq}{eB}}\sqrt{2N+1}\right)`$
$`\times `$ $`\delta _{\sigma \sigma ^{}}\delta _{x_0;x_0^{}\mathrm{}cq/eB},`$
$`N_s`$ $`=`$ $`{\displaystyle \underset{\nu \nu ^{}}{}}{\displaystyle \frac{f_\nu f_\nu ^{}}{E_\nu E_\nu ^{}}}P^{|2s|}(p_z)`$ (55)
$`\times `$ $`J_{l2s}\left({\displaystyle \frac{\mathrm{}cq}{eB}}\sqrt{2N+1}\right)n_{\nu ^{}\nu }(𝐪);`$
$`N_s^{}`$ $`=`$ $`{\displaystyle \underset{\nu \nu ^{}}{}}{\displaystyle \frac{f_\nu f_\nu ^{}}{E_\nu E_\nu ^{}}}P^{|2s|}(p_z)`$ (56)
$`\times `$ $`J_{l2s}\left({\displaystyle \frac{\mathrm{}cq}{eB}}\sqrt{2N+1}\right)\sigma n_{\nu ^{}\nu }(𝐪).`$
The expressions for the renormalized matrix elements $`n_{\nu \nu ^{}}^{}(𝐪)`$ have the form
$$n_{\nu \nu ^{}}^{}(𝐪)=n_{Nlp_z}^{}(𝐪)\delta _{\sigma \sigma ^{}}\delta _{x_0;x_0^{}\mathrm{}cq/eB}\delta _{p_zp_z^{}},$$
where the quantity $`n_{Nlp_z}^{}(q)`$ satisfies the relation, following from the definition (7):
$`n_{Nlp_z}^{}(𝐪)=n_{Nlp_z}(𝐪)`$ (57)
$`+`$ $`{\displaystyle \frac{1}{g_0}}{\displaystyle \underset{s}{}}P^{|2s|}(p_z)J_{l2s}\left({\displaystyle \frac{\mathrm{}cq}{eB}}\sqrt{2N+1}\right)I_s`$
$`+`$ $`{\displaystyle \frac{\sigma }{g_0}}{\displaystyle \underset{s}{}}Q^{|2s|}(p_z)J_{l2s}\left({\displaystyle \frac{\mathrm{}cq}{eB}}\sqrt{2N+1}\right)I_s^{}.`$
When the Fermi surface is axisymmetric the matrix elements $`n_{Nlp_z}^{}(q)`$ do not depend on $`p_z:`$
$$n_{Nlp_z}(𝐪)=J_l\left(\frac{\mathrm{}cq}{eB}\sqrt{2N+1}\right).$$
(58)
The average quantities $`I_s`$ and $`I_s^{}`$ satisfy a system of linear equations which follows from the Eq. (57):
$$\{\begin{array}{c}I_s+\underset{s^{}}{}(g_{ss^{}}I_s^{}+r_{ss^{}}I_s^{}^{})=N_s,\hfill \\ I_s^{}+\underset{s^{}}{}(r_{s^{}s}I_s^{}+q_{ss^{}}I_s^{}^{})=N_s^{}.\hfill \end{array}$$
(59)
Using (24) we get the expressions for the coefficients of the system (59):
$`g_{ss^{}}`$ $`=`$ $`\stackrel{~}{A}_{2s}\delta _{ss^{}}+{\displaystyle \underset{m}{}}\mathrm{\Delta }_mP^{|2s|}(p_m)P^{|2s^{}|}(p_m)`$ (60)
$`\times `$ $`J_{2s}(qR_m)J_{2s^{}}(qR_m),`$
$`q_{ss^{}}`$ $`=`$ $`\stackrel{~}{B}_{2s}\delta _{ss^{}}+{\displaystyle \underset{m}{}}\mathrm{\Delta }_mQ^{|2s|}(p_m)Q^{|2s^{}|}(p_m)`$ (61)
$`\times `$ $`J_{2s}(qR_m)J_{2s^{}}(qR_m),`$
$`r_{ss^{}}`$ $`=`$ $`{\displaystyle \underset{m}{}}\mathrm{\Delta }_m^{}P^{|2s|}(p_m)Q^{|2s^{}|}(p_m)`$ (62)
$`\times `$ $`J_{2s}(qR_m)J_{2s^{}}(qR_m),`$
where $`R_m`$ in the radius of the cyclotron orbit corresponding to the $`m`$-th extremal cross section of the FS:
$`\stackrel{~}{A}_{2s}`$ $`=`$ $`{\displaystyle \frac{m_{}}{2\pi ^2\mathrm{}^3}}{\displaystyle \frac{1}{g_0}}{\displaystyle [P^{|2s|}(p_z)]^2𝑑p_z};`$
$`\stackrel{~}{B}_{2s}`$ $`=`$ $`{\displaystyle \frac{m_{}}{2\pi ^2\mathrm{}^3}}{\displaystyle \frac{1}{g_0}}{\displaystyle [Q^{|2s|}(p_z)]^2𝑑p_z}.`$ (63)
In obtaining Eqs. (60)-(62) we used the identity concerning the Bessel functions
$$\underset{l=\mathrm{}}{\overset{\mathrm{}}{}}J_{lm}(x)J_l(x)=\delta _{l0}.$$
(64)
When the FS has the unique extremal cross-section (at $`p_z=0)`$ Eq. (64) takes the form:
$$\{\begin{array}{ccc}\frac{N_s}{1+\stackrel{~}{A}_{2s}}\hfill & =I_s\hfill & +\mathrm{\Delta }\frac{P^{|2s|}(0)J_{2s}(qR_{ex})}{1+\stackrel{~}{A}_{2s}}X\hfill \\ & & +\mathrm{\Delta }^{}\frac{P^{|2s|}(0)J_{2s}(qR_{ex})}{1+\stackrel{~}{A}_{2s}}Y,\hfill \\ & & \\ \frac{N_s^{}}{1+\stackrel{~}{B}_{2s}}\hfill & =I_s^{}\hfill & +\mathrm{\Delta }^{}\frac{Q^{|2s|}(0)J_{2s}(qR_{ex})}{1+\stackrel{~}{B}_{2s}}X\hfill \\ & & +\mathrm{\Delta }\frac{Q^{|2s|}(0)J_{2s}(qR_{ex})}{1+\stackrel{~}{B}_{2s}}Y.\hfill \end{array}$$
Here:
$`X={\displaystyle \underset{s}{}}P^{|2s|}(0)J_{2s}(qR_{ex})I_s,`$ (65)
$`Y={\displaystyle \underset{s}{}}Q^{|2s|}(0)J_{2s}(qR_{ex})I_s^{}.`$ (66)
We can find the quantities $`X`$ and $`Y`$ solving the system of equations:
$$\{\begin{array}{c}X(1+\alpha _q\mathrm{\Delta })+Y\alpha _q\mathrm{\Delta }^{}=\underset{s}{}N_s\frac{P^{|2s|}(0)}{1+\stackrel{~}{A}_{2s}}J_{2s}(qR_{ex}),\hfill \\ X\beta _q\mathrm{\Delta }^{}+Y(1+\beta _q\mathrm{\Delta })=\underset{s}{}N_s^{}\frac{Q^{|2s|}(0)}{1+\stackrel{~}{B}_{2s}}J_{2s}(qR_{ex}).\hfill \end{array}$$
(67)
Here
$`\alpha _q`$ $`=`$ $`{\displaystyle \underset{s}{}}{\displaystyle \frac{A_{2s}^{}}{1+\stackrel{~}{A}_{2s}}}J_{2s}^2(qR_{ex}),`$
$`\beta _q`$ $`=`$ $`{\displaystyle \underset{s}{}}{\displaystyle \frac{B_{2s}^{}}{1+\stackrel{~}{B}_{2s}}}J_{2s}^2(qR_{ex})`$
$`A_{2s}^{}`$ $`=`$ $`[P^{|2s|}(0)]^2;B_{2s}^{}=[Q^{|2s|}(0)]^2.`$ (68)
As a result we arrive at the expressions:
$`X`$ $`=`$ $`{\displaystyle \frac{1+\beta \mathrm{\Delta }}{D}}{\displaystyle \underset{s}{}}N_s{\displaystyle \frac{P^{|2s|}(0)}{1+\stackrel{~}{A}_{2s}}}J_{2s}(qR_{ex})`$
$``$ $`{\displaystyle \frac{\alpha \mathrm{\Delta }^{}}{D}}{\displaystyle \underset{s}{}}N_s^{}{\displaystyle \frac{Q^{|2s|}(0)}{1+\stackrel{~}{B}_{2s}}}J_{2s}(qR_{ex}),`$
$`Y`$ $`=`$ $`{\displaystyle \frac{1+\alpha \mathrm{\Delta }}{D}}{\displaystyle \underset{s}{}}N_s^{}{\displaystyle \frac{Q^{|2s|}(0)}{1+\stackrel{~}{B}_{2s}}}J_{2s}(qR_{ex})`$ (69)
$``$ $`{\displaystyle \frac{\beta \mathrm{\Delta }^{}}{D}}{\displaystyle \underset{s}{}}N_s{\displaystyle \frac{P^{|2s|}(0)}{1+\stackrel{~}{A}_{2s}}}J_{2s}(qR_{ex}),`$
where the determinant of the system equals:
$$D=1+(\alpha _q+\beta _q)\mathrm{\Delta }+\alpha _q\beta _q(\mathrm{\Delta }^2\mathrm{\Delta }^2).$$
(70)
Sustituting the expressions (70) into the system (65), we can calculate the averages:
$`I_s`$ $`=`$ $`{\displaystyle \frac{N_s}{1+\stackrel{~}{A}_{2s}}}{\displaystyle \frac{\mathrm{\Delta }+\beta _q(\mathrm{\Delta }^2\mathrm{\Delta }^2)}{D}}{\displaystyle \frac{P^{|2s|}(0)}{1+\stackrel{~}{A}_{2s}}}J_{2s}(qR_{ex})`$ (71)
$`\times `$ $`{\displaystyle \underset{s^{}}{}}N_s^{}{\displaystyle \frac{P^{|2s^{}|}(0)}{1+\stackrel{~}{A}_{2s^{}}}}J_{2s^{}}(qR_{ex}){\displaystyle \frac{\mathrm{\Delta }^{}}{D}}{\displaystyle \frac{P^{|2s|}(0)}{1+\stackrel{~}{A}_{2s}}}`$
$`\times `$ $`J_{2s}(qR_{ex}){\displaystyle \underset{s^{}}{}}N_s^{}^{}{\displaystyle \frac{Q^{|2s^{}|}(0)}{1+\stackrel{~}{B}_{2s^{}}}}J_{2s^{}}(qR_{ex}),`$
$`I_s^{}`$ $`=`$ $`{\displaystyle \frac{N_s^{}}{1+\stackrel{~}{B}_{2s}}}{\displaystyle \frac{\mathrm{\Delta }+\alpha _q(\mathrm{\Delta }^2\mathrm{\Delta }^2)}{D}}{\displaystyle \frac{Q^{|2s|}(0)}{1+\stackrel{~}{A}_{2s}}}J_{2s}(qR_{ex})`$ (72)
$`\times `$ $`{\displaystyle \underset{s^{}}{}}N_s^{}^{}{\displaystyle \frac{Q^{|2s^{}|}(0)}{1+\stackrel{~}{A}_{2s^{}}}}J_{2s^{}}(qR_{ex}){\displaystyle \frac{\mathrm{\Delta }^{}}{D}}{\displaystyle \frac{Q^{|2s|}(0)}{1+\stackrel{~}{B}_{2s}}}`$
$`\times `$ $`J_{2s}(qR_{ex}){\displaystyle \underset{s^{}}{}}N_s^{}{\displaystyle \frac{P^{|2s^{}|}(0)}{1+\stackrel{~}{A}_{2s^{}}}}J_{2s^{}}(qR_{ex}).`$
The asymptotic expressions for the quantities $`N_s`$ and $`N_s^{}`$ also can be found by the formula (24). For the FS with the unique extremal cross-section we have:
$`N_s`$ $`=`$ $`{\displaystyle \frac{m_{}}{2\pi ^2\mathrm{}^3}}{\displaystyle 𝑑p_zP^0(p_z)\delta _{s0}}g_0\mathrm{\Delta }J_0(qR_{ex})`$
$`\times `$ $`P^{|2s|}(0)J_{2s}(qR_{ex}),`$
$`N_s^{}`$ $`=`$ $`g_0\mathrm{\Delta }^{}J_0(qR_{ex})Q^{|2s|}(0)J_{2s}(qR_{ex}).`$ (73)
Using the Eqs. (57) and (72)-(74), we obtain
$``$ $`{\displaystyle \underset{\nu \nu ^{}}{}}{\displaystyle \frac{f_\nu f_\nu ^{}}{E_\nu E_\nu ^{}}}n_{\nu \nu ^{}}^{}(𝐪)n_{\nu ^{}\nu }(𝐪)=g_0[1\alpha _0`$ (74)
$`+`$ $`{\displaystyle \frac{[(1\overline{\alpha }_0)^2J_0^2(qR_{ex})][\mathrm{\Delta }+\beta _q(\mathrm{\Delta }^2\mathrm{\Delta }^2)]}{1+(\alpha _q+\beta _q)\mathrm{\Delta }+\alpha _q\beta _q(\mathrm{\Delta }^2\mathrm{\Delta }^2)}}].`$
Here we introduced additional notations:
$`\alpha _0={\displaystyle \frac{A_0}{1+\stackrel{~}{A}_0}},\overline{\alpha }_0={\displaystyle \frac{\overline{A}_0}{1+\stackrel{~}{A}_0}},`$
$`A_0=\left({\displaystyle \frac{m_{}}{2\pi ^2\mathrm{}^3g_0}}{\displaystyle P^0(p_z)𝑑p_z}\right)^2,`$
$`\overline{A}_0={\displaystyle \frac{P^0(0)m_{}}{2\pi ^2\mathrm{}^3g_0}}{\displaystyle P^0(p_z)𝑑p_z}.`$ (75)
Passing to the limit $`q0`$ in the Eq. (75) we obtain the final result for $`N_\zeta ^{}.`$
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# Contents
## 1 Introduction
A natural problem is to understand what a typical element of the finite general linear group $`GL(n,q)`$ “looks like”. Many of the interesting properties of a random matrix depend only on its conjugacy class. The following list of questions one could ask are of this type:
1. How many Jordan blocks are there in the rational canonical form of a random matrix?
2. What is the distribution of the order of a random matrix?
3. What is the probability that the characteristic polynomial of a random matrix has no repeated factors?
4. What is the probability that the characteristic polynomial of a random matrix is equal to its minimal polynomial?
5. What is the probability that a random matrix is semisimple (i.e. diagonalizable over the algebraic closure $`\overline{F_q}`$ of the field of $`q`$ elements)?
As Section 2 will indicate, answers to these questions have applications to the study of random number generators, to the analysis of algorithms in computational group theory, and to other parts of group theory. Section 2 describes a unified approach to answering such probability questions using cycle index generating functions. As an example of its power, it is proved independently in \[F1\] and \[W2\] that the $`n\mathrm{}`$ limit of answer to Question 3 is $`(1\frac{1}{q^5})/(1+\frac{1}{q^3})`$. There is (at present) no other method for deriving this result and generating functions give effective bounds on the convergence rate to the limit. Extensions of the cycle index method to the set of all matrices and to other finite classical groups are sketched.
Section 3 gives a purely probabilistic picture of what the conjugacy class of a random element of $`GL(n,q)`$ looks like. The main object of study is a probability measure $`M_{GL,u,q}`$ on the set of all partitions of all natural numbers. This measure is connected with the Hall-Littlewood symmetric functions. Exploiting this connection leads to several methods for growing random partitions distributed as $`M_{GL,u,q}`$ and gives insightful probabilistic proofs of group theoretic results. We hope to convince the reader that the interplay between probability and symmetric functions is beautiful and useful. A method is given for sampling from $`M_{GL,u,q}`$ conditioned to live on partitions of a fixed size (which amounts to studying Jordan form of unipotent elements) and for sampling from a $`q`$-analog of Plancherel measure (which is related to the longest increasing subsequence problem of random permutations).
Section 3 goes on to describe a probabilistic approach to $`M_{GL,u,q}`$ using Markov chains. This connection is quite surprising, and we indicate how it leads to a simple and motivated proof of the Rogers-Ramanujan identities. The measure $`M_{GL,u,q}`$ has analogs for the finite unitary, symplectic, and orthogonal groups. As this is somewhat technical these results are omitted and pointers to the literature are given. However we remark now that while the analogs of the symmetric function theory viewpoint are unclear for the finite symplectic and orthogonal groups, the connections with Markov chains carry over. Thus there is a coherent probabilistic picture of the conjugacy classes of the finite classical groups.
Section 4 surveys probabilistic aspects of conjugacy classes in $`T(n,q)`$, the group of $`n\times n`$ upper triangular matrices over the field $`F_q`$ with $`1`$’s along the main diagonal. Actually a simpler object is studied, namely the Jordan form of randomly chosen elements of $`T(n,q)`$. From work of Borodin and Kirillov, one can sample from the corresponding measures on partitions. We link their results with symmetric function theory and potential theory on Bratteli diagrams.
The field surveyed in this article is young and evolving. The applications to computational group theory call for extensions of probability estimates discussed in Section 2 to maximal subgroups of finite classical groups. It would be marvellous if the program surveyed here carries over; this happens for the finite affine groups \[F9\]. The first step is understanding conjugacy classes and partial results can be found in the thesis \[Mu\].
We close with a final motivation for the study of conjugacy classes of random matrices over finite fields. The past few years have seen an explosion of interest in eigenvalues of random matrices from compact Lie groups. For the unitary group $`U(n,C)`$ over the complex numbers, two matrices are in the same conjugacy class if and only if they have the same set of eigenvalues. Hence, at least in this case, which is related to the zeroes of the Riemann zeta function \[KeaSn\], the study of eigenvalues is the same as the study of conjugacy classes.
As complements to this article, the reader may enjoy the surveys \[Sh2\],\[Py1\],\[Py2\]. These articles discuss probabilistic and enumerative questions in group theory and have essentially no overlap with the program surveyed here.
## 2 Cycle Index Techniques
Before describing cycle index techniques for the finite classical groups, we mention that the cycle index techniques here are modelled on similar techniques for the study of conjugacy class functions on the symmetric groups. For a permutation $`\pi `$, let $`n_i(\pi )`$ be the number of length $`i`$ cycles of $`\pi `$. The cycle index of a subgroup $`G`$ of $`S_n`$ is defined as
$$\frac{1}{|G|}\underset{\pi G}{}\underset{i1}{}x_i^{n_i(\pi )}$$
and is called a cycle index because it stores information about the cycle structure of elements of $`G`$. Applications of the cycle index to graph theory and chemical compounds are exposited in \[PoRe\]. It is standard to refer to the generating function
$$1+\underset{n1}{}\frac{u^n}{n!}\underset{\pi S_n}{}\underset{i1}{}x_i^{n_i(\pi )}$$
as the cycle index or cycle index generating function of the symmetric groups. From the fact that there are $`\frac{n!}{_in_i!i^{n_i}}`$ elements in $`S_n`$ with $`n_i`$ cycles of length $`i`$, one deduces Polya’s result that this generating function is equal to $`_{m1}e^{\frac{x_mu^m}{m}}`$. This allows one to study conjugacy class functions of random permutations (e.g. number of fixed points, number of cycles, the order of a permutation, length of the longest cycle) by generating functions. We refer the reader to \[Ko\] for results in this direction using analysis and to \[ShLl\] for results about cycle structure proved by a probabilistic interpretation of the cycle index generating function.
Subsection 2.1 reviews the conjugacy classes of $`GL(n,q)`$ and then discusses cycles indices for $`GL(n,q)`$ and $`Mat(n,q)`$, the set of all $`n\times n`$ matrices with entries in the field of $`q`$ elements. Subsection 2.2 describes applications of cycle index techniques. Subsection 2.3 discusses generalizations of cycle indices to the finite classical groups.
It is necessary to recall some standard notation. Let $`\lambda `$ be a partition of some non-negative integer $`|\lambda |`$ into integer parts $`\lambda _1\lambda _2\mathrm{}0`$. We will also write $`\lambda n`$ if $`\lambda `$ is a partition of $`n`$. Let $`m_i(\lambda )`$ be the number of parts of $`\lambda `$ of size $`i`$, and let $`\lambda ^{}`$ be the partition dual to $`\lambda `$ in the sense that $`\lambda _i^{}=m_i(\lambda )+m_{i+1}(\lambda )+\mathrm{}`$. Let $`n(\lambda )`$ be the quantity $`_{i1}(i1)\lambda _i`$ and let $`(\frac{u}{q})_i`$ denote $`(1\frac{u}{q})\mathrm{}(1\frac{u}{q^i})`$.
### 2.1 The General Linear Groups
To begin we follow Kung \[Kun\] in defining a cycle index for $`GL(n,q)`$. First it is necessary to understand the conjugacy classes of $`GL(n,q)`$. As is explained in Chapter 6 of the textbook \[Her\], an element $`\alpha GL(n,q)`$ has its conjugacy class determined by its rational canonical form. This form corresponds to the following combinatorial data. To each monic non-constant irreducible polynomial $`\varphi `$ over $`F_q`$, associate a partition (perhaps the trivial partition) $`\lambda _\varphi `$ of some non-negative integer $`|\lambda _\varphi |`$. Let $`deg(\varphi )`$ denote the degree of $`\varphi `$. The only restrictions necessary for this data to represent a conjugacy class are that $`|\lambda _z|=0`$ and $`_\varphi |\lambda _\varphi |deg(\varphi )=n.`$
An explicit representative of this conjugacy class may be given as follows. Define the companion matrix $`C(\varphi )`$ of a polynomial $`\varphi (z)=z^{deg(\varphi )}+\alpha _{deg(\varphi )1}z^{deg(\varphi )1}+\mathrm{}+\alpha _1z+\alpha _0`$ to be:
$$\left(\begin{array}{ccccc}0& 1& 0& \mathrm{}& 0\\ 0& 0& 1& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& 0& \mathrm{}& 1\\ \alpha _0& \alpha _1& \mathrm{}& \mathrm{}& \alpha _{deg(\varphi )1}\end{array}\right)$$
Let $`\varphi _1,\mathrm{},\varphi _k`$ be the polynomials such that $`|\lambda _{\varphi _i}|>0`$. Denote the parts of $`\lambda _{\varphi _i}`$ by $`\lambda _{\varphi _i,1}\lambda _{\varphi _i,2}\mathrm{}`$. Then a matrix corresponding to the above conjugacy class data is
$$\left(\begin{array}{cccc}R_1& 0& 0& 0\\ 0& R_2& 0& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& 0& R_k\end{array}\right)$$
where $`R_i`$ is the matrix
$$\left(\begin{array}{ccc}C(\varphi _i^{\lambda _{\varphi _i,1}})& 0& 0\\ 0& C(\varphi _i^{\lambda _{\varphi _i,2}})& 0\\ 0& 0& \mathrm{}\end{array}\right)$$
For example, the identity matrix has $`\lambda _{z1}`$ equal to $`(1^n)`$ and all other $`\lambda _\varphi `$ equal to the emptyset. An elementary transvection with $`a0`$ in the $`(1,2)`$ position, ones on the diagonal and zeros elsewhere has $`\lambda _{z1}`$ equal to $`(2,1^{n2})`$ and all other $`\lambda _\varphi `$ equal to the emptyset. For a given matrix only finitely many $`\lambda _\varphi `$ are non-empty. Many algebraic properties of a matrix can be stated in terms of the data parameterizing its conjugacy class. For instance the characteristic polynomial of $`\alpha GL(n,q)`$ is equal to $`_\varphi \varphi ^{|\lambda _\varphi (\alpha )|}`$ and the minimal polynomial of $`\alpha `$ is equal to $`_\varphi \varphi ^{|\lambda _{\varphi ,1}(\alpha )|}`$. Furthermore $`\alpha `$ is semisimple (diagonalizable over the algebraic closure $`\overline{F_q}`$) precisely when all $`\lambda _\varphi (\alpha )`$ have largest part at most 1.
To define the cycle index for $`Z_{GL(n,q)}`$, let $`x_{\varphi ,\lambda }`$ be variables corresponding to pairs of polynomials and partitions. Define
$$Z_{GL(n,q)}=\frac{1}{|GL(n,q)|}\underset{\alpha GL(n,q)}{}\underset{\varphi :|\lambda _\varphi (\alpha )|>0}{}x_{\varphi ,\lambda _\varphi (\alpha )}.$$
Note that the coefficient of a monomial is the probability of belonging to the corresponding conjugacy class, and is therefore equal to one over the order of the centralizer of a representative. It is well known (e.g. easily deduced from page 181 of \[Mac\]) that one over the order of the centralizer of conjugacy class of $`GL(n,q)`$ corresponding to the data $`\{\lambda _\varphi \}`$ is
$$\frac{1}{_\varphi q^{deg(\varphi )_i(\lambda _{\varphi ,i}^{})^2}_{i1}(\frac{1}{q^{deg(\varphi )}})_{m_i(\lambda _\varphi )}}.$$
The formulas given for conjugacy class size in \[Kun\] and \[St1\] are written in different form; for the reader’s benefit they have been expressed here in the form most useful to us. It follows that
$$1+\underset{n=1}{\overset{\mathrm{}}{}}Z_{GL(n,q)}u^n=\underset{\varphi z}{}\left[1+\underset{n1}{}\underset{\lambda n}{}x_{\varphi ,\lambda }\frac{u^{ndeg(\varphi )}}{q^{deg(\varphi )_i(\lambda _i^{})^2}_{i1}\left(\frac{1}{q^{deg(\varphi )}}\right)_{m_i(\lambda _\varphi )}}\right].$$
This is called the cycle index generating function.
Let $`Mat(n,q)`$ be the set of all $`n\times n`$ matrices over the field $`F_q`$. Define
$$Z_{Mat(n,q)}=\frac{1}{|GL(n,q)|}\underset{\alpha Mat(n,q)}{}\underset{\varphi :|\lambda _\varphi (\alpha )|>0}{}x_{\varphi ,\lambda _\varphi (\alpha )}.$$
Analogous arguments \[St1\] show that
$$1+\underset{n=1}{\overset{\mathrm{}}{}}Z_{Mat(n,q)}u^n=\underset{\varphi }{}\left[1+\underset{n1}{}\underset{\lambda n}{}x_{\varphi ,\lambda }\frac{u^{ndeg(\varphi )}}{q^{deg(\varphi )_i(\lambda _i^{})^2}_{i1}\left(\frac{1}{q^{deg(\varphi )}}\right)_{m_i(\lambda _\varphi )}}\right].$$
This will be used in Subsection 2.2. Note that the denominator in $`Z_{Mat(n,q)}`$ is $`|GL(n,q)|`$, not $`|Mat(n,q)|`$, since the formula follows from a formula for the size of the orbits of $`GL(n,q)`$ acting on $`Mat(n,q)`$ by conjugation. This makes no essential difference for applications.
### 2.2 Applications
This subsection describes applications of cycle indices. The first example is treated in detail and results for the other examples are sketched.
Example 1: Cyclic and Separable Matrices
Recall that a matrix $`\alpha Mat(n,q)`$ operating on a vector space $`V`$ is called cyclic if there is a vector $`v_0V`$ such that $`v_0,v_0\alpha ,v_0\alpha ^2,\mathrm{}`$ span $`V`$. As is explained in \[NP2\], this is equivalent to the condition that the characteristic and minimal polynomials of $`\alpha `$ are equal.
The need to estimate the proportion of cyclic matrices arose from \[NP1\] in connection with analyzing the running time of an algorithm for deciding whether or not the group generated by a given set of matrices in $`GL(n,q)`$ contains the special linear group $`SL(n,q)`$. Cyclic matrices also arise in recent efforts to improve upon the MeatAxe algorithm for computing modular characters \[NP4\] and in Example 8 below. John Thompson has asked if every matrix is the product of a cyclic matrix and a permutation matrix, suggesting that the answer could have applications to finite projective planes.
Letting $`c_M(n,q)`$ be the proportion of cyclic elements of $`Mat(n,q)`$, the paper \[NP2\] proves that
$$\frac{1}{q^2(q+1)}<1c_M(n,q)<\frac{1}{(q^21)(q1)}.$$
The cycle index approach is also informative, yielding a formula for the $`n\mathrm{}`$ limit of $`C_M(n,q)`$, denoted by $`c_M(\mathrm{},q)`$, together with convergence rates. For the argument two lemmas are useful, as is some notation. Let $`N_d(q)`$ be the number of degree $`d`$ irreducible polynomials over the field $`F_q`$. In all that follows $`\varphi `$ will denote a monic irreducible polynomial over $`F_q`$. Given a power series $`f(u)`$, let $`[u^n]f(u)`$ denote the coefficient of $`u^n`$ in $`f(u)`$.
###### Lemma 1
$$\underset{\varphi }{}(1\frac{u^{deg(\varphi )}}{q^{deg(\varphi )}})=1u$$
Proof: Expanding $`\frac{1}{1\frac{u^{deg(\varphi )}}{q^{deg(\varphi )}}}`$ as a geometric series and using unique factorization in $`F_q[x]`$, one sees that the coefficient of $`u^d`$ in the reciprocal of the left hand side is $`\frac{1}{q^d}`$ times the number of monic polynomials of degree $`d`$, hence 1. Comparing with the reciprocal of the right hand side completes the proof. $`\mathrm{}`$
###### Lemma 2
If the Taylor series of $`f`$ around 0 converges at $`u=1`$, then
$$lim_n\mathrm{}[u^n]\frac{f(u)}{1u}=f(1).$$
Proof: Write the Taylor expansion $`f(u)=_{n=0}^{\mathrm{}}a_nu^n`$. Then observe that $`[u^n]\frac{f(u)}{1u}=_{i=0}^na_i`$. $`\mathrm{}`$
Theorem 1 calculates $`c_M(\mathrm{},q)`$.
###### Theorem 1
(\[F1\],\[W2\])
$$c_M(\mathrm{},q)=(1\frac{1}{q^5})\underset{r=3}{\overset{\mathrm{}}{}}(1\frac{1}{q^r})$$
Proof: Recall that $`\alpha `$ is cyclic precisely when its characteristic polynomial and minimal polynomials are equal. From Subsection 2.1, these polynomials are equal when all $`\lambda _\varphi `$ have at most one part. In the cycle index for $`Mat(n,q)`$ set $`x_{\varphi ,\lambda }=1`$ if $`\lambda `$ has at most 1 part and $`x_{\varphi ,\lambda }=0`$ otherwise. It follows that
$$c_M(n,q)=\frac{|GL(n,q)|}{q^{n^2}}[u^n]\underset{\varphi }{}(1+\underset{j=1}{\overset{\mathrm{}}{}}\frac{u^{jdeg(\varphi )}}{q^{(j1)deg(\varphi )}(q^{deg(\varphi )}1)}).$$
By Lemma 1 this equation can be rewritten as
$`c_M(n,q)`$ $`=`$ $`{\displaystyle \frac{|GL(n,q)|}{q^{n^2}}}[u^n]{\displaystyle \frac{_\varphi (1\frac{u^{deg(\varphi )}}{q^{deg(\varphi )}})(1+_{j=1}^{\mathrm{}}\frac{u^{jdeg(\varphi )}}{q^{(j1)deg(\varphi )}(q^{deg(\varphi )}1)})}{1u}}`$
$`=`$ $`{\displaystyle \frac{|GL(n,q)|}{q^{n^2}}}[u^n]{\displaystyle \frac{_\varphi (1+\frac{u^{deg(\varphi )}}{q^{deg(\varphi )}(q^{deg(\varphi )}1)})}{1u}}`$
$`=`$ $`{\displaystyle \frac{|GL(n,q)|}{q^{n^2}}}[u^n]{\displaystyle \frac{_{d1}(1+\frac{u^d}{q^d(q^d1)})^{N_d(q)}}{1u}}.`$
Recall that a product $`_{n=1}^{\mathrm{}}(1+a_n)`$ converges absolutely if the series $`_{n1}|a_n|`$ converges. Thus using the crude bound $`N_d(q)q^d`$
$$\underset{d1}{}(1+\frac{u^d}{q^d(q^d1)})^{N_d(q)}$$
is analytic in a disc of radius greater than $`1`$. Lemma 2 implies that
$`c_M(\mathrm{},q)`$ $`=`$ $`lim_n\mathrm{}{\displaystyle \frac{|GL(n,q)|}{q^{n^2}}}[u^n]{\displaystyle \frac{_{d1}(1+\frac{u^d}{q^d(q^d1)})^{N_d(q)}}{1u}}`$
$`=`$ $`{\displaystyle \underset{r=1}{\overset{\mathrm{}}{}}}(1{\displaystyle \frac{1}{q^r}}){\displaystyle \underset{d1}{}}(1+{\displaystyle \frac{1}{q^d(q^d1)}})^{N_d(q)}.`$
Applying Lemma 1 (with $`u=\frac{1}{q}`$, $`u=\frac{1}{q^2}`$ and then $`u=\frac{1}{q^5}`$) gives that
$`c_M(\mathrm{},q)`$ $`=`$ $`{\displaystyle \underset{r=3}{\overset{\mathrm{}}{}}}(1{\displaystyle \frac{1}{q^r}}){\displaystyle \underset{d1}{}}((1+{\displaystyle \frac{1}{q^d(q^d1)}})(1{\displaystyle \frac{1}{q^{2d}}})(1{\displaystyle \frac{1}{q^{3d}}}))^{N_d(q)}`$
$`=`$ $`{\displaystyle \underset{r=3}{\overset{\mathrm{}}{}}}(1{\displaystyle \frac{1}{q^r}}){\displaystyle \underset{d1}{}}(1{\displaystyle \frac{1}{q^{6d}}})^{N_d(q)}`$
$`=`$ $`(1{\displaystyle \frac{1}{q^5}}){\displaystyle \underset{r=3}{\overset{\mathrm{}}{}}}(1{\displaystyle \frac{1}{q^r}}).`$
$`\mathrm{}`$
The next challenge is to bound the convergence rate of $`c_M(n,q)`$ to $`c_{\mathrm{}}(n,q)`$. Wall \[W2\] found a strikingly simple way of doing this by relating the cycle index of cyclic matrices to the cycle index of the set of matrices whose characteristic polynomial is squarefree (these matrices are termed separable in \[NP2\]). To state the result, let $`s_M(n,q)`$ be the probability that an $`n\times n`$ matrix is separable. Next let $`C_M(u,q)`$ and $`S_M(u,q)`$ be the generating functions defined as
$$C_M(u,q)=1+\underset{n1}{}\frac{u^nq^{n^2}}{|GL(n,q)|}c_M(n,q)$$
$$S_M(u,q)=1+\underset{n1}{}\frac{u^nq^{n^2}}{|GL(n,q)|}s_M(n,q).$$
###### Lemma 3
(\[W2\])
$$(1u)C_M(u,q)=S_M(u/q,q).$$
Proof: The proof of Theorem 1 shows that
$$(1u)C_M(u,q)=\underset{d1}{}(1+\frac{u^d}{q^d(q^d1)})^{N_d(q)}.$$
A matrix is separable if and only if all $`\lambda _\varphi `$ have size $`0`$ or $`1`$. Hence
$$S_M(u,q)=\underset{d1}{}(1+\frac{u^d}{q^d1})^{N_d(q)}.$$
The result follows. $`\mathrm{}`$
###### Corollary 1
(\[W2\])
$$0<|c_M(n,q)c_M(\mathrm{},q)|<\frac{1}{q^{n+1}(11/q)}$$
Proof: Taking coefficients of $`u^{n+1}`$ on both sides of Lemma 3 gives the relation
$$c_M(n+1,q)c_M(n,q)=\frac{s_M(n+1,q)c_M(n,q)}{q^{n+1}}.$$
Since $`0|s_M(n+1,q)c_M(n,q)|1`$ for all $`n`$, it follows that
$$|c_M(n,q)c_M(\mathrm{},q)|\underset{i=n}{\overset{\mathrm{}}{}}|c_M(i+1,q)c_M(i,q)|\underset{i=n}{\overset{\mathrm{}}{}}\frac{1}{q^{i+1}},$$
as desired. $`\mathrm{}`$
Remarks:
1. As mentioned in the introduction, an argument similar to that of Theorem 1 shows that the $`n\mathrm{}`$ probability that an element of $`GL(n,q)`$ is cyclic is $`(1\frac{1}{q^5})/(1+\frac{1}{q^3})`$. For large $`q`$ this goes like $`11/q^3`$. The reason for this is a result of Steinberg \[Stei\] stating that the set of non-regular elements in an algebraic group has co-dimension 3. In type $`A`$, regular (i.e. centralizer of minimum dimension) and cyclic elements coincide, but not always. For more discussion on this point, see \[NP2\], \[FNP\].
2. The generating functions $`S_M(u,q)`$ and $`C_M(u,q)`$ have intriguing analytical properties. It is proved in \[W2\] that
$$S_M(u,q)=\frac{_{d=1}^{\mathrm{}}(1\frac{u^d(u^d1)}{q^d(q^d1)})^{N_d(q)}}{1u}.$$
Thus $`S_M(u,q)`$ has a pole at 1 and $`S_M(u,q)\frac{1}{1u}`$ can be analytically extended to the circle of radius $`q`$. Analogous properties hold for $`C_M(u,q)`$ by means of Lemma 3.
3. The limits $`s_M(\mathrm{},q)`$ and $`s_{GL}(\mathrm{},q)`$ are in \[F1\],\[W2\]. Bounding the rate of convergence of $`s_M(n,q)`$ to $`s_M(\mathrm{},q)`$ leads to interesting number theory. Let $`p(d)`$ be the number of partitions of $`d`$ and let $`p_2(d)=_{i=0}^dp(d)`$. It is proved in \[W2\] that
$$|\frac{s_M(n,q)q^{n^2}}{|GL(n,q)|}1|\underset{d=n+1}{\overset{\mathrm{}}{}}(p_2(d)+qp(d2))q^d\frac{1}{3}(\frac{4q+27}{2q3})(\frac{2}{3}q)^n.$$
4. Lehrer \[Leh\] expresses $`s_M(n,q)`$ and $`s_{GL}(n,q)`$ as inner products of characters in the symmetric group and proves a stability result about their expansions in power of $`q^1`$. See also \[W2\] and \[LehSe\].
5. The results of \[F2\] and \[W2\] surveyed above are extended to the finite classical groups in \[FNP\]. The paper \[FlJ\] gives (intractable) formulas for the chance of being separable in groups such as $`SL(n,q)`$ (i.e. semisimple and simply connected).
Example 2: Eigenvalue free matrices
The paper \[NP3\] studies eigenvalues free matrices (i.e. matrices without fixed lines) over finite fields as a step in obtaining estimates of cyclic probabilities in orthogonal groups \[NP5\]. It is interesting that the study of eigenvalue free matrices was one of the motivations for the original papers \[Kun\],\[St1\], the latter of which proves that the $`n,q\mathrm{}`$ limit of the chance that an element of $`GL(n,q)`$ has no eigenvalues is $`\frac{1}{e}`$.
The $`n\mathrm{}`$ probability that a random element of $`S_n`$ has no fixed points is also $`\frac{1}{e}`$. This is not coincidence; in general the $`q\mathrm{}`$ limit of the chance that the characteristic polynomial of a random element of $`Mat(n,q)`$ factors into $`n_i`$ degree $`i`$ irreducible factors is the same as the probability that an element of $`S_n`$ factors into $`n_i`$ cycles of degree $`i`$. This is proved at the end of \[St1\] and is extended to finite Lie groups in \[F1\] using the combinatorics of maximal tori. There is another interesting line of argument which should be mentioned. It is easy to see from the cycle index that the factorization type of the characteristic polynomial of a random element of $`Mat(n,q)`$ and the factorization type of a random degree $`n`$ polynomial over $`F_q`$ have the same distribution as $`q\mathrm{}`$. Now the factorization type of a random degree $`n`$ polynomial over $`F_q`$ has same distribution as the cycle type of a random permutation distributed as a $`q`$-shuffle on $`n`$ cards \[DiaMcPi\], and as $`q\mathrm{}`$ a $`q`$-shuffle converges to a random permutation. The connection of Lie theory with card shuffling may seem adhoc, but is really the tip of a deep iceberg \[F10\].
Example 3: Characteristic polynomials
The previous example is a special case of the problem of studying the degrees of the factors of the characteristic polynomial of a random matrix. Many results in this direction (all proved used cycle indices) can be found in Stong’s paper \[St1\]. Hansen and Schmutz \[HSchm\] use cycle index manipulations to prove that if one ignores factors of small degree, then the factorization type of the characteristic polynomial of a random element of $`GL(n,q)`$ is close to the factorization type of a random degree $`n`$ polynomial over $`F_q`$. More precisely, let $`A_{n,l}`$ be the set of sequences $`(\alpha _{l+1},\mathrm{},\alpha _n)`$ where $`\alpha _i`$ is the number of degree $`i`$ factors of a random polynomial chosen from some measure. Let $`Q_n^{(1)}`$ be the measure on polynomials arising from characteristic polynomials of random elements of $`GL(n,q)`$ and let $`Q_n^{(2)}`$ be the measure arising from choosing a degree $`n`$ polynomial over $`F_q`$ uniformly at random. They prove
###### Theorem 2
(\[HSchm\]) There exists constants $`c_1,c_2`$ such that for all $`l`$ with $`c_1log(n)ln`$ and $`BN^{nl}`$,
$$|Q_n^{(1)}(A_n(B))Q_n^{(2)}(A_n(B))|<c_2/l.$$
The final section of their paper uses this principle to prove results about characteristic polynomials of random matrices using known results about random polynomials. A useful reference on the distribution of degrees of random polynomials over finite fields is \[ArBarT\].
Example 4: Generating transvections
Recall that the motivation behind Example 1 was a group recognition problem, i.e. trying to determine whether or not the group generated by a given set $`X`$ of matrices in $`GL(n,q)`$ contains the special general linear group $`SL(n,q)`$. However the problem still remains of making the recognition algorithm constructive. For instance if the group generated by $`X`$ is $`GL(n,q)`$ it would be desirable to write any element of $`GL(n,q)`$ as a word in $`X`$.
The paper \[CeLg\] proposes such a constructive recognition algorithm. An essential step involves constructing a transvection, that is a non-identity element of $`SL(n,q)`$ which has an $`n1`$ dimensional fixed space. This in turn is done in two steps. First, find an element $`\alpha `$ of $`GL(n,q)`$ conjugate to diag$`(C((z\tau )^2),R)`$ where $`C`$ is the companion matrix as in Subsection 2.1 and $`R`$ is semisimple without $`\tau `$ as an eigenvalue. Second, one checks that raising $`\alpha `$ to the least common multiple of the orders of $`\tau `$ and $`R`$ gives a transvection.
Thus it necessary to bound the number of feasible $`\alpha `$ in the first step. Such $`\alpha `$ have conjugacy class data $`\lambda _{z\tau }=(2)`$ and all other $`\lambda _\varphi `$ have largest part at most 1. The cycle index approach gives bounds improving on those in \[CeLg\]; see \[FNP\] for the details.
Example 5: Semisimple matrices
A fundamental problem in computational group theory is to construct an element of order $`p`$. Given a group element $`g`$ with order a multiple of $`p`$, this can be done by raising $`g`$ to an appropriate power. It is proved in \[IsKanSp\] that if $`G`$ is a permutation group of degree $`n`$ with order divisible by $`p`$, then the probability that a random element of $`G`$ has order divisible by $`p`$ is at least $`\frac{1}{n}`$.
Their proof reduces the assertion to simple groups and then uses the classification of simple groups. Let us consider the group $`GL(n,q)`$, which is close enough to simple to be useful for the applications at hand. When $`p`$ is the characteristic of the field of definition of $`GL(n,q)`$, an element has order prime to $`p`$ precisely when it is semisimple. Thus the problem is to study the probability that an element of $`GL(n,q)`$ is semisimple. The paper \[GuLub\] shows that if $`G`$ is a simple Chevalley group, then the probability of not being semisimple is at most $`3/(q1)+2/(q1)^2`$ and thus at most $`c/q`$ for some constant $`c`$ as conjectured by Kantor.
As mentioned earlier, a matrix $`\alpha `$ is semisimple if and only if all $`\lambda _\varphi (\alpha )`$ have largest part size at most 1. Stong \[St1\] used cycle indices to obtain crude asymptotic bounds for the probability that an element of $`GL(n,q)`$ is semisimple. The thesis \[F1\] used the Rogers-Ramanujan identities to prove that the $`n\mathrm{}`$ probability that an element of $`GL(n,q)`$ is semisimple is
$$\underset{\genfrac{}{}{0pt}{}{r=1}{r=0,\pm 2(mod5)}}{\overset{\mathrm{}}{}}\frac{(1\frac{1}{q^{r1}})}{(1\frac{1}{q^r})}.$$
The paper \[FNP\] gives effective bounds for finite $`n`$.
Example 6: Order of a matrix
A natural problem is to study the order of a random matrix. This has been done in \[St2\] and \[Schm\]; see also the remarks in Subsection 3.3 and the very preliminary calculations for other classical groups in \[F1\]. Shalev \[Sh1\] uses facts about the distribution of the order of a random matrix together with Aschbacher’s study of maximal subgroups of classical groups \[As\] as key tools in studying the probability that a random element of $`GL(n,q)`$ belongs to an irreducible subgroup of $`GL(n,q)`$ that does not contain $`SL(n,q)`$. As explained in \[Sh1\] this has a number of appications; for instance it leads to a proof that if $`x`$ is any non-trivial element of $`PSL(n,q)`$ then the probability that $`x`$ and a randomly chosen element $`y`$ generate $`PSL(n,q)`$ tends to $`1`$ as $`q\mathrm{}`$. Shalev \[Sh1\] asks for extensions of these results to other finite classical groups.
It is also useful to count elements of given orders (e.g. $`2`$ or $`3`$) in classical groups and their maximal subgroups. The recent paper \[CTY\] uses cycle indices to perform such enumerations. One motivation for such enumerations is the study of finite simple quotients of $`PSL(2,Z)`$; a group $`G`$ is a quotient of $`PSL(2,Z)`$ if and only if $`G=<x,y>`$ with $`x^2=y^3=1`$. For further discussion, see \[Sh2\].
Example 7: Random number generators
We follow \[Mar\],\[MarTs\] in indicating the relevance of random matrix theory to the study of random number generators. Suppose one wants to test a mechanism for generating a random integer between $`0`$ and $`2^{33}1`$. In base $`2`$ these are length $`33`$ binary vectors. Generating say $`n`$ of these and listing them gives an $`n\times 33`$ matrix. If the random generator were perfect, the arising matrix would be random. One could choose a statistic such as the rank of a matrix and compare the generation method with theory. They report that shift-register generators will fail such tests but that congruential generators usually pass. It would be interesting to see how various random number generators perform when tested using other conjugacy class functions of random matrices.
Diaconis and Graham \[DiaGr\] analyze random walks of the form $`X_n=AX_{n1}+ϵ_n`$ where $`X_i`$ is a length $`d`$ $`01`$ vector, $`A`$ is an element of $`GL(n,2)`$, and $`ϵ_n`$ is a random vector of disturbance terms. For more general $`A`$ (in $`GL(n,q)`$) this includes the problem of running a psuedo-random number generator with recurrence $`Y_n=a_1Y_{n1}+\mathrm{}+a_dY_{nd}+ϵ_n`$ with $`Y_iF_q`$. They show that the rational canonical form of $`A`$ is related in a subtle way to the convergence rate of the walk. It would be interesting to understand what happens when $`A`$ is a random matrix.
Example 8: Product replacement algorithm
In recent years finite group theory has become much more computational. Given a generating set $`S`$ of a finite group $`G`$, it is natural to seek random elements of $`G`$. One approach, implemented in the computer systems GAP and MAGMA, is the product replacement algorithm \[CeLgMuNiOb\]. Fixing $`G`$ and some $`k`$, one performs a random walk on $`k`$-tuples $`(g_1,\mathrm{},g_k)`$ of elements of $`G`$ which generate the group. The walk proceeds by picking an ordered pair $`(i,j)`$ with $`1ijn`$ uniformly at random and applying one of the following four operations with equal probability:
$$R_{i,j}^\pm :(g_1,\mathrm{},g_i,\mathrm{},g_k)(g_1,\mathrm{},g_ig_j^\pm ,\mathrm{},g_k)$$
$$L_{i,j}^\pm :(g_1,\mathrm{},g_i,\mathrm{},g_k)(g_1,\mathrm{},g_j^\pm g_i,\mathrm{},g_k).$$
These moves map generating $`k`$-tuples to generating $`k`$-tuples. One starts from any generating $`k`$-tuple, applies the algorithm for $`r`$ steps, and then outputs a random entry of the resulting $`k`$-tuple (i.e. a group element).
The product replacement algorithm has superb practical performance (often converging more rapidly than random walk on the Cayley graph), in spite of the theoretical defects that a random entry of a random generating $`k`$-tuple does not have the same distribution as a random element of $`G`$, and that the convergence rate of the chain on $`k`$-tuples to its stationary distribution is unknown. The paper \[CeLgMuNiOb\], aware of these issues, tests the algorithm against theory, using conjugacy class statistics such as the order of an element, the number of factors of the characteristic polynomial of a random matrix, the degree of the largest irreducible factor of the characteristic polynomial of a random matrix, and the proportion of cyclic matrices in the finite classical groups. In short, understanding properties of random matrices is crucial to their analysis.
A recent effort to understand the performance of the product replacement algorithm uses Kazhdan’s property (T) from the representation theory of Lie groups \[LubPa\]; the paper \[Pa\] is a useful survey. Much remains to be done.
Example 9: Running times of algorithms
One of the main approaches to computing determinants and permanents of integer matrices involves doing the computations for reductions mod prime powers. Section 4.6.4 of \[Kn\] gives a detailed discussion with references to literature on upper bounds of running times. If one believes that typical matrices one encounters in the real world are like random matrices, this motivates studying random matrices over finite fields. In fact von Neumann’s interest in eigenvalues of random matrices with independent normal entries arose from the same heuristic applied to questions in numerical analysis (the introduction of \[E\] gives further discussion of this point).
Examples of algorithms in which properties of random matrices were really needed to bound running times include recognizing when a group generated by a set of matrices contains $`SL(n,q)`$ \[NP1\] and the MeatAxe algorithm for computing modular characters \[NP4\].
Example 10: Isometry classes of linear codes
Fripertinger \[Frip1\], \[Frip2\] considers cycle indices (in the permutation sense) of matrix groups acting on lines. His interest was in understanding properties of random isometry classes of linear codes–a harder problem than understanding random linear codes. The cycle indices he obtains seem quite intractable for theorem proving, but are useful in conjunction with computers. He also gives references to the switching function literature.
Curiously, understanding the permutation action of random matrices of lines comes up in another context. Wieand \[Wi\] has shown that the eigenvalues of random permutation matrices possess a structure similar to the eigenvalues of matrices from compact Lie groups. Persi Diaconis has suggested that the eigenvalues of high dimensional representations of finite groups of Lie type (such as the permutation action on lines) may possess similar structure; see \[F5\] for more in this direction.
### 2.3 Generalization to the Classical Groups
This subsection will focus on the finite unitary groups, with remarks about symplectic and orthogonal groups at the end. These cycle indices were derived in \[F1\],\[F2\] and were applied to the problem of estimating proportions of cyclic, separable, and semisimple matrices (these terms were defined in Subsection 2.2) in \[FNP\].
The unitary group $`U(n,q)`$ can be defined as the subgroup of $`GL(n,q^2)`$ preserving a non-degenerate skew-linear form. Recall that a skew-linear form on an $`n`$ dimensional vector space $`V`$ over $`F_{q^2}`$ is a bilinear map $`<,>:V\times VF_{q^2}`$ such that $`<\stackrel{}{x},\stackrel{}{y}>=<\stackrel{}{y},\stackrel{}{x}>^q`$ (raising to the $`q`$th power is an involution in a field of order $`q^2`$). One such form is given by $`<\stackrel{}{x},\stackrel{}{y}>=_{i=1}^nx_iy_i^q`$. Any two non-degenerate skew-linear forms are equivalent, so that $`U(n,q)`$ is unique up to isomorphism.
Wall \[W1\] parametrized the conjugacy classes of the finite unitary groups and computed their sizes. To describe his result, an involution on polynomials with non-zero constant term is needed. Given a polynomial $`\varphi `$ with coefficients in $`F_{q^2}`$ and non vanishing constant term, define a polynomial $`\stackrel{~}{\varphi }`$ by:
$$\stackrel{~}{\varphi }=\frac{z^{deg(\varphi )}\varphi ^q(\frac{1}{z})}{[\varphi (0)]^q}$$
where $`\varphi ^q`$ raises each coefficient of $`\varphi `$ to the $`q`$th power. Writing this out, a polynomial $`\varphi (z)=z^{deg(\varphi )}+\alpha _{deg(\varphi )1}z^{deg(\varphi )1}+\mathrm{}+\alpha _1z+\alpha _0`$ with $`\alpha _00`$ is sent to $`\stackrel{~}{\varphi }(z)=z^{deg(\varphi )}+(\frac{\alpha _1}{\alpha _0})^qz^{deg(\varphi )1}+\mathrm{}+(\frac{\alpha _{deg(\varphi )1}}{\alpha _0})^qz+(\frac{1}{\alpha _0})^q`$. An element $`\alpha U(n,q)`$ associates to each monic, non-constant, irreducible polynomial $`\varphi `$ over $`F_{q^2}`$ a partition $`\lambda _\varphi `$ of some non-negative integer $`|\lambda _\varphi |`$ by means of rational canonical form. The restrictions necessary for the data $`\lambda _\varphi `$ to represent a conjugacy class are that $`|\lambda _z|=0`$, $`\lambda _\varphi =\lambda _{\stackrel{~}{\varphi }}`$, and that $`_\varphi |\lambda _\varphi |deg(\varphi )=n.`$
Using formulas for conjugacy class sizes from \[W1\] together with some combinatorial manipulations, the following unitary group cycle index generating function was derived in \[F1\]. The products in the theorem are as always over monic irreducible polynomials.
###### Theorem 3
$`1+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{u^n}{|U(n,q)|}}{\displaystyle \underset{\alpha U(n,q)}{}}{\displaystyle \underset{\varphi :|\lambda _\varphi (\alpha )|>0}{}}x_{\varphi ,\lambda _\varphi (\alpha )}`$
$`=`$ $`{\displaystyle \underset{\varphi z,\varphi =\stackrel{~}{\varphi }}{}}\left[1+{\displaystyle \underset{n1}{}}{\displaystyle \underset{\lambda n}{}}x_{\varphi ,\lambda }{\displaystyle \frac{(u)^{ndeg(\varphi )}}{(q)^{deg(\varphi )_i(\lambda _i^{})^2}_{i1}\left(\frac{1}{(q)^{deg(\varphi )}}\right)_{m_i(\lambda )}}}\right]`$
$`{\displaystyle \underset{\{\varphi ,\stackrel{~}{\varphi }\},\varphi \stackrel{~}{\varphi }}{}}[1+{\displaystyle \underset{n1}{}}{\displaystyle \underset{\lambda n}{}}x_{\varphi ,\lambda }x_{\stackrel{~}{\varphi },\lambda }{\displaystyle \frac{u^{2ndeg(\varphi )}}{q^{2deg(\varphi )_i(\lambda _i^{})^2}_{i1}\left(\frac{1}{q^{2deg(\varphi )}}\right)_{m_i(\lambda )}}}]`$
One interesting theoretical result concerning the cycle index of $`U(n,q)`$ is the following functional equation. Letting $`C_{GL}(u,q)`$ and $`C_U(u,q)`$ be the cycle index generating functions for cyclic matrices in the general linear and unitary groups respectively, the functional equation states that
$$C_{GL}(u,q)C_U(u,q)=C_{GL}(u^2,q^2).$$
The paper \[FNP\] proves that this relation holds whenever the condition on the partitions $`\lambda _\varphi `$ is independent of the polynomial $`\varphi `$. In the current example, a matrix is cyclic if and only if all $`\lambda _\varphi `$ have at most one row. This condition is independent of $`\varphi `$.
Cycle indices for the symplectic and orthogonal groups are a bit trickier to establish from Wall’s formulas. To the treatment in \[F1\],\[F2\] we add a remark which should be very helpful to anyone trying to use those cycle indices. The paper \[F2\] only wrote out an explicit formula for the cycle index for the sum of $`+`$,$``$ type orthogonal groups. To solve for an individual orthogonal group, it is necessary to average that formula with a formula for the difference of $`+`$,$``$ type orthogonal groups (this procedure is carried out in a special case in \[FNP\]). In general, the formula for the difference of orthogonal groups is obtained from the formula for the sum of orthogonal groups as follows. First, for the polynomials $`z\pm 1`$, replace terms corresponding to partitions with an odd number of odd parts by their negatives. Second, for polynomials invariant under $`\stackrel{~}{}`$, replace terms corresponding to partitions of odd size by their negatives.
### 2.4 Limitations and Other Methods
Cycle index techniques, while very useful, also have their limitations and are not always the best way to proceed, as the following examples demonstrate.
Example 1: Primitive prime divisor elements
For integers $`b,e>1`$ a primitive prime divisor of $`b^e1`$ is a prime dividing $`b^e1`$ but not dividing $`b^i1`$ for any $`i`$ with $`1i<e`$. An element of $`GL(n,q)`$ is called a primitive prime divisor (ppd) element if its order is divisible by a primitive prime divisor of $`q^e1`$ with $`n/2<en`$. The analysis in \[NiP\] derives elegant bounds on the proportions of ppd elements in the finite classical groups and applies them to the group recognition problem for classical groups over finite fields (determining when a group generated by a set of matrices contains $`SL(n,q)`$). We do not see how to get comparable bounds using generating function techniques.
Example 2: Proportions of semisimple elements in exceptional groups
Although Example 5 of Section 2.2 was estimating proportions of semisimple matrices, this was only for the finite classical groups, where the index $`n`$ can take an infinite number of values. Cycle indices don’t seem useful unless there is a tower of groups of varying rank available.
Fortunately the computer package CHEVIE permits calculations precisely in finite rank cases such as the exceptional groups. Indeed this is how \[GuLub\] obtained estimates of the proportions of semisimple elements in the exceptional groups.
Example 3: Non-uniform distributions on matrices
The cycle indices give useful information about conjugacy class functions when the matrix is chosen uniformly at random. However there are other distributions on matrices which one could study and for which cycle index methods (at present) can not be applied.
One example is random $`n\times n`$ matrices where the matrix entries are chosen independently according to a given probability distribution on $`F_q`$. Charlap, Rees, and Robbins \[ChReRo\] show that if the probability distribution is not concentrated on any proper affine subspace of $`F_q`$, then as $`n\mathrm{}`$ the probability that the matrix is invertible is the same as for a uniform matrix. They use Moebius inversion on the lattice of subspaces of an $`n`$ dimensional vector space and the Poisson summation formula. Is the same true for other natural conjugacy class functions? We expect that the answer is yes, which can be regarded as a type of “universality” result for the asymptotic description of random elements of $`GL(n,q)`$ to be given in Subsection 3.1. Analogous universality results are known for matrices with complex entries \[So\]. For further information on the rank of random $`01`$ matrices, see \[BKW\] for sparse matrices, \[Bo\] for a survey of results on the rank over the real numbers, and also the discussion of work of Rudvalis and Shinoda in Subsection 3.2.
It is conceivable that cycle index techniques will be able to handle certain natural non-uniform distributions on $`GL(n,q)`$. This happens for the symmetric groups, where natural non-uniform measures such as performing a $`q`$-riffle shuffle on a deck of cards has a useful cycle index \[DiaMcPi\],\[F10\].
## 3 Running Example: General Linear Groups
The purpose of this section is to give different ways of understanding the conjugacy class of a random element of $`GL(n,q)`$. The analogous theory for other finite classical groups is mentioned in passing but is not treated in detail as many of the main ideas can be communicated using $`GL(n,q)`$. Subsection 3.1 will show how this leads naturally to the study of certain probability measures $`M_{GL,u,q}`$ on the set of all partitions of all natural numbers. Connections with symmetric function theory lead to several ways of growing random partitions distributed according to $`M_{GL,u,q}`$. One consequence is a motivated proof of the Rogers-Ramanujan identities.
### 3.1 Measures on Partitions
The goal is to obtain a probabilistic description of the conjugacy class of a random element of $`GL(n,q)`$. The ideas are based on \[F1\]. For this the following definition will be fundamental.
Definition: The measure $`M_{GL,u,q}`$ on the set of all partitions of all natural numbers is defined by
$$M_{GL,u,q}(\lambda )=\underset{r=1}{\overset{\mathrm{}}{}}(1\frac{u}{q^r})\frac{u^{|\lambda |}}{q^{_i(\lambda _i^{})^2}_i(\frac{1}{q})_{m_i(\lambda )}}.$$
The motivation for this definition will be clear from Theorem 4. The measure $`M_{GL,u,q}`$, while seemingly complicated, does have some nice combinatorial properties. For instance for partitions of a fixed size, this measure respects the dominance order on partitions (in this partial order $`\lambda \mu `$ if and only if $`\lambda _1+\mathrm{}\lambda _i\mu _1+\mathrm{}+\mu _i`$ for all $`i1`$). In work with Bob Guralnick we actually needed this property.
Lemma 4 proves that for $`q>1`$ and $`0<u<1`$, the measure $`M_{GL,u,q}`$ is in fact a probability measure. There are at least three other proofs of this fact: an argument using $`q`$ series, specializing an identity about Hall-Littlewood polynomials, or a slick argument using Markov chains and an identity of Cauchy. This third argument will be given in Subsection 3.4.
###### Lemma 4
If $`q>1`$ and $`0<u<1`$, then $`M_{GL,u,q}`$ defines a probability measure.
Proof: $`M_{GL,u,q}`$ is clearly non-negative when $`q>1`$ and $`0<u<1`$. Stong \[St1\] established an equation which is equivalent to the sought identity
$$\underset{\lambda }{}\frac{u^{|\lambda |}}{q^{_i(\lambda _i^{})^2}_i(\frac{1}{q})_{m_i(\lambda )}}=\underset{r=1}{\overset{\mathrm{}}{}}(\frac{1}{1\frac{u}{q^r}}).$$
As some effort is required to see this equivalence, we derive the identity directly using Stong’s line of reasoning.
First observe that unipotent elements of $`GL(n,q)`$ corresponding to nilpotent $`n\times n`$ matrices (subtract the identity matrix), and that the number of nilpotent $`n\times n`$ matrices is $`q^{n(n1)}`$ by the Fine-Herstein theorem \[FeinHer\]. The number of unipotent elements in $`GL(n,q)`$ can be evaluated in another way using the cycle index of the general linear groups. Namely set $`x_{\varphi ,\lambda }=1`$ if $`\varphi =z1`$ and set $`x_{\varphi ,\lambda }=0`$ otherwise. One concludes that
$$\underset{\lambda n}{}\frac{1}{q^{_i(\lambda _i^{})^2}_i(\frac{1}{q})_{m_i(\lambda )}}=\frac{q^{n(n1)}}{|GL(n,q)|}.$$
Now multiply both sides by $`u^n`$, sum in $`n`$, and apply Euler’s identity
$$\underset{n=0}{\overset{\mathrm{}}{}}\frac{u^nq^{\left(\genfrac{}{}{0pt}{}{n}{2}\right)}}{(q^n1)\mathrm{}(q1)}=\underset{r=1}{\overset{\mathrm{}}{}}(\frac{1}{1\frac{u}{q^r}}).$$
$`\mathrm{}`$
The measure $`M_{GL,u,q}`$ is a fundamental object for understanding the probability theory of conjugacy classes of $`GL(n,q)`$. This emerges from Theorem 4.
###### Theorem 4
1. Fix $`u`$ with $`0<u<1`$. Then choose a random natural number $`N`$ with probability of getting $`n`$ equal to $`(1u)u^n`$. Choose $`\alpha `$ uniformly in $`GL(N,q)`$. Then as $`\varphi `$ varies, the random partitions $`\lambda _\varphi (\alpha )`$ are independent random variables, with $`\lambda _\varphi `$ distributed according to the measure $`M_{GL,u^{deg(\varphi )},q^{deg(\varphi )}}`$.
2. Choose $`\alpha `$ uniformly in $`GL(n,q)`$. Then as $`n\mathrm{}`$, the random partitions $`\lambda _\varphi (\alpha )`$ converge in finite dimensional distribution to independent random variables, with $`\lambda _\varphi `$ distributed according to the measure $`M_{GL,1,q^{deg(\varphi )}}`$.
Proof: Recall the cycle index factorization
$$1+\underset{n=1}{\overset{\mathrm{}}{}}Z_{GL(n,q)}u^n=\underset{\varphi z}{}\left[1+\underset{n1}{}\underset{\lambda n}{}x_{\varphi ,\lambda }\frac{u^{ndeg(\varphi )}}{_\varphi q^{deg(\varphi )_i(\lambda _i^{})^2}_{i1}(\frac{1}{q^{deg(\varphi )}})_{m_i}}\right].$$
Setting all $`x_{\varphi ,\lambda }`$ equal to $`1`$ and using Lemma 4 shows that
$$\frac{1}{1u}=\underset{\varphi z}{}\underset{r=1}{\overset{\mathrm{}}{}}(\frac{1}{1\frac{u^{deg(\varphi )}}{q^{rdeg(\varphi )}}}).$$
Taking reciprocals and multiplying by the cycle index factorization shows that
$$(1u)+\underset{n=1}{\overset{\mathrm{}}{}}Z_{GL(n,q)}(1u)u^n=\underset{\varphi z}{}\left(M_{GL,u^{deg(\varphi )},q^{deg(\varphi )}}(\mathrm{})+\underset{\lambda :|\lambda |>0}{}M_{GL,u^{deg(\varphi )},q^{deg(\varphi )}}(\lambda )x_{\varphi ,\lambda }\right).$$
This proves the first assertion of the theorem. For the second assertion, use Lemma 2 from Subsection 2.2. $`\mathrm{}`$
Remarks:
1. Theorem 4 has an analog for the symmetric groups \[ShLl\]. The statement is as follows. Fix $`u`$ with $`0<u<1`$. Then choose a random natural number $`N`$ with probability of getting $`n`$ equal to $`(1u)u^n`$. Choose $`\pi `$ uniformly in $`S_N`$. Letting $`n_i`$ be the number of $`i`$-cycles of $`\pi `$, the random variables $`n_i`$ are independent, with $`n_i`$ distributed as a Poisson with mean $`\frac{u^i}{i}`$. Furthermore if one chooses $`\pi `$ uniformly in $`S_n`$ and lets $`n\mathrm{}`$, then the random variables $`n_i`$ are independent random variables, with $`n_i`$ distributed as a Poisson$`(\frac{1}{i})`$.
2. The idea of performing an auxiliary randomization of $`n`$ is a mainstay of statistical mechanics, known as the grand canonical ensemble. For a clear discussion see Sections 1.7, 1.9, and 4.3 of \[Fey\].
### 3.2 Symmetric Function Theory and Sampling Algorithms
The aim of this subsection is two-fold. First, the measures $`M_{GL,u,q}`$ are connected with the Hall-Littlewood symmetric functions. Then we indicate how this connection can be exploited to give probabilistic methods for growing random partitions distributed as $`M_{GL,u,q}`$. The purpose is not to drown the reader in formulas, but rather to show that the connection between symmetric functions and probability is deep, beautiful, and useful in both directions. The results on this section are based on \[F1\] and \[F3\], except for the remark on how to make the algorithms terminate in finite time, which is joint with Mark Huber.
To begin, we recall the Hall-Littlewood symmetric functions, which arise in many parts of mathematics: enumeration of $`p`$ groups, representation theory of $`GL(n,q)`$, and counting automorphisms of modules. The basic references for Hall-Littlewood polynomials $`P_\lambda `$ is Chapter 3 of \[Mac\], which offers the following definition
$$P_\lambda (x_1,\mathrm{},x_n;t)=\left[\frac{1}{_{i0}_{r=1}^{m_i(\lambda )}\frac{1t^r}{1t}}\right]\underset{wS_n}{}w\left(x_1^{\lambda _1}\mathrm{}x_n^{\lambda _n}\underset{i<j}{}\frac{x_itx_j}{x_ix_j}\right).$$
Here $`w`$ is a permutation acting on the $`x`$-variables by sending $`x_i`$ to $`x_{w(i)}`$. Recall that $`m_i(\lambda )`$ is the number of parts of $`\lambda `$ of size $`i`$. At first glance it is not obvious that these are polynomials, but the denominators cancel out after the symmetrization. The Hall-Littlewood polynomials interpolate between the Schur functions ($`t=0`$) and the monomial symmetric functions ($`t=1`$).
Theorem 5 relates the measures $`M_{GL,u,q}`$ to the Hall-Littlewood polynomials. Recall that $`n(\lambda )=_i(i1)\lambda _i=_i\left(\genfrac{}{}{0pt}{}{\lambda _i^{}}{2}\right)`$.
###### Theorem 5
$$M_{GL,u,q}(\lambda )=\underset{i=1}{\overset{\mathrm{}}{}}(1\frac{u}{q^i})\frac{P_\lambda (\frac{u}{q},\frac{u}{q^2},\mathrm{};\frac{1}{q})}{q^{n(\lambda )}}$$
Proof: From the above formula for Hall-Littlewood polynomials, it is clear that the only surviving term in the specialization $`P_\lambda (\frac{u}{q},\frac{u}{q^2},\mathrm{};\frac{1}{q})`$ is the term when $`w`$ is the identity. The rest is a simple combinatorial verification. (Alternatively, one could use “principal specialization” formulas for Macdonald polynomials on page 337 of \[Mac\]). $`\mathrm{}`$
Remark: The paper \[F3\] gives symmetric function theoretic generalizations of the measure $`M_{GL,u,q}`$ on partitions. In the case of Schur functions $`s_\lambda `$, this measure depends on two infinite sets of variables $`x_i,y_i`$ and assigns a partition $`\lambda `$ mass equal to $`s_\lambda (x_i)s_\lambda (y_i)_{i,j}(1x_iy_j)`$. It is remarkable that precisely this measure arose in work of the random matrix community relating the distribution of the lengths of increasing subsequences of random permutations to the distribution of eigenvalues of random GUE matrices (these matrices have complex entries). To elaborate, the Robinson-Schensted-Knuth correspondence associates a random partition of size $`n`$ to a random permutation of size $`n`$ and the shape of the partition encodes information about the longest increasing subsequence of the permutation. Choosing the size of the symmetric group randomly (according to a Poisson distribution) gives a probability measure on the set of all partitions of all natural numbers which is a special case of the above Schur function measure. Then the coordinate change $`h_j=\lambda _1^{}+\lambda _jj`$ maps the set of row lengths $`\{\lambda _j\}`$ of the partition to a set of distinct integers $`\{h_j\}`$. These $`h_j`$ can be viewed as positions of electrostatic charges repelling each other, and from this viewpoint the measure on subsets of the integers bears a striking resemblance to the eigenvalue density of a random GUE matrix. This fantastic heuristic can be made precise and led to a solution of the long-standing conjecture relating lengths of increasing subsequences of permutations to eigenvalues of random matrices. For these developments see \[BOOl\],\[Jo\] and the many references therein.
Now we return to the measure $`M_{GL,u,q}`$ and describe an algorithm for growing random partitions according to this measure.
The Young Tableau Algorithm
Start with $`N=1`$ and $`\lambda `$ the empty partition. Also start with a collection of coins indexed by the natural numbers, such that coin $`i`$ has probability $`\frac{u}{q^i}`$ of heads and probability $`1\frac{u}{q^i}`$ of tails.
Flip coin $`N`$.
If coin $`N`$ comes up tails, leave $`\lambda `$ unchanged, set $`N=N+1`$ and go to Step 1.
If coin $`N`$ comes up heads, choose an integer $`S>0`$ according to the following rule. Set $`S=1`$ with probability $`\frac{q^{N\lambda _1^{}}1}{q^N1}`$. Set $`S=s>1`$ with probability $`\frac{q^{N\lambda _s^{}}q^{N\lambda _{s1}^{}}}{q^N1}`$. Then increase the size of column $`s`$ of $`\lambda `$ by 1 and go to Step 1.
As an example of the Young Tableau Algorithm, suppose we are at Step 1 with $`\lambda `$ equal to the following partition:
$$\begin{array}{cccc}& & & \\ & & & \\ & & & \end{array}$$
Suppose also that $`N=4`$ and that coin 4 had already come up heads once, at which time we added to column 1, giving $`\lambda `$. We flip coin 4 again and get heads, going to Step 2b. We add a box to column $`1`$ with probability $`\frac{q1}{q^41}`$, to column $`2`$ with probability $`\frac{q^2q}{q^41}`$, to column $`3`$ with probability $`\frac{q^3q^2}{q^41}`$, to column $`4`$ with probability $`0`$, and to column $`5`$ with probability $`\frac{q^4q^3}{q^41}`$. We then return to Step 1.
###### Theorem 6
For $`0<u<1`$ and $`q>1`$, the Young Tableau Algorithm generates partitions which are distributed according to the measure $`M_{GL,u,q}`$.
To give insight into the proof of Theorem 6, we remark that it was deduced by proving a stronger result (Theorem 7) inductively and then taking the $`N\mathrm{}`$ limit. As is clear from the statement of Theorem 7, the connection with Hall-Littlewood polynomials (in particular the ability to truncate them) was crucial. It is unlikely that the Young Tableau Algorithm would have been discovered without this connection.
###### Theorem 7
Let $`P^N(\lambda )`$ be the probability that the algorithm outputs $`\lambda `$ when coin $`N`$ comes up tails. Then
$$P^N(\lambda )=\{\begin{array}{cc}\frac{u^{|\lambda |}(\frac{u}{q})_N(\frac{1}{q})_N}{(\frac{1}{q})_{N\lambda _1^{}}}\frac{P_\lambda (\frac{1}{q},\mathrm{},\frac{1}{q^N},0,\mathrm{};0,\frac{1}{q})}{q^{n(\lambda )}}\hfill & \text{if }\lambda _1^{}N\hfill \\ 0\hfill & \text{if }\lambda _1^{}>N.\hfill \end{array}$$
Next we explain why the Young Tableau Algorithm is called that. A standard Young tableau $`T`$ of size $`n`$ is a partition of $`n`$ with each box filled by one of $`\{1,\mathrm{},n\}`$ such that each of $`\{1,\mathrm{},n\}`$ appears exactly once and the numbers increase in each row and column of $`T`$. For instance,
$$\begin{array}{ccccc}\text{1}& \text{3}& \text{5}& \text{6}& \\ \text{2}& \text{4}& \text{7}& & \\ \text{8}& \text{9}& & & \end{array}$$
is a standard Young tableau. Standard Young tableaux are important in combinatorics and representation theory. The Young Tableau Algorithm is so named because numbering the boxes in the order in which they are created gives a standard Young tableau. Thus although our initial interest was in the measure $`M_{GL,u,q}`$ on partitions, the Young Tableau Algorithm yields more: a probability measure on standard Young tableaux. One consequence of this is a (new) representation of prinicipally specialized Hall-Littlewood polynomials as a sum of certain weights over standard Young tableaux.
Let us indicate an application of this probability measure on standard Young tableaux. Rudvalis and Shinoda \[RuShi\] studied the distribution of fixed vectors for the classical groups over finite fields. Let $`G=G(n)`$ be a classical group (i.e. one of $`GL`$,$`U`$,$`Sp`$, or $`O`$) acting on an $`n`$ dimensional vector space $`V`$ over a finite field $`F_q`$ (in the unitary case $`F_{q^2}`$) in its natural way. Let $`P_{G,n}(k,q)`$ be the chance that an element of $`G`$ fixes a $`k`$ dimensional subspace and let $`P_{G,\mathrm{}}(k,q)`$ be the $`n\mathrm{}`$ limit of $`P_{G,n}(k,q)`$. They found (in a 76 page unpublished work) beautiful formulas for $`P_{G,\mathrm{}}(k,q)`$. Their formulas are (setting $`x=\frac{1}{q}`$):
1. $`P_{GL,\mathrm{}}(k,q)=\left[_{r=1}^{\mathrm{}}(1x^r)\right]\frac{x^{k^2}}{(1x)^2\mathrm{}(1x^k)^2}`$
2. $`P_{U,\mathrm{}}(k,q)=\left[_{r=1}^{\mathrm{}}\frac{1}{1+x^{2r1}}\right]\frac{x^{k^2}}{(1x^2)\mathrm{}(1x^{2k})}`$
3. $`P_{Sp,\mathrm{}}(k,q)=\left[_{r=1}^{\mathrm{}}\frac{1}{1+x^r}\right]\frac{x^{\frac{k^2+k}{2}}}{(1x)\mathrm{}(1x^k)}`$
4. $`P_{O,\mathrm{}}(k,q)=\left[_{r=0}^{\mathrm{}}\frac{1}{1+x^r}\right]\frac{x^{\frac{k^2k}{2}}}{(1x)\mathrm{}(1x^k)}`$.
From a probabilistic perspective, it is very natural to try to interpret the factorizations in these formulas as certain random variables being independent (the paper \[RuShi\] gives no insight as to why these formulas have a product form). The Young tableau algorithm leads to such an understanding for the finite general linear and unitary groups; see \[F3\] for details.
Remarks
1. A skew diagram is the set theoretic difference between paritions $`\mu ,\lambda `$ with $`\mu \lambda `$ and a horizontal strip is a skew diagram with at most one square in each column. There is another algorithm for growing random partitions distributed according to $`M_{GL,u,q}`$ in which one tosses coins and adds horizontal strips (as opposed to a box at a time). Details are in \[F3\].
2. (Joint with Mark Huber) We indicate how to make the Young Tableau Algorithm run on a computer, so as to terminate in finite time (clearly one can’t flip infinitely many coins). Let $`a_N`$ be the number of times that coin $`N`$ comes up heads; the idea is to first determine the random vector $`(a_1,a_2,\mathrm{})`$ and then grow the partitions as in Step 2b of the Young Tableau Algorithm. So let us explain how to determine $`(a_1,a_2,\mathrm{})`$. For $`N1`$ let $`t^{(N)}`$ be the probability that all tosses of all coins numbered $`N`$ or greater are tails. For $`N1`$ and $`j0`$ let $`t_j^{(N)}`$ be the probability that some toss of a coin numbered $`N`$ or greater is a head and that coin $`N`$ comes up heads $`j`$ times. It is simple to write down expressions for $`t^{(N)},t_0^{(N)},t_1^{(N)},\mathrm{}`$ and clearly $`t^{(N)}+_{j0}t_j^{(N)}=1`$.
The basic operation a computer can perform is to produce a random variable $`U`$ distributed uniformly in the interval $`[0,1]`$. By dividing $`[0,1]`$ into intervals of length $`t^{(1)},t_0^{(1)},t_1^{(1)},\mathrm{}`$ and seeing where $`U`$ is located, one arrives at the value of $`a_1`$. Furthermore, if $`U`$ landed in the interval of length $`t^{(1)}`$ then all coins come up tails and the algorithm is over. Otherwise, move on to coin 2, dividing $`[0,1]`$ into intervals of length $`t^{(2)},t_0^{(2)},t_1^{(2)},\mathrm{}`$ and so on.
For $`0<u<1`$ and $`q`$ the size of a finite field, this algorithm terminates quickly. The probability of the algorithm stopping after the generation of the first uniform in $`[0,1]`$ is $`_{i=1}^{\mathrm{}}(1u/q^i)_{i=1}^{\mathrm{}}(11/q^i)>(11/q)^21/4`$ where the second inequality is Corollary 3.6 of \[NP2\]. Should it be necessary to generate future uniforms, the same argument shows that the algorithm stops after each one with probalility at least $`1/2`$.
### 3.3 Sampling for a Given Size: Unipotent Elements
An element of $`GL(n,q)`$ is called unipotent if all of its eigenvalues are $`1`$; a theorem of Steinberg asserts that the number of unipotent elements in $`GL(n,q)`$ is $`q^{n(n1)}`$ (this is the square of the order of a $`q`$-Sylow subgroup if $`q`$ is prime). Unipotent elements are interesting because any element $`\alpha `$ in $`GL(n,q)`$ can be written uniquely as the product $`\alpha _s\alpha _u`$ where $`\alpha _s`$ is semisimple and $`\alpha _u`$ is unipotent.
Thus it is natural to study the random partition $`\lambda _{z1}`$ for unipotent elements in $`GL(n,q)`$. This is the same as conditioning the measure $`M_{GL,u,q}`$ to live on partitions of size $`n`$. This subsection explains how to modify the sampling method of Subsection 3.2 to sample from this conditioned version of $`M_{GL,u,q}`$ and also from a $`q`$-analog of Plancharel measure (related to the longest increasing subsequence problem). These results are joint with Mark Huber.
Algorithm for Sampling from $`M_{GL,u,q}`$ given that $`|\lambda |=n`$
Start with $`N=1`$ and $`\lambda `$ the empty partition.
If $`n=0`$ then stop. Otherwise set $`h=1\frac{1}{q^n}`$.
Flip a coin with probability of heads $`h`$.
If the toss of Step 2 came up tails, increase the value of $`N`$ by $`1`$ and go to Step 2.
If the toss of Step 2 comes up heads, decrease the value of $`n`$ by $`1`$, increase $`\lambda `$ according to the rule of Step 2b of the Young Tableau Algorithm (which depends on $`N`$), and then go to Step 1.
Theorem 8 will show that the above algorithm samples from $`M_{GL,u,q}`$ conditioned to live on partitions of size $`n`$. It is perhaps surprising that unlike the Young Tableau Algorithm, the probability of a coin coming up heads is independent of the coin number; it depends only on the number of future boxes needed to get a partition of size $`n`$.
###### Lemma 5
Let $`N_i`$ be the number of times that coin $`i`$ comes up heads in the Young Tableau Algorithm with $`u=1`$ and let $`\stackrel{}{N_i}`$ be the infinite vector with $`i`$th component $`N_i`$.
1. The probability that $`\stackrel{}{N_i}=\stackrel{}{n_i}`$ is $`\frac{_{r=1}^{\mathrm{}}(1\frac{1}{q^i})}{q^{_iin_i}}`$.
2. $$\underset{\stackrel{}{n_i}:{\scriptscriptstyle n_i}=a}{}\frac{1}{q^{_iin_i}}=\frac{1}{q^a(\frac{1}{q})_a}.$$
Proof: The first assertion is clear. The second assertion is well known in the theory of partitions, but we argue probabilistically. Multiply both sides by $`_{r=1}^{\mathrm{}}(1\frac{1}{q^i})`$. Then note from the first assertion that the left hand side is the $`M_{GL,1,q}`$ chance of having a partition of size $`a`$. Now use the second equation in the proof of Lemma 4 in Subsection 3.1. $`\mathrm{}`$
For Theorem 8 the notation Prob. is shorthand for the probability of an event.
###### Theorem 8
The algorithm for sampling from $`M_{GL,u,q}`$ conditioned to live on partitions on size $`n`$ is valid.
Proof: From the formula for $`M_{GL,u,q}`$, the conditioned measure for $`M_{GL,u,q}`$ is the same as for $`M_{GL,1,q}`$. Now let $`n_i`$ be the number of times that coin $`i`$ comes up heads in the Young Tableau Algorithm. Letting $`|`$ denote conditioning, it suffices to show that
$$Prob.(n_i1|\underset{ji}{}n_j=s)=1\frac{1}{q^s}.$$
In fact (for reasons to be explained later) we compute a bit more, namely the conditional probability that $`n_i=a`$ given that $`_{ji}n_j=s`$. By definition this conditional probability is the ratio
$$\frac{Prob.(n_i=a,_{ji}n_j=s)}{Prob.(_{ji}n_j=s)}.$$
The numerator and denominator are computed using Lemma 6 as follows:
$`Prob.(n_i=a,{\displaystyle \underset{ji}{}}n_j=s)`$ $`=`$ $`{\displaystyle \underset{a_{i+1}+\mathrm{}=sa}{}}{\displaystyle \frac{_{r=i}^{\mathrm{}}(11/q^r)}{q^{ia}q^{_{ji+1}ja_j}}}`$
$`=`$ $`{\displaystyle \underset{a_{i+1}+\mathrm{}=sa}{}}{\displaystyle \frac{_{r=i}^{\mathrm{}}(11/q^r)}{q^{is}q^{_{ji+1}(ji)a_j}}}`$
$`=`$ $`{\displaystyle \underset{a_1+\mathrm{}=sa}{}}{\displaystyle \frac{_{r=i}^{\mathrm{}}(11/q^r)}{q^{is}q^{_{j1}a_j}}}`$
$`=`$ $`{\displaystyle \frac{_{r=i}^{\mathrm{}}(11/q^r)}{q^{is}q^{sa}(\frac{1}{q})_{sa}}}.`$
$`Prob.({\displaystyle \underset{ji}{}}n_j=s)`$ $`=`$ $`{\displaystyle \underset{a_i+\mathrm{}=s}{}}{\displaystyle \frac{_{r=i}^{\mathrm{}}(11/q^r)}{q^{_{ji}ja_j}}}`$
$`=`$ $`{\displaystyle \underset{a_i+\mathrm{}=s}{}}{\displaystyle \frac{_{r=i}^{\mathrm{}}(11/q^r)}{q^{(i1)s+_{ji}(j(i1))a_j}}}`$
$`=`$ $`{\displaystyle \underset{a_1+\mathrm{}=s}{}}{\displaystyle \frac{_{r=i}^{\mathrm{}}(11/q^r)}{q^{(i1)s+_{j1}ja_j}}}`$
$`=`$ $`{\displaystyle \frac{_{r=i}^{\mathrm{}}(11/q^r)}{q^{is}(\frac{1}{q})_s}}.`$
Thus $`Prob.(n_i=0|_{ji}n_j=s)=\frac{1}{q^s}`$ and the result follows. $`\mathrm{}`$
As mentioned in Subsection 3.2 there is a natural measure $`M_{Pl,q}`$ on the set of all partitions of all integers which when conditioned to live on partitions of a given size gives a $`q`$-analog of Plancherel measure, which is related to longest increasing subsequence in non-uniform random permutations \[F3\]. In what follows $`J_a(q)`$ is the polynomial discussed on pages 52-54 of \[F1\], $`h(s)`$ denotes the hook-length of a dot in $`\lambda `$ \[Mac\] and $`[n]=\frac{q^n1}{q1}`$ is the $`q`$-analog of the number $`n`$. Recall that a skew diagram is the set theoretic difference between paritions $`\mu ,\lambda `$ with $`\mu \lambda `$ and that a horizontal strip is a skew diagram with at most one square in each column.
Algorithm for Sampling from $`M_{Pl,q}`$ for $`q>1`$ given that $`|\lambda |=n`$
Start with $`\lambda `$ the empty partition.
If $`n=0`$ then stop. Otherwise choose $`a`$ with $`0an`$ with probability
$$\frac{q^{n^2}(1\frac{1}{q^{na+1}})^2\mathrm{}(1\frac{1}{q^n})^2}{q^{(na)^2+n}(\frac{1}{q})_a}\frac{J_{na}(q)}{J_n(q)}.$$
Then increase $`\lambda `$ to $`\mathrm{\Lambda }`$ with probability
$$(1\frac{1}{q})\mathrm{}(1\frac{1}{q^a})\frac{q^{n(\lambda )}_{s\lambda }(1\frac{1}{q^{h(s)}})}{q^{n(\mathrm{\Lambda })}_{s\mathrm{\Lambda }}(1\frac{1}{q^{h(s)}})}$$
if $`\mathrm{\Lambda }\lambda `$ is a horizontal strip of size $`a`$ and with probability $`0`$ otherwise. Finally replace $`n`$ by $`na`$ and repeat Step 1.
Using Lemma 6, Theorem 9 proves that the algorithm for sampling from $`M_{Pl,q}`$ conditioned to live on $`|\lambda |=n`$ works. We omit the details, which (given the background material in \[F1\]) are analogous to the case of $`M_{GL,u,q}`$.
###### Lemma 6
Let $`N_i`$ be the number of times that coin $`i`$ comes up heads in the algorithm from \[F3\] for sampling from the measure $`M_{Pl,q}`$ and let $`\stackrel{}{N_i}`$ be the infinite vector with $`i`$th component $`N_i`$.
1. The probability that $`\stackrel{}{N_i}=\stackrel{}{n_i}`$ is $`\frac{_{r=1}^{\mathrm{}}_{j=r}^{\mathrm{}}(1\frac{1}{q^j})}{q^{_iin_i}_i(\frac{1}{q})_{n_i}}`$.
2. $$\underset{\stackrel{}{n_i}:_{n_i}=a}{}\frac{1}{q^{_iin_i}_i(\frac{1}{q})_{n_i}}=\frac{J_a(q)}{q^{a^2}(1\frac{1}{q})^2\mathrm{}(1\frac{1}{q^a})^2}.$$
###### Theorem 9
The algorithm given for sampling from $`M_{Pl,q}`$ with $`q>1`$ conditioned to live on $`|\lambda |=n`$ is valid.
### 3.4 Markov Chain Approach
The main result in this subsection is a third method for understanding the measure $`M_{GL,u,q}`$ probabilistically (\[F7\]). The idea is to build up the random partition a column at a time; if the current column has size $`a`$, then the next column will have size $`b`$ (with $`ba`$) with probability $`K(a,b)`$. The surprise is that this transition rule turns out to be independent of the columns, yielding a Markov on the natural numbers. This Markov chain is diagonalizable with eigenvalues $`1,\frac{u}{q},\frac{u^2}{q^4},\mathrm{}`$. It will be used to give a probabilistic proof of the Rogers-Ramanujan identities in Subsection 3.5.
It is convenient to set $`\lambda _0^{}`$ (the height of an imaginary zeroth column) equal to $`\mathrm{}`$. For the entirety of this subsection, let $`P(a)`$ be the $`M_{GL,u,q}`$ probability that $`\lambda _1^{}=a`$. Theorem 10, which makes the connection with Markov chains, is proved in a completely elementary way. The argument reproves that $`M_{GL,u,q}`$ is a probability measure (Lemma 4 of Subsection 3.1), shows that the asserted Markov transition probabilities add to one, and gives a formula for $`P(a)`$.
###### Theorem 10
Starting with $`\lambda _0^{}=\mathrm{}`$, define in succession $`\lambda _1^{},\lambda _2^{},\mathrm{}`$ according to the rule that if $`\lambda _i^{}=a`$, then $`\lambda _{i+1}^{}=b`$ with probability
$$K(a,b)=\frac{u^b(\frac{1}{q})_a(\frac{u}{q})_a}{q^{b^2}(\frac{1}{q})_{ab}(\frac{1}{q})_b(\frac{u}{q})_b}.$$
Then the resulting partition is distributed according to $`M_{GL,u,q}`$.
Proof: Suppose we know that $`M_{GL,u,q}`$ is a probability measure and that
$$P(a)=\frac{u^a(\frac{u}{q})_{\mathrm{}}}{q^{a^2}(\frac{1}{q})_a(\frac{u}{q})_a}.$$
Then the $`M_{GL,u,q}`$ probability of choosing a partition with $`\lambda _i^{}=r_i^{}`$ for all $`i`$ is
$$Prob.(\lambda _0^{}=\mathrm{})\frac{Prob.(\lambda _0^{}=\mathrm{},\lambda _1^{}=r_1)}{Prob.(\lambda _0^{}=\mathrm{})}\underset{i=1}{\overset{\mathrm{}}{}}\frac{Prob.(\lambda _0^{}=\mathrm{},\lambda _1^{}=r_1,\mathrm{},\lambda _{i+1}^{}=r_{i+1})}{Prob.(\lambda _0^{}=\mathrm{},\lambda _1^{}=r_1,\mathrm{},\lambda _i^{}=r_i)}.$$
Thus it is enough to prove the (surprising) assertion that
$$\frac{Prob.(\lambda _0^{}=\mathrm{},\lambda _1^{}=r_1,\mathrm{},\lambda _{i1}^{}=r_{i1},\lambda _i^{}=a,\lambda _{i+1}^{}=b)}{Prob.(\lambda _0^{}=\mathrm{},\lambda _1^{}=r_1,\mathrm{},\lambda _{i1}^{}=r_{i1},\lambda _i^{}=a)}=\frac{u^b(\frac{1}{q})_a(\frac{u}{q})_a}{q^{b^2}(\frac{1}{q})_{ab}(\frac{1}{q})_b(\frac{u}{q})_b},$$
for all $`i,a,b,r_1,\mathrm{},r_{i1}`$. One calculates that
$$\underset{\genfrac{}{}{0pt}{}{\lambda :\lambda _1^{}=r_1,\mathrm{},\lambda _{i1}^{}=r_{i1}}{\lambda _i^{}=a}}{}M_{GL,u,q}(\lambda )=\frac{u^{r_1+\mathrm{}+r_{i1}}}{q^{r_1^2+\mathrm{}+r_{i1}^2}(\frac{1}{q})_{r_1r_2}\mathrm{}(\frac{1}{q})_{r_{i2}r_{i1}}(\frac{1}{q})_{r_{i1}a}}P(a).$$
Similarly, observe that
$$\underset{\genfrac{}{}{0pt}{}{\lambda :\lambda _1^{}=r_1,\mathrm{},\lambda _{i1}^{}=r_{i1}}{\lambda _i^{}=a,\lambda _{i+1}^{}=b}}{}M_{GL,u,q}(\lambda )=\frac{u^{r_1+\mathrm{}+r_{i1}+a}}{q^{r_1^2+\mathrm{}+r_{i1}^2+a^2}(\frac{1}{q})_{r_1r_2}\mathrm{}(\frac{1}{q})_{r_{i2}r_{i1}}(\frac{1}{q})_{r_{i1}a}(\frac{1}{q})_{ab}}P(b).$$
Thus the ratio of these two expressions is
$$\frac{u^b(\frac{1}{q})_a(\frac{u}{q})_a}{q^{b^2}(\frac{1}{q})_{ab}(\frac{1}{q})_b(\frac{u}{q})_b},$$
as desired. Note that the transition probabilities must sum to 1 because
$$\underset{ba}{}\frac{_{\genfrac{}{}{0pt}{}{\lambda :\lambda _1^{}=r_1,\mathrm{},\lambda _{i1}^{}=r_{i1}}{\lambda _i^{}=a,\lambda _{i+1}^{}=b}}M_{GL,u,q}(\lambda )}{_{\genfrac{}{}{0pt}{}{\lambda :\lambda _1^{}=r_1,\mathrm{},\lambda _{i1}^{}=r_{i1}}{\lambda _i^{}=a}}M_{GL,u,q}(\lambda )}=1$$
for any measure $`M_{GL,u,q}`$ on partitions.
Thus to complete the proof, it must be shown that $`M_{GL,u,q}`$ is a probability measure and that
$$P(a)=\frac{u^a(\frac{u}{q})_{\mathrm{}}}{q^{a^2}(\frac{1}{q})_a(\frac{u}{q})_a}.$$
Since
$$\frac{_{\genfrac{}{}{0pt}{}{\lambda :\lambda _1^{}=r_1,\mathrm{},\lambda _{i1}^{}=r_{i1}}{\lambda _i^{}=a,\lambda _{i+1}^{}=b}}M_{GL,u,q}(\lambda )}{_{\genfrac{}{}{0pt}{}{\lambda :\lambda _1^{}=r_1,\mathrm{},\lambda _{i1}^{}=r_{i1}}{\lambda _i^{}=a}}M_{GL,u,q}(\lambda )}=\frac{P(b)u^a}{P(a)q^{a^2}(\frac{1}{q})_{ab}}$$
it follows that
$$\underset{ba}{}\frac{P(b)u^a}{P(a)q^{a^2}(\frac{1}{q})_{ab}}=1.$$
From this recursion and the fact that $`P(0)=(\frac{u}{q})_{\mathrm{}}`$, one solves for $`P(a)`$ inductively, finding that
$$P(a)=\frac{u^a(\frac{u}{q})_{\mathrm{}}}{q^{a^2}(\frac{1}{q})_a(\frac{u}{q})_a}.$$
Cauchy’s identity (page 20 of \[A1\]) gives that $`_aP(a)=1`$, so that $`M_{GL,u,q}`$ is a probability measure. $`\mathrm{}`$
Theorem 11 diagonalizes the transition matrix $`K`$, finding a basis of eigenvectors, which is fundamental for understanding the Markov chain (part 3 is stated as a Lemma in \[A2\]). Since the matrix $`K`$ is upper triangular with distinct eigenvalues, this is straightforward.
###### Theorem 11
1. Let $`C`$ be the diagonal matrix with $`(i,i)`$ entry $`(\frac{1}{q})_i(\frac{u}{q})_i`$. Let $`M`$ be the matrix $`\left(\frac{u^j}{q^{j^2}(\frac{1}{q})_{ij}}\right)`$. Then $`K=CMC^1`$, which reduces the problem of diagonalizing $`K`$ to that of diagonalizing $`M`$.
2. Let $`A`$ be the matrix $`\left(\frac{1}{(\frac{1}{q})_{ij}(\frac{u}{q})_{i+j}}\right)`$. Then the columns of $`A`$ are eigenvectors of $`M`$ for right multiplication, the $`j`$th column having eigenvalue $`\frac{u^j}{q^{j^2}}`$.
3. The inverse matrix $`A^1`$ is $`\left(\frac{(1u/q^{2i})(1)^{ij}(\frac{u}{q})_{i+j1}}{q^{\left(\genfrac{}{}{0pt}{}{ij}{2}\right)}(\frac{1}{q})_{ij}}\right)`$.
Corollary 2 (immediate from Theorem 11) will be useful for the proof of the Rogers-Ramanujan identities in Section 3.5. In the case $`L\mathrm{}`$ and $`j=0`$, it is the so called Rogers-Selberg identity.
###### Corollary 2
Let $`E`$ be the diagonal matrix with $`(i,i)`$ entry $`\frac{u^i}{q^{i^2}}`$. Then $`K^r=CAE^rA^1C^1`$. More explicitly,
$$K^r(L,j)=\frac{(\frac{1}{q})_L(\frac{u}{q})_L}{(\frac{1}{q})_j(\frac{u}{q})_j}\underset{n=0}{\overset{\mathrm{}}{}}\frac{u^{rn}(1u/q^{2n})(1)^{nj}(\frac{u}{q})_{n+j1}}{q^{rn^2}(\frac{1}{q})_{Ln}(\frac{u}{q})_{L+n}q^{\left(\genfrac{}{}{0pt}{}{nj}{2}\right)}(\frac{1}{q})_{nj}}.$$
Proof: This is immediate from Theorem 11. $`\mathrm{}`$
Remarks
1. One of our motivations for seeking a Markov chain description of $`M_{GL,u,q}`$ is work of Fristedt \[Fris\], who had a Markov chain approach for the measure $`P_q`$ on the set of all partitions of all natural numbers defined by $`P_q(\lambda )=_{i=1}^{\mathrm{}}(1q^i)q^{|\lambda |}`$ where $`q<1`$. Fristedt’s interest was in studying what a uniformly chosen partition of an integer looks like, and conditioning $`P_q`$ to live on partitions of size $`n`$ gives a uniform partition. The measure $`P_q`$ is related to to vertex operators \[O1\] and to the enumeration of ramified coverings of the torus \[Dij\]. In this regard the papers \[O1\] and \[BlO\] prove that the $`k`$ point correlation function
$$F(t_1,\mathrm{},t_k)=\underset{\lambda }{}q^{|\lambda |}\underset{k=1}{\overset{n}{}}\underset{i=1}{\overset{\mathrm{}}{}}t_k^{\lambda _ii+\frac{1}{2}}$$
is a sum of determinants involving genus 1 theta functions and their derivatives and give connections with quasi-modular forms. It would be marvellous if the measure $`M_{GL,u,q}`$ (being related to modular forms via the Rogers-Ramanujan identities) is also related to enumerative questions in algebraic geometry.
2. As mentioned in the introduction, the Markov chain approach gives a unified description of conjugacy classes of the finite classical groups. For the symplectic and orthogonal groups it is necessary to use two Markov chains $`K_1`$ and $`K_2`$. For the symplectic case, steps with column number $`i`$ odd use $`K_1`$ and steps with column number $`i`$ even use $`K_2`$. For the orthogonal case, steps with column number $`i`$ odd use $`K_2`$ and steps with column number $`i`$ even use $`K_1`$. The Markov chains $`K_1,K_2`$ are the same for both cases! Details are in \[F6\]. The Markov chain approach is also related to quivers \[F7\].
### 3.5 Rogers-Ramanujan Identities
The Rogers-Ramanujan identities \[Ro\]
$$1+\underset{n=1}{\overset{\mathrm{}}{}}\frac{q^{n^2}}{(1q)(1q^2)\mathrm{}(1q^n)}=\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{(1q^{5n1})(1q^{5n4})}$$
$$1+\underset{n=1}{\overset{\mathrm{}}{}}\frac{q^{n(n+1)}}{(1q)(1q^2)\mathrm{}(1q^n)}=\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{(1q^{5n2})(1q^{5n3})}$$
are among the most interesting partition identities in number theory and combinatorics, with connections to Lie theory and statistical mechanics (see the discussions in \[A2\] and \[F7\] for many references). One ongoing challenge in the subject (posed by Hardy) has been to find a proof of the Rogers-Ramanujan identities which is both motivated and simple. The purpose of this subsection is to describe such a proof (\[F7\]), which is also the first probabilistic proof of the Rogers-Ramanujan identities.
To illustrate the idea we give the proof of the following generalization of the first Rogers-Ramanujan identity (called the Andrews-Gordon identity \[A3\],\[Gor\]):
$$\underset{n_1,\mathrm{},n_{k1}0}{}\frac{1}{q^{N_1^2+\mathrm{}+N_{k1}^2}(1/q)_{n_1}\mathrm{}(1/q)_{n_{k1}}}=\underset{\genfrac{}{}{0pt}{}{r=1}{r0,\pm k(mod2k+1)}}{\overset{\mathrm{}}{}}\frac{1}{1(1/q)^r}$$
where $`N_i=n_i+\mathrm{}+n_{k1}`$.
The idea is simple. We study the distribution of the length of the first row of a random partition distributed as $`M_{GL,1,q}`$. From the definition of $`M_{GL,1,q}`$ the probability that the first row has length less than $`k`$ is equal to
$$\underset{r=1}{\overset{\mathrm{}}{}}(1\frac{1}{q^r})\underset{\lambda :\lambda _k^{}=0}{}\frac{1}{q^{(\lambda _1^{})^2+\mathrm{}+(\lambda _{k1}^{})^2}(1/q)_{\lambda _1^{}\lambda _2^{}}\mathrm{}(1/q)_{\lambda _{k1}^{}\lambda _k^{}}}.$$
Letting $`n_i`$ denote $`\lambda _i^{}\lambda _{i+1}^{}`$ and $`N_i`$ denote $`\lambda _i^{}`$, this becomes
$$\underset{r=1}{\overset{\mathrm{}}{}}(1\frac{1}{q^r})\underset{n_1,\mathrm{},n_{k1}0}{}\frac{1}{q^{N_1^2+\mathrm{}+N_{k1}^2}(1/q)_{n_1}\mathrm{}(1/q)_{n_{k1}}}$$
which is a essentially the left hand side of the Andrews-Gordon identity. On the other hand the probability that the first row has length less than $`k`$ is equal to the probability that the Markov chain of Section 3.4 is absorbed at $`0`$ at time $`k`$. Since we diagonalized the matrix associated to this Markov chain, it is straightforward to compute this probability. To get it into product form it is necessary to apply Jacobi’s triple product identity which has a simple combinatorial proof \[A1\]. Further details are in \[F7\].
Next we argue that this proof is motivated. Certainly the measure $`M_{GL,1,q}`$ is a natural object to study, given that it is the $`n\mathrm{}`$ limit law of $`\lambda _{z1}`$ for a random element of $`GL(n,q)`$. It was natural to try to build up the random partitions $`\lambda `$ column by column as in Section 3.4. Observing that the resulting Markov chain is absorbing at $`0`$ with probability one, the time to absorption (equivalent to the distribution of the length of the first row) is the most natural quantity one could examine. The final step is applying Jacobi’s triple product identity, and thus going from a “sum = sum” identity to a “sum = product” identity. As mentioned above Jacobi’s triple product identity is easy to verify, but one still wants a motivation for trying to write the left hand side of the Andrews-Gordon identity in product form. One motivation is Baxter’s work on statistical mechanics (surveyed in \[A2\],\[Bax1\],\[Bax2\]) in which he really needed “sum = product” identities and was led to conjecture analogs of Rogers-Ramanujan type identities. Although a proof of the Rogers-Ramanujan identities doesn’t emerge from his work, it is clearly one of the truly great accomplishments in mathematics and his book \[Bax1\] has been very influential. A second motivation is our work on the $`n\mathrm{}`$ asymptotic probability that an element of $`GL(n,q)`$ is semisimple. The argument, recorded in \[F1\] or the more readily available \[F4\] needed a “sum = product” identity. The corresponding computation in \[F9\] for the finite affine groups needed both Rogers-Ramanujan identities.
Andrews’ paper \[A4\] notes that many proofs of the Rogers-Ramanujan identities make use of the following result called Bailey’s Lemma, alluded to in \[Bai\] and stated explicity in \[A3\]. A pair of sequences $`\{\alpha _L\}`$ and $`\{\beta _L\}`$ are called a Bailey pair if
$$\beta _L=\underset{r=0}{\overset{L}{}}\frac{\alpha _r}{(1/q)_{Lr}(u/q)_{L+r}}.$$
Bailey’s Lemma states that if $`\alpha _L^{}=\frac{u^L}{q^{L^2}}\alpha _L`$ and $`\beta _L^{}=_{r=0}^L\frac{u^r}{q^{r^2}(1/q)_{Lr}}\beta _r`$, then $`\{\alpha _L^{}\}`$ and $`\{\beta _L^{}\}`$ are a Bailey pair. From the viewpoint of Markov chains, this case of Bailey’s Lemma is clear. To explain, let $`A,D,M`$ be as in Theorem 11 (recall that $`M=ADA^1`$). Viewing $`\alpha =\stackrel{}{\alpha _L}`$ and $`\beta =\stackrel{}{\beta _L}`$ as column vectors, the notion of a Bailey pair means that $`\beta =A\alpha `$. This case of Bailey’s Lemma follows because
$$\beta ^{}=M\beta =ADA^1\beta =AD\alpha =A\alpha ^{}.$$
As Andrews explains in \[A2\], the power of Bailey’s lemma lies in its ability to be iterated and gives a short proof of the Rogers-Selberg identity (Corollary 2 in Section 3.4). From the remarks in this paragraph it is clear that iterating Bailey’s lemma corresponds to taking several according to the Markov chain $`K`$. This demystifies the Bailey’s Lemma proofs of the Rogers-Ramanujan identities, which strike this author as unmotivated. The fact that the Markov chain approach has analogs for other finite classical groups and for quivers is further evidence of its naturality.
## 4 Upper Triangular Matrices
This section surveys probabilistic aspects of conjugacy classes in the group $`T(n,q)`$ of upper triangular matrices over finite fields with $`1`$’s along the main diagonal. At present little is known about conjugacy in $`T(n,q)`$. The papers \[VAr\], \[VArV\] study the number of conjugacy classes. Kirillov \[Kir\] calls for an extension of his method of coadjoint orbits for groups over real, complex, or $`p`$-adic fields to the group $`T(n,q)`$ and gives premilinary connections with statistical physics; the paper \[IsKar\] gives a counterexample to one of his conjectures. As we do not see how to further develop those results or improve on their exposition, we instead focus on a simpler problem: the probabilistic study of Jordan form of elements of $`T(n,q)`$.
Subsection 4.1 describes a probabilistic growth algorithm for the Jordan form of upper triangular matrices over a finite field. This is linked with symmetric function theory and potential theory on Bratteli diagrams in Subsection 4.2.
### 4.1 Growth Algorithm for Jordan Form
Theorem 12 gives a probabilistic growth algorithm for the Jordan form of random elements of $`T(n,q)`$. Its proof uses elementary reasoning from linear algebra.
###### Theorem 12
(\[Kir\],\[B\]) The Jordan form of a uniformly chosen element of $`T(n,q)`$ can be sampled from by stopping the following procedure after $`n`$ steps:
Starting with the empty partition, at each step transition from a partition $`\lambda `$ to a partition $`\mathrm{\Lambda }`$ by adding a box to column $`i`$ chosen according to the rules
* $`i=1`$ with probability $`\frac{1}{q^{\lambda _1^{}}}`$
* $`i=j>1`$ with probability $`\frac{1}{q^{\lambda _j^{}}}\frac{1}{q^{\lambda _{j1}^{}}}`$
Theorem 12 leads to the following central limit theorem about the asymptotic Jordan form of an element of $`T(n,q)`$.
###### Theorem 13
(\[B\]) Let $`\lambda `$ be the partition corresponding to the Jordan form of a random element of $`T(n,q)`$. Let $`Prob^n`$ denote probability under the uniform measure on $`T(n,q)`$ and let $`p_i=\frac{1}{q^{i1}}\frac{1}{q^i}`$. Then
$$lim_n\mathrm{}Prob^n(\frac{\lambda _ip_in}{\sqrt{n}}x_i,i=1,\mathrm{},k)=(2\pi )^{\frac{k}{2}}_{\mathrm{}}^{x_1}\mathrm{}_{\mathrm{}}^{x_k}e^{\frac{1}{2}<Qt,t>}𝑑t$$
for any $`(x_1,\mathrm{},x_k)R^k`$, where the covariance matrix equals
$$Q=diag(p_1,\mathrm{},p_k)(p_ip_j)_{i,j=1}^k.$$
### 4.2 Symmetric Functions and Potential Theory
Given the usefulness of symmetric functions in the probabilistic study of the measure $`M_{GL,u,q}`$, it is natural to seek an analogous understanding of Theorem 12. That is the topic of the present subsection. The ideas here are from the report \[F8\].
The first step is to link the probability that an element of $`T(n,q)`$ has Jordan form of type $`\mathrm{\Lambda }`$ with symmetric function theory. For the rest of this section, $`P_\mathrm{\Lambda }(q,t)`$ denotes a Macdonald polynomial, $`K_{\mu \mathrm{\Lambda }}(q,t)`$ denotes a Kostka-Foulkes polynomial, and $`f^\mu `$ is the dimension of the irreducible representation of $`S_n`$ corresponding to the partition $`\mu `$ (see \[Mac\] for background). Note that when $`q=0`$ the Macdonald polynomial is our friend, a Hall-Littlewood polynomial.
###### Theorem 14
(\[F5\]) The probability that a random element of $`T(n,q)`$ has Jordan form of type $`\mathrm{\Lambda }`$ is
$$P_\mathrm{\Lambda }(1\frac{1}{q},\frac{1}{q}\frac{1}{q^2},\mathrm{};0,\frac{1}{q})\underset{\mu n}{}f^\mu K_{\mu \mathrm{\Lambda }}(0,q).$$
Next we give some background on potential theory on Bratteli diagrams. This is a beautiful subject, with connections to probability and representation theory. We recommend \[Ke1\] for background on potential theory with many examples and \[BOl\] for a survey of recent developments. The basic set-up is as follows. One starts with a Bratteli diagram; that is an oriented graded graph $`\mathrm{\Gamma }=_{n0}\mathrm{\Gamma }_n`$ such that
1. $`\mathrm{\Gamma }_0`$ is a single vertex $`\mathrm{}`$.
2. If the starting vertex of an edge is in $`\mathrm{\Gamma }_i`$, then its end vertex is in $`\mathrm{\Gamma }_{i+1}`$.
3. Every vertex has at least one outgoing edge.
4. All $`\mathrm{\Gamma }_i`$ are finite.
For two vertices $`\lambda ,\mathrm{\Lambda }\mathrm{\Gamma }`$, one writes $`\lambda \mathrm{\Lambda }`$ if there is an edge from $`\lambda `$ to $`\mathrm{\Lambda }`$. Part of the underlying data is a multiplicity function $`\kappa (\lambda ,\mathrm{\Lambda })`$. Letting the weight of a path in $`\mathrm{\Gamma }`$ be the product of the multiplicities of its edges, one defines the dimension $`dim(\mathrm{\Lambda })`$ of a vertex $`\mathrm{\Lambda }`$ to be the sum of the weights over all maximal length paths from $`\mathrm{}`$ to $`\mathrm{\Lambda }`$ (this definition clearly extend to intervals). Given a Bratteli diagram with a multiplicity function, one calls a function $`\varphi `$ harmonic if $`\varphi (0)=1`$, $`\varphi (\lambda )0`$ for all $`\lambda \mathrm{\Gamma }`$, and
$$\varphi (\lambda )=\underset{\mathrm{\Lambda }:\lambda \mathrm{\Lambda }}{}\kappa (\lambda ,\mathrm{\Lambda })\varphi (\mathrm{\Lambda }).$$
An equivalent concept is that of coherent probability distributions. Namely a set $`\{M_n\}`$ of probability distributions $`M_n`$ on $`\mathrm{\Gamma }_n`$ is called coherent if
$$M_{n1}(\lambda )=\underset{\mathrm{\Lambda }:\lambda \mathrm{\Lambda }}{}\frac{dim(\lambda )\kappa (\lambda ,\mathrm{\Lambda })}{dim(\mathrm{\Lambda })}M_n(\mathrm{\Lambda }).$$
The formula allowing one to move between the definitions is $`\varphi (\lambda )=\frac{M_n(\lambda )}{dim(\lambda )}`$.
One reason the set-up is interesting from the viewpoint of probability theory is the fact that every harmonic function can be written as a Poisson integral over the set of extreme harmonic functions (which is often the Martin boundary). For the Pascal lattice (vertices of $`\mathrm{\Gamma }_n`$ are pairs $`(k,n)`$ with $`k=0,1,\mathrm{},n`$ and $`(k,n)`$ is connected to $`(k,n+1)`$ and $`(k+1,n+1)`$), this fact is the simplest instance of de Finetti’s theorem. When the multiplicity function $`\kappa `$ is integer valued, one can define a sequence of algebras $`A_n`$ associated to the Bratteli diagram, and harmonic functions correspond to certain characters of the inductive limit of the algebras $`A_n`$.
Next we define a branching for which the probability that an element of $`T(n,q)`$ has Jordan type $`\mathrm{\Lambda }`$ is a harmonic function. First some notation is needed. For $`\lambda \mathrm{\Lambda }`$, let $`R_{\mathrm{\Lambda }/\lambda }`$ (resp. $`C_{\mathrm{\Lambda }/\lambda }`$) be the boxes of $`\lambda `$ in the same row (resp. column) as the boxes removed from $`\lambda `$ to get $`\mathrm{\Lambda }`$. This notation differs from that in \[Mac\]. Let $`a_\lambda (s)`$, $`l_\lambda (s)`$ be the number of dots in $`\lambda `$ strictly to the east and south of $`s`$, and let $`h_\lambda (s)=a_\lambda (s)+l_\lambda (s)+1`$.
Definition 1: For $`0q<1`$ and $`0<t<1`$, the underlying Bratteli diagram $`\mathrm{\Gamma }`$ has as level $`\mathrm{\Gamma }_n`$ all partitions $`\lambda `$ of $`n`$. Letting $`i`$ be the column number of the dot removed to go from $`\lambda `$ to $`\mathrm{\Lambda }`$, for $`\lambda \mathrm{\Lambda }`$, define the multiplicty function as
$$\kappa (\lambda ,\mathrm{\Lambda })=\frac{1}{t^{\mathrm{\Lambda }_i^{}1}}\underset{sR_{\mathrm{\Lambda }/\lambda }}{}\frac{1q^{a_\mathrm{\Lambda }(s)+1}t^{l_\mathrm{\Lambda }(s)}}{1q^{a_\lambda (s)+1}t^{l_\lambda (s)}}\underset{sC_{\mathrm{\Lambda }/\lambda }}{}\frac{1q^{a_\mathrm{\Lambda }(s)}t^{l_\mathrm{\Lambda }(s)+1}}{1q^{a_\lambda (s)}t^{l_\lambda (s)+1}}.$$
Equation I.10 of \[GarsH\] proves that
$$dim(\mathrm{\Lambda })=\frac{1}{t^{n(\mathrm{\Lambda })}}\underset{\mu n}{}f^\mu K_{\mu \mathrm{\Lambda }}(q,t).$$
Definition 2: For $`0q<1,0<t<1`$ and $`0x_1,x_2,\mathrm{}`$ such that $`x_i=1`$, define a family $`\{M_n\}`$ of probability measures on partitions of size $`n`$ by
$`M_n(\mathrm{\Lambda })`$ $`=`$ $`{\displaystyle \frac{(1q)^{|\mathrm{\Lambda }|}P_\mathrm{\Lambda }(x;q,t)_{\mu n}f^\mu K_{\mu \mathrm{\Lambda }}(q,t)}{_{s\mathrm{\Lambda }}(1q^{a_\mathrm{\Lambda }(s)+1}t^{l_\mathrm{\Lambda }(s)})}}`$
$`=`$ $`{\displaystyle \frac{(1q)^{|\mathrm{\Lambda }|}P_\mathrm{\Lambda }(x;q,t)t^{n(\mathrm{\Lambda })}dim(\mathrm{\Lambda })}{_{s\mathrm{\Lambda }}(1q^{a_\mathrm{\Lambda }(s)+1}t^{l_\mathrm{\Lambda }(s)})}}`$
Consider the specialization that $`q=0`$ and $`t=\frac{1}{q}`$, where this second $`q`$ is the size of a finite field. Further, set $`x_i=\frac{1}{q^{i1}}\frac{1}{q^i}`$. Then Theorem 14 implies that $`M_n(\mathrm{\Lambda })`$ is the probability that a uniformly chosen element of $`T(n,q)`$ has Jordan type $`\mathrm{\Lambda }`$. The multiplicities have a simple description; letting $`i`$ be the column to which one adds in order to go from $`\lambda `$ to $`\mathrm{\Lambda }`$, it follows that $`\kappa (\lambda ,\mathrm{\Lambda })=q^{\lambda _i^{}}+q^{\lambda _i^{}1}+\mathrm{}+q^{\lambda _{i+1}^{}}`$. Second, $`dim(\mathrm{\Lambda })`$ reduces to a Green’s polynomial $`Q^\mathrm{\Lambda }(q)=Q_{(1^n)}^\mathrm{\Lambda }(q)`$ as in Section 3.7 of \[Mac\]. These polynomials are important in the representation theory of the finite general linear groups. This specialization was the motivation for Definition 2.
The connection with potential theory is given by the following result.
###### Theorem 15
(\[F8\]) The measures of Definition 2 are harmonic with respect to the branching of Definition 1.
It is elementary and well-known that if one starts at the empty partition and transitions from $`\lambda `$ to $`\mathrm{\Lambda }`$ with probability $`\frac{\kappa (\lambda ,\mathrm{\Lambda })M_n(\mathrm{\Lambda })dim(\lambda )}{M_{n1}(\lambda )dim(\mathrm{\Lambda })}`$, one gets samples from any coherent family of measures $`\{M_n\}`$. Applying this principle to the above specialization in which $`M_n(\mathrm{\Lambda })`$ is $`T(n,q)`$ and using Macdonald’s principal specialization formula (page 337 of \[Mac\]) gives the advertised proof of Theorem 12 by means of symmetric functions and potential theory.
Remarks:
1. As indicated in \[F8\], the example of Schur functions ($`q=t<1`$) is also interesting. The measure $`M_n(\mathrm{\Lambda })`$ reduces to $`s_\mathrm{\Lambda }f^\mathrm{\Lambda }`$, where $`s_\mathrm{\Lambda }`$ is a Schur function. Setting $`x_1=\mathrm{}=x_n=\frac{1}{n}`$ and letting $`n\mathrm{}`$, one obtains Plancherel measure, which is important in representation theory and random matrix theory. Letting $`x_1=\mathrm{}=x_n`$ satisfy $`x_i=1`$ (all other $`x_j=0`$) gives a natural deformation of Plancherel measure, studied for instance by \[ItTWi\]. Stanley \[Sta\] shows that this measure on partitions also arises by applying the Robinson-Schensted-Knuth algorithm to a random permutation distributed after a biased riffle shuffle (in other words, this measure encodes information about the longest increasing subsequences of permutations distributed as shuffles).
2. It has been pointed out to the author that the branchings $`\kappa (\lambda ,\mathrm{\Lambda })`$ of Definition 1 are related to the branchings $`\tau (\lambda ,\mathrm{\Lambda })`$ of \[Ke2\] by the formula
$$\kappa (\lambda ,\mathrm{\Lambda })=f(\lambda )\tau (\lambda ,\mathrm{\Lambda })f(\mathrm{\Lambda })^1,$$
for a certain positive function $`f(\lambda )`$ on the set of vertices, which implies by \[Ke3\] that the boundaries of these two branchings are homeomorphic and that the branchings of Definition 1 are multiplicative. Kerov \[Ke2\] has a conjectural description of the boundary. It has been verified for Schur functions \[T\], Kingman branching \[Kin\], and Jack polynomials \[KeOOl\], but remains open for the general case of Macdonald polynomials. In particular, it is open for Hall-Littlewood polynomials, the case related to $`T(n,q)`$. It is interesting that the $`\kappa (\lambda ,\mathrm{\Lambda })`$ of Definition 1 are integers for Hall-Littlewood polynomials, whereas the $`\tau (\lambda ,\mathrm{\Lambda })`$ of \[Ke2\] are not.
## Acknowledgements
The author’s greatest thanks go to Persi Diaconis (his former thesis advisor) for years of friendship, encouragement, and inspiration. He was very helpful in the preparation of this article. We thank Peter M. Neumann and Cheryl E. Praeger for countless conversations about conjugacy classes and computational group theory, and Mark Huber for permission to survey some joint unpublished results. The author received the financial support of an NSF Postdoctoral Fellowship.
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# ROSAT HRI monitoring of X-ray variability in the NLS1 galaxy PKS 0558–504
## 1 Introduction
Narrow-line Seyfert 1 galaxies are identified by their optical emission line properties: the ratio \[O III\]/H$`\beta `$ is less than 3 and FWHM H$`\beta `$ is less than 2000$`\mathrm{km}\mathrm{s}^1`$ (Osterbrock & Pogge 1985, Goodrich 1989). Their optical spectra are also characterized by the presence of strong permitted Fe II, Ca II, O I $`\lambda `$ 8446 lines (Persson 1988). NLS1 exhibit characteristic features at other wavelengths as well: they are seldom radio loud (Ulvestad et al. 1995, Siebert et al. 1999, Grupe et al. 1999, 2000) and they are usually strong infrared emitters (Moran et al. 1996). In X-rays NLS1 have been generally found to have extreme spectral and variability properties that might be related to an extreme value of a fundamental physical parameter, originating from the vicinity of a supermassive black hole (e.g. Brandt & Boller 1998).
PKS 0558–504 ($`z=0.137,m_\mathrm{B}=14.97`$) is one of the few radio-loud NLS1 galaxies ($`R_\mathrm{L}=f_{5\mathrm{G}\mathrm{H}\mathrm{z}}/f_\mathrm{B}27`$, Siebert et al. 1999). It was optically identified on the basis of X-ray positions from the High Energy Astronomy Observatory (HEAO-1, Remillard et al. 1986). A Ginga observation (Remillard et al. 1991) showed an increase of the X-ray flux by 67% in 3 minutes, implying that the apparent luminosity must be enhanced by relativistic beaming. Further X-ray observations with different satellites have confirmed the steep X-ray spectrum and high luminosity of this source, but no more relativistic flares have been presented in the literature. It is important to search for such flares with an X-ray imaging detector to definitively rule out the possibility that the Ginga data suffered from source confusion. Tab. 1 summarizes the luminosities, observed by previous X-ray instruments, converted to the ROSAT soft X-ray band. The conversion to luminosities in the 0.2–2.4 keV energy band was performed using PIMMS, assuming Galactic absorption ($`N_\mathrm{H}`$ = $`4.39\times 10^{20}\mathrm{cm}^2`$, Dickey & Lockman 1990) and a power law spectral model with photon index $`\mathrm{\Gamma }`$ ranging between 2.1 (Remillard et al. 1991) and 3.1 (Brinkmann et al. 1997). However, possible deviations from a single power law or long term spectral changes can lead to systematic uncertainties. The measured soft X-ray spectrum is rather steep and the medium energy power laws are considerably flatter, but the sparse data do not allow determination of whether the source shows spectral steepening towards lower energies or whether long term spectral changes occur during intensity variations. The luminosities were calculated by assuming a Friedman cosmology with $`H_0=70\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$, $`q_0=0.5`$ and isotropic emission.
In this paper we report the results of two ROSAT HRI observation campaigns taken five months apart (in November 1997 and April 1998) and the survey PSPC data (September 1990), with the purpose to check whether the strong X-ray variability is persistent and whether a nearby source contributes to the X-ray flux. In section 2 we present the observations and the spatial analysis. Section 3 deals with the variability of PKS 0558–504. Section 4 contains the main conclusions.
## 2 Observations and spatial analysis
PKS 0558–504 was observed with the ROSAT HRI on November 18 1997, with an effective exposure of 2.14 ksec, and eleven times between April 19–25, 1998, with exposures ranging between 860 sec to 4.4 ksec. All individual observations of April 1998 have been merged with a final total exposure of 21.52 ksec (see Table 2). The data analysis was performed using standard routines within the EXSAS environment (Zimmermann et al. 1994). The count rates (vignetting and dead time corrected) quoted in Table 2, as well as the light curves, were obtained by extracting photons from a circle with 150″ radius around the source center and subtracting a background from a source-free region. In order to reduce the uncertainties from an extrapolation of a steep power law spectrum to low energies (0.1–0.2 keV) where the Galactic absorption is important, we base our discussion on the luminosities in the 0.2–2.4 keV band. As a result the fluxes and luminosities quoted in Table 2 represent lower limits only. For completeness, we also mention in the text the values obtained for the 0.1–2.4 keV energy band, which are typically a factor of two higher. The conversion factor between HRI count rates and luminosities, $`A_1=3.6\times 10^{45}\mathrm{erg}\mathrm{count}^1`$ ($`A_2=7.4\times 10^{45}\mathrm{erg}\mathrm{count}^1`$ for 0.1–2.4 keV), was calculated by assuming a power law spectral model with the best fit parameters ($`\mathrm{\Gamma }=2.99`$, $`N_\mathrm{H}`$ = $`4.53\times 10^{20}\mathrm{cm}^2`$) of the PSPC spectrum.
A peculiar property displayed by PKS 0558–504 is the unusually high ratio of X-ray to radio luminosity (Brinkmann et al. 1997), that might imply a contribution to the X-ray flux from a nearby source. To perform a spatial analysis we used all the April 1998 observations. A contour plot of the X-ray emission overlaid onto the optical image is shown in Fig. 1. The photons were binned in 2″$`\times `$2″ pixels and smoothed with a Gaussian with $`\sigma =6`$″. The surface brightness profile is well fitted by the original HRI PSF-model convolved with an additional Gaussian to allow for the known smearing of the PSF by residual wobble motion, which can vary between different observations (Morse 1994). The equatorial coordinates of the centroid in the HRI image, computed from a Gaussian fit to the spatial distribution, are RA(2000)=$`5^h59^m47\stackrel{s}{.}6`$, DEC(2000)=$`50^\mathrm{o}26\mathrm{}48\mathrm{}`$, in good agreement with the optical position, taking into account that the internal HRI position error is of the order of 5″. The only other X-ray source visible in the field of view is in the south-west of PKS 0558–504 at the position RA(2000)=$`5^h59^m21\stackrel{s}{.}2`$, DEC(2000)=$`50^\mathrm{o}28\mathrm{}23\mathrm{}`$, with a mean count rate of 0.0021 $`\mathrm{counts}\mathrm{s}^1`$, which corresponds to a flux of $`1.4\times 10^{13}\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$ (assuming Galactic absorption and a power law spectral model with $`\mathrm{\Gamma }=2`$). This object is classified as ‘Stellar’ with a magnitude of $`m_\mathrm{B}=19.01`$ in the digitized COSMOS UKST southern sky survey. According to Maccacaro et al. (1988) its X-ray to optical flux ratio suggests that it is an AGN. This only additional X-ray source in the HRI field of view is nearly a factor 1000 fainter than PKS 0558–504, and no other strong X-ray source is found in the ROSAT survey within a radius of 2<sup>o</sup>. Therefore contributions from a previously unknown nearby source to the high luminosity and the strong variability can be ruled out.
## 3 X-ray variability
On the basis of the luminosities quoted in Tab. 1, no long-term X-ray variability by more than a factor of 4 has been seen from PKS 0558–504. Recent SAX and RXTE observations seem to confirm this picture (A. Comastri and K. Leighly, private communication). However, spectral variations and the extrapolation of steep power law spectrum to low energies can lead to uncertain luminosities, as pointed out by Brandt et al. (1999).
More reliable results can be obtained by comparing data from the same instrument. For instance, by comparing the mean count rates of the two HRI observations taken five months apart, we find an increase of the count rates by a factor 2.4, corresponding to a luminosity variation of $`\mathrm{\Delta }L_{0.22.4\mathrm{keV}}3.4\times 10^{45}\mathrm{erg}\mathrm{s}^1`$ ($`\mathrm{\Delta }L_{0.12.4\mathrm{keV}}7\times 10^{45}\mathrm{erg}\mathrm{s}^1`$).
In Fig. 2 we show the total light curve for PKS 0558–504 during April 1998. The data points are binned into bins of 400 s, in order to avoid spurious count rate variations due to the ROSAT wobble. To characterize quantitatively the variability in the light curve, we calculated the excess variance (Nandra et al. 1997), $`\sigma _{\mathrm{rms}}^2=(5.43\pm 2.60)10^2`$. At first sight the most extreme count rate variation seems to occur during April 23 and 24, with $`\mathrm{\Delta }\mathrm{cts}/\mathrm{\Delta }t=(1.34\pm 0.09)`$ $`\times 10^5\mathrm{counts}\mathrm{s}^2`$, corresponding to $`\mathrm{\Delta }L_{0.22.4\mathrm{keV}}/\mathrm{\Delta }t=(4.8\pm 0.3)\times 10^{40}\mathrm{erg}\mathrm{s}^2`$ ($`\mathrm{\Delta }L_{0.12.4\mathrm{keV}}/\mathrm{\Delta }t=(9.9\pm 0.7)\times 10^{40}\mathrm{erg}\mathrm{s}^2`$), calculated by performing a linear least square fit to that part of the light curve. However, if we consider the steep increase of the count rate on April 24 only, we obtain an even more extreme value of $`\mathrm{\Delta }\mathrm{cts}/\mathrm{\Delta }t=(4.13\pm 1.34)\times 10^5\mathrm{counts}\mathrm{s}^2`$, leading to $`\mathrm{\Delta }L_{0.22.4\mathrm{keV}}/\mathrm{\Delta }t=(1.7\pm 0.5)\times 10^{41}\mathrm{erg}\mathrm{s}^2`$ ($`\mathrm{\Delta }L_{0.12.4\mathrm{keV}}/\mathrm{\Delta }t=(3.5\pm 1.1)\times 10^{41}\mathrm{erg}\mathrm{s}^2`$). This value can be used to estimate the lower limit of the radiative efficiency: $`\eta >4.8\times 10^{43}\mathrm{\Delta }L/\mathrm{\Delta }t`$ (Fabian 1979). Straightforward application of the limit gives $`\eta >0.08\pm 0.02`$ ($`\eta >0.17\pm 0.05`$ for the 0.1–2.4 keV energy band), which exceeds the theoretical maximum for accretion onto a Schwarzschild black hole, but not onto a maximally rotating Kerr black hole. The extremely high efficiency, $`\eta \stackrel{>}{}2`$, derived from the Ginga 2–10 keV luminosities (Remillard et al. 1991) however, strongly indicates that some approximations used in the calculation of the efficiency limit must be relaxed, allowing uniform radiation release and relativistic effects in the vicinity of the black hole (e.g. Brandt et al. 1999).
PKS 0558–504 was observed in the ROSAT All-Sky Survey between 1990 September 8 (10:58:02 UT) and 1990 September 14 (22:14:44 UT), with an effective exposure of 1036 s and an average count rate of $`5.5\mathrm{counts}\mathrm{s}^1`$. The light curve, shown in Fig. 3 ($`\sigma _{\mathrm{rms}}^2=(6.32\pm 0.59)10^2)`$, exhibits variability by more than a factor of three with a maximum $`\mathrm{\Delta }\mathrm{cts}/\mathrm{\Delta }t=(14.3\pm 2.4)`$ $`\times 10^5\mathrm{counts}\mathrm{s}^2`$. From the PSPC count rates and the spectral parameters given below we obtain $`\mathrm{\Delta }L_{0.22.4\mathrm{keV}}/\mathrm{\Delta }t=(1.8\pm 0.3)\times 10^{41}\mathrm{erg}\mathrm{s}^2`$ and $`\mathrm{\Delta }L_{0.12.4\mathrm{keV}}/\mathrm{\Delta }t=(3.6\pm 0.6)\times 10^{41}\mathrm{erg}\mathrm{s}^2`$, which are very similar to the values of the most extreme HRI event, but on a different time scale (the rest frame interval is $`\mathrm{\Delta }t7`$h for the survey data and $`\mathrm{\Delta }t1.5`$h for the HRI flare). As a consequence the values derived for the radiative efficiency are similar: $`\eta >0.09\pm 0.01`$ for the 0.2–2.4 keV band and $`\eta >0.17\pm 0.03`$ for the 0.1–2.4 keV band.
The spectrum can be fit with a single power law with photon index $`\mathrm{\Gamma }=2.99\pm 0.09`$ and free absorption of $`N_H=(4.53\pm 0.44)\times 10^{20}`$cm<sup>-2</sup>, in excellent agreement with the Galactic value. The residuals of this fit, given in Fig. 4, do not give strong evidence for deviations from a simple power law although there are indications at $`E1`$keV for some spectral changes. With these parameters the resulting unabsorbed fluxes during the survey observations are $`1.2\times 10^{10}\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$ and $`2.5\times 10^{10}\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$, for the 0.2–2.4 keV and 0.1–2.4 keV energy bands, respectively, corresponding to $`L_{0.22.4\mathrm{keV}}=5.9\times 10^{45}\mathrm{erg}\mathrm{s}^1`$ and $`L_{0.12.4\mathrm{keV}}=1.2\times 10^{46}\mathrm{erg}\mathrm{s}^1`$.
For several AGN, light curves with large amplitude flares have been interpreted as indication for non-linear processes (Green 1993, Boller et al. 1997, Leighly & O’Brien 1997). Given that the April 1998 light curve of PKS 0558–504 presents at least two large flares, we searched for non-Gaussianity (and possibly for non-linearity; see Leighly 1999a for a detailed discussion), by adopting the Green (1993) method: a time series is non-Gaussian if the ratio of its standard deviation to its mean is larger than unity. Using the data points in Fig. 2 we find that the unweighted mean count rate is 1.61 $`\mathrm{counts}\mathrm{s}^1`$ and the standard deviation 0.39 $`\mathrm{counts}\mathrm{s}^1`$ ($`\sigma /\overline{x}=0.24`$). As a result we do not find evidence for non-Gaussian variability. However this method assumes that the sample mean and standard deviation used are accurate representations of the true mean and standard deviation, and this might not be true in our case, due to the limited number of observation intervals. The same conclusions were reached by Leighly (1999a) from ASCA data, using a different method based on the skewness of the flux distribution.
## 4 Conclusions
We have presented ROSAT HRI observations of the radio-loud NLS1 galaxy PKS 0558–504. The main results can be summarized as follows:
From the spatial analysis, no other strong X-ray sources have been detected in the neighborhood of PKS 0558–504, therefore external contributions to the high luminosity and to the strong variability from a nearby source are ruled out. By comparing the X-ray observations throughout the last decade, it is evident that the strong X-ray variability of PKS 0558–504 occurs persistently. During the ROSAT HRI observations, PKS 0558–504 shows strong variability, both on medium (months) and short (days, hours) time scales. The most extreme variation implies a radiative efficiency larger than the theoretical maximum for accretion onto a Schwarzschild black hole, and our findings generally support those of Remillard et al. (1991) where a relativistic $`(\eta \stackrel{>}{}2)`$ flare was discovered. As PKS 0558–504 is a radio-loud object, beamed emission from a jet could be the cause for the brightness and variability in X-rays. However, it is worth noting that the radio-quiet NLS1 PHL 1092 has also shown a relativistic flare $`(\eta \stackrel{>}{}0.6)`$, and radio-quiet NLS1 more generally show enhanced X-ray variability. The soft X-ray spectrum is rather steep with a power law photon index of $`\mathrm{\Gamma }3.0`$ and shows no strong indications for spectral breaks. The obtained medium energy power laws are considerably flatter ($`\mathrm{\Gamma }2.2`$) but the sparse data and the limited energy bands of the different instruments do not allow determination of whether the source shows a spectral steepening towards lower energies or whether long term spectral changes occur during intensity variations. An answer to these vital questions can only be given by the current broad band X-ray missions like SAX, XMM-Newton or Chandra.
###### Acknowledgements.
The ROSAT project is supported by the Bundesministerium für Bildung, Wissenschaft, Forschung und Technologie (BMBF) and the Max-Planck-Gesellschaft. MG acknowledges support from the European Commission under contract number ERBFMRX-CT98-0195 (TMR network “Accretion onto black holes, compact stars and protostars”). WNB acknowledges support from NASA LTSA grant NAG5-8107.
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# Introduction
## Introduction
0.1. Let $`k`$ be a field of characteristic zero. To a nonsingular (irreducible) variety $`X`$ and a morphism $`f:X𝔸^1`$, both defined over $`k`$, was associated the invariant motivic (Igusa) zeta function by Denef and Loeser \[DL2\]. By definition it lives in a power series ring in one variable over the ring $`_L`$, where $``$ is the Grothendieck ring of algebraic varieties over $`k`$, $`L`$ is the class of $`𝔸^1`$ in $``$, and $`_L`$ denotes localization. When $`X=𝔸^d`$ this invariant specializes to both the usual $`p`$–adic Igusa zeta function and the topological zeta function associated to a polynomial $`f`$. (In fact in \[DL2\] the authors treat an even more general invariant, involving motives instead of varieties, from which also the whole Hodge spectrum of $`f`$ at any point of $`f^1\{0\}`$ can be deduced.) This notion of motivic zeta function can easily be extended to an effective divisor $`D`$ instead of just a morphism $`f`$.
The authors were inspired by Kontsevich’s idea of motivic integration. In \[Kon\] Kontsevich associated to a nonsingular irreducible variety $`X`$ and an effective divisor $`D`$ on $`X`$ an invariant $`(D)`$, living by definition in an appropriate completion $`\widehat{}`$ of $`_L`$. He used this invariant to show that birationally equivalent (smooth, projective) Calabi–Yau varieties have the same Hodge numbers.
0.2. There are important formulas for these invariants in terms of an embedded resolution (with strict normal crossings) $`h:YX`$ of supp $`D`$. Let $`\mathrm{dim}X=d`$ and denote by $`E_i,iT`$, the irreducible components of $`h^1(\mathrm{supp}D)`$. To the $`E_i`$ are associated natural multiplicities $`N_i`$ and $`\nu _i`$ defined by $`h^{}D=_{iT}N_iE_i`$ and $`\mathrm{div}(h^{}dx)=_{iT}(\nu _i1)E_i`$, where $`dx`$ is a local generator of the sheaf of regular differential $`d`$-forms on $`X`$. Also we partition $`Y`$ into the locally closed strata $`E_I^{}:=(_{iI}E_i)(_\mathrm{}IE_{\mathrm{}})\text{ for }IT.`$
We denote the class of a variety $`V`$ in $``$ by $`[V]`$, and by analogy with the usual $`p`$–adic Igusa zeta function we denote the variable of the power series ring over $`_L`$ formally by $`L^s`$. Then the motivic zeta function $`𝒵(D,s)`$ of $`D`$ is given by the formula
$$𝒵(D,s)=L^d\underset{IT}{}[E_I^{}]\underset{iI}{}\frac{(L1)L^{\nu _i}(L^s)^{N_i}}{1L^{\nu _i}(L^s)^{N_i}}$$
and so it lives already in a localization of the polynomial ring $`_L[L^s]`$. Kontsevich’s invariant for $`D`$ is given by
$$(D)=L^d\underset{IT}{}[E_I^{}]\underset{iI}{}\frac{L1}{L^{\nu _i+N_i}1}$$
and can thus in some sense be derived from $`𝒵(D,s)`$ by ‘substituting $`s=1`$’.
0.3. One can specialize $`𝒵(D,s)`$ and $`(D)`$ to more ‘concrete’ invariants, involving instead of the class $`[V]`$ of a variety $`V`$ in $``$ other additive invariants as the Hodge polynomial $`H(V)`$ or the Euler characteristic $`\chi (V)`$ of $`V`$. With a little work one obtains for instance from $`𝒵(D,s)`$ the topological zeta function
$$z(D,s)=\underset{IT}{}\chi (E_I^{})\underset{iI}{}\frac{1}{\nu _i+sN_i}(s),$$
which was introduced in \[DL1\] for $`X=𝔸^d`$ and $`k=`$, and the invariant
$$e(D)=\underset{IT}{}\chi (E_I^{})\underset{iI}{}\frac{1}{\nu _i+N_i}.$$
0.4. Can the invariants above be generalized to singular (normal) varieties $`X`$ such that analogous formulas in terms of an embedded resolution are valid ? The main problem is whether these formulas are independent of the chosen resolution. Let $`D`$ be an effective Weil divisor on $`X`$ and $`h:YX`$ an embedded resolution of $`X_{\text{sing}}\mathrm{supp}D`$ with irreducible components $`E_i,iT`$, of $`h^1(X_{\text{sing}}\mathrm{supp}D)`$. Can we generalize the multiplicities $`N_i`$ and $`\nu _i`$ ? When $`D`$ is Cartier (or $``$–Cartier) the same expression $`h^{}D=_{iT}N_iE_i`$ makes sense. We think that the most natural generalization of the $`\nu _i`$ are the log discrepancies given by $`K_Y=h^{}K_X+_{iT}(\nu _i1)E_i`$, where $`K_{}`$ is the canonical divisor. To this end we need in general $`X`$ to be Gorenstein (or $``$–Gorenstein). Up to now the following generalizations appeared (with $`k=`$).
(a) In dimension 2 these multiplicities are defined for arbitrary Weil divisors on normal surfaces. In \[V3\] we introduced a topological zeta function and a motivic zeta function for effective divisors on normal surface germs. We could have done this as well globally, associating to an effective Weil divisor $`D`$ on a normal surface $`X`$ for which $`X_{\text{sing}}\mathrm{supp}D`$ the zeta functions $`𝒵(D,s)`$ and $`z(D,s)`$, given by the same formulas as above.
(b) In arbitrary dimension Batyrev \[B2\] considered the case $`D=0`$ and associated ‘Kontsevich–like’ invariants to a log terminal $`X`$ on the level of Hodge polynomials and Euler characteristics. The last one, which he called stringy Euler number, is given by the formula for $`e(D)`$ in (0.3) with all $`N_i=0`$. The invariant on Hodge polynomial level was used in \[B2\] to define stringy Hodge numbers for projective canonical Gorenstein varieties, and to formulate a topological mirror duality test for canonical Calabi–Yau varieties.
(c) Batyrev \[B3\] also extended his construction to Kawamata log terminal pairs $`(X,D)`$, i.e. pairs such that $`K_X+D`$ is $``$–Cartier and all $`a_i>0`$ in the expression $`K_Y=h^{}(K_X+D)+_{iT}(a_i1)E_i`$. On the Euler characteristic level this invariant is given by the formula
$$e\left((X,D)\right)=\underset{IT}{}\chi (E_I^{})\underset{iI}{}\frac{1}{a_i}.$$
In \[B3\] these invariants are used to prove a version of Reid’s McKay correspondence conjecture.
We should mention that Batyrev is naturally restricted to the log terminality conditions above (all $`\nu _i>0`$ and all $`a_i>0`$, respectively) by applying motivic integration techniques to show that the formulas above are independent of the chosen resolution; see \[B2, Theorem 6.28\].
We also want to remark that $`(D)`$ is generalized in \[DL3\] in a different way (see 3.5).
0.5. In this paper we extend the invariants above beyond the log terminal case to the following general situation. Now let $`X`$ be any normal $``$–Gorenstein variety and $`D`$ an effective $``$–Cartier divisor with $`X_{\text{sing}}\mathrm{supp}D`$. We associate first to these data zeta functions $`𝒵(D,s),Z(D,s)`$ and $`z(D,s)`$ on ‘motivic’ level, Hodge polynomial level and Euler characteristic level, respectively, such that the same formulas as in (0.2) and (0.3) are valid. Then we define ‘Kontsevich’ invariants $`(D),E(D)`$ and $`e(D)`$ on the analogous levels by taking the limit for $`s1`$ in the associated zeta functions (admitting the value $`\mathrm{}`$). In particular when all $`\nu _i+N_i0`$ the formulas in (0.2) and (0.3) are again valid.
Furthermore taking the limit for $`s1`$ in the zeta functions we obtain invariants $`\left((X,D)\right),E\left((X,D)\right)`$ and $`e\left((X,D)\right)`$ of the pair $`(X,D)`$ on the same levels, the last one given by the same formula as in (0.4).
In fact we can relax our condition $`X_{\text{sing}}\mathrm{supp}D`$ to $`LCS(X)\mathrm{supp}D`$, where $`LCS(X)`$ is the locus of log canonical singularities of $`X`$. In particular this locus is empty when $`X`$ is log terminal; so we really generalize the invariants of \[B2\].
0.6. In §1 we recall the motivic zeta function of Denef and Loeser and the invariant of Kontsevich on smooth varieties $`X`$, generalizing the first one to effective divisors instead of regular functions. As an introduction to singular varieties we treat the easy case of a canonical $`X`$ in §2; there we also consider an application to minimal models. For $``$–Gorenstein varieties $`X`$ the zeta functions $`Z(D,s)`$ and $`z(D,s)`$ on the level of Hodge polynomials and Euler characteristics, respectively, are constructed in an elementary way in §3. We provide some examples in §4. The ‘motivic’ version requires more work. In §5 we first introduce a motivic zeta function $`𝒵(D,J,s)`$ on a smooth $`X`$, associated to both an effective divisor $`D`$ and an invertible subsheaf $`J`$ of the sheaf of regular differential forms on $`X`$. (This can be compared with associating a $`p`$–adic Igusa zeta function to both a polynomial and a differential form.) Then we use this object to define the motivic zeta function $`𝒵(D,s)`$ for a $``$–Gorenstein $`X`$ in §6. We include an appendix indicating how to extend the original Kontsevich invariant on smooth $`X`$ to $``$–divisors instead of (ordinary) divisors, needing a finite extension of $`\widehat{}`$.
0.7. Remark. After this work was finished we learned about the proofs of Włodarczyk \[Wł\] and of Abramovich et al \[AKMW\] of the weak factorization conjecture for birational maps. Using weak factorization we can give another proof that the zeta functions in this paper are well defined.
## 1. Smooth varieties
1.1. Let $`k`$ be a field of characteristic zero; the varieties and morphisms we will consider are assumed to be defined over $`k`$. (A variety is a reduced separated scheme of finite type over $`k`$, not necessarily irreducible.)
We fix some terminology concerning resolution. A resolution of an irreducible variety $`X`$ is a proper birational morphism $`h:YX`$ from a smooth variety $`Y`$, which is an isomorphism outside the set $`X_{\text{sing}}`$ of singular points of $`X`$. A log resolution or embedded resolution of an irreducible variety $`X`$ is a resolution $`h:YX`$ of $`X`$ for which $`h^1(X_{\text{sing}})`$ is a divisor with strict normal crossings, i.e. with smooth irreducible components intersecting transversely. A log resolution or embedded resolution of a reduced Weil divisor $`D`$ on a normal variety $`X`$ is a proper birational morphism $`h:YX`$ from a smooth $`Y`$, which is an isomorphism outside $`X_{\text{sing}}D`$, and such that $`h^1(X_{\text{sing}}D)`$ is a divisor with strict normal crossings.
We denote by $``$ the Grothendieck ring of (algebraic) varieties over $`k`$. This is the free abelian group generated by the symbols $`[V]`$, where $`[V]`$ is a variety, subject to the relations $`[V]=[V^{}]`$ if $`VV^{}`$ and $`[V]=[VV^{}]+[V^{}]`$ if $`V^{}`$ is closed in $`V`$. Its ring structure is given by $`[V][V^{}]:=[V\times V^{}]`$. We abbreviate $`L:=[𝔸^1]`$ and denote by $`_L=[L^1]`$ the localization of $``$ w.r.t. the multiplicative set $`\{L^n,n\}`$.
1.2. For $`[V]`$ we denote by $`H(V)[u,v]`$ its Hodge polynomial and by $`\chi (V)`$ its Euler characteristic. We briefly explain these notions.
Let first $`k=`$. Then for a variety $`V`$ we denote by $`h^{p,q}(H_c^i(V,))`$ the rank of the $`(p,q)`$–Hodge component of its $`i`$-th cohomology group with compact support and by $`e^{p,q}(V):=_{i0}(1)^ih^{p,q}(H_c^i(V,))`$ its Hodge numbers. The Hodge polynomial of $`V`$ is $`H(V)=H(V;u,v):=_{p,q}e^{p,q}(V)u^pv^q[u,v]`$.
Precisely by the defining relations of $``$ there is a well defined ring morphism $`H:[u,v]`$ determined by $`[V]H(V)`$.
We denote by $`\chi (V)`$ the topological Euler characteristic of $`V`$, i.e. the alternating sum of the ranks of its Betti or de Rham cohomology groups. Clearly $`\chi (V)=H(V;1,1)`$ and we also obtain a ring morphism $`\chi :`$ determined by $`[V]\chi (V)`$.
For arbitrary $`k`$ (of characteristic zero) we choose an embedding of the field of definition of the variety $`V`$ into $``$. Then we can define the same morphisms $`H`$ and $`\chi `$ on $``$ starting from the $`e^{p,q}(V)`$; they are independent of the chosen embedding since for a smooth projective $`V`$ we have that $`e^{p,q}(V)=(1)^{p+q}\mathrm{dim}_kH^q(V,\mathrm{\Omega }_V^p)`$.
1.3. Till the end of this section we let $`X`$ be a smooth irreducible variety of dimension $`d`$ and $`W`$ a subvariety of $`X`$.
In \[DL2\] Denef and Loeser associate to $`WX`$ and a morphism $`f:X𝔸^1`$ an invariant named motivic Igusa zeta function. We recall here briefly its definition but generalize immediately to effective divisors $`D`$ (instead of functions $`f`$). We refer to \[DL2\] for more details and motivation, and for the relation with the usual $`p`$–adic Igusa zeta function.
We denote by $`(X)`$ the scheme of germs of arcs on $`X`$. It is a scheme over $`k`$ whose $`k`$–rational points are the morphisms $`\mathrm{Spec}k[[t]]X`$ (called the germs of arcs on $`X`$). In fact $`(X)`$ is defined as the projective limit $`\underset{}{\mathrm{lim}}_n(X)`$ of the schemes of truncated arcs $`_n(X)`$, whose $`k`$–rational points are the morphisms $`\mathrm{Spec}(k[t]/t^{n+1}k[t])X`$ (see \[DL2\] and \[BLR, p.276\]). There are canonical morphisms $`\pi _n:(X)_n(X)`$, induced by truncation. Remark also that $`_0(X)=X`$.
Now let $`D`$ be an effective divisor on $`X`$. For $`n`$ we define $`Y_{n,D,W}`$ as the subscheme of $`(X)`$ whose $`K`$–rational points, for any field $`Kk`$, are the morphisms $`\phi :\mathrm{Spec}K[[t]]X`$ satisfying the following conditions :
(i) $`\phi `$ sends the closed point of $`\mathrm{Spec}K[[t]]`$ to a point $`P`$ in $`W`$;
(ii) if $`f`$ is a local equation of $`D`$ at $`P`$, then the power series in $`t`$ given by $`f\phi `$ must be exactly of order $`n`$. (This is clearly independent of the choice of $`f`$.) We then denote by $`X_{n,D,W}`$ the image of $`Y_{n,D,W}`$ in $`_n(X)`$, viewed as a reduced subscheme of $`_n(X)`$. The motivic zeta function of $`D`$ (and $`WX`$) is
$$𝒵_W(D,s)=𝒵_W(X,D,s):=\underset{n}{}[X_{n,D,W}]L^{(n+1)dns}_L[[L^s]].$$
Here $`L^s`$ is just a variable and in the power series ring $`_L[[L^s]]`$ we abbreviate $`L^a(L^s)^b`$ by $`L^{asb}`$ for $`a`$ and $`b`$. (When $`D`$ is given by a global function $`f`$ on $`X`$ Denef and Loeser denoted this invariant by $`_W^{}f^s`$ in \[DL2\].)
One can think here mainly about $`W`$ as being $`X`$ itself, the divisor $`\{f=0\}`$, or a point of $`\{f=0\}`$. This $`W`$–formalism enables us to treat these cases together, and the greater generality is also useful.
1.3.1. We briefly compare this with the classical $`p`$–adic situation. Let $`f_p[x]=_p[x_1,\mathrm{},x_d]`$ and denote by $`|z|=p^{\mathrm{ord}_pz}`$ the $`p`$–adic absolute value of $`z_p`$. Igusa’s local zeta function of $`f`$ is
$$Z_p(f,s):=_{_p^d}|f(x)|^s|dx|$$
for $`s`$ with $`\mathrm{}(s)>0`$, where $`|dx|`$ denotes the Haar measure on $`_p^d`$ such that $`_p^d`$ has measure $`1`$. When $`f_p[x]`$ it is not difficult to verify that
$$Z_p(f,s)=\underset{n}{}\mathrm{card}(X_{n,f})p^{(n+1)dns},$$
where $`X_{n,f}`$ is the image in $`(_p/p^{n+1}_p)^d`$ of $`Y_{n,f}=\{x_p^d|\mathrm{ord}_pf(x)=n\}`$. See \[D2\] for an introduction and an overview on Igusa’s local zeta function.
1.4. There is an important formula for $`𝒵_W(D,s)`$ in terms of a log resolution of $`\mathrm{supp}D`$. In particular it implies the rationality result that $`𝒵_W(D,s)`$ belongs in fact already to a certain localization of the polynomial ring $`_L[L^s]`$.
Let $`h:YX`$ be a log resolution of $`\mathrm{supp}D`$. We denote by $`E_i,iT`$, the irreducible components of $`h^1(\mathrm{supp}D)`$ and by $`N_i`$ and $`\nu _i1`$ the multiplicities of $`E_i`$ in $`h^{}D`$ and the divisor of $`h^{}\text{dx}`$, respectively, where dx is a local generator of the sheaf $`\mathrm{\Omega }_X^d`$ of regular differential $`d`$–forms. We partition $`Y`$ into the locally closed strata $`E_I^{}:=(_{iI}E_i)(_\mathrm{}IE_{\mathrm{}})`$ for $`IT`$. (Here $`E_\varphi =Y_\mathrm{}TE_{\mathrm{}}`$.)
###### Theorem
We have the formula
$$𝒵_W(D,s)=L^d\underset{IT}{}[E_I^{}h^1W]\underset{iI}{}\frac{L1}{L^{\nu _i+sN_i}1}$$
(where $`\frac{1}{L^{\nu _i+sN_i}1}:=\frac{L^{\nu _isN_i}}{1L^{\nu _isN_i}})`$. So $`𝒵_W(D,s)`$ belongs already to the localization $`_L[L^s]_{(1L^{nNs})_{n,N\{0\}}}`$ of the polynomial ring $`_L[L^s]`$.
1.4.1. One should compare this formula with the classical formula of Denef \[D1, Theorems 2.4 and 3.1\] for Igusa’s local zeta function $`Z_p(f,s)`$ of $`f[x_1,\mathrm{},x_d]`$ in terms of a resolution $`h:Y𝔸^d`$ of $`\{f=0\}`$. Using the notation above we have for all but finitely many $`p`$ that
$$Z_p(f,s)=p^d\underset{IT}{}\mathrm{\#}(E_I^{})_{𝔽_p}\underset{iI}{}\frac{p1}{p^{\nu _i+sN_i}1},$$
where $`\mathrm{\#}()_{𝔽_p}`$ denotes the number of $`𝔽_p`$–rational points of the reduction$`modp`$. See \[D1, D2\] for more details.
1.5. Here we generalize $`𝒵_W(D,s)`$ to effective $``$–divisors on $`X`$. Now let $`D`$ be an effective $``$–divisor on $`X`$ and say that $`rD`$ is a divisor for $`r\{0\}`$. We define $`𝒵_W(D,s):=𝒵_W(rD,s/r)`$, meaning by this the motivic zeta function of 1.3 for the divisor $`rD`$, where the variable $`L^s`$ is replaced by a variable $`(L^s)^{1/r}`$. This definition is easily checked to be independent of the chosen $`r`$, using Theorem 1.4. Moreover Theorem 1.4 is still valid in this context. The only difference is that the $`N_i,iT`$, are now rational numbers (of the form $`a/r`$ with $`a\{0\})`$, and one should consider $`L^{sN_i}`$ as an abbreviation of $`((L^s)^{1/r})^{rN_i}`$.
1.6. One can specialize the motivic zeta functions $`𝒵_W(D,s)`$ to more ‘concrete’ invariants on the level of Hodge polynomials and on the level of Euler characteristics.
($`i`$) Let $`D`$ be an effective divisor on $`X`$. Since the Hodge polynomial $`H(𝔸^1)=uv`$ the morphism $`H:[u,v]`$ extends naturally to a ring morphism $`H:_L[u,v]_{uv}=[u,v][(uv)^1]`$ (and further to a morphism on power series rings over these rings). We define
$$Z_W(D,s)=Z_W(X,D,s):=H(𝒵_W(D,s))=\underset{n}{}H(X_{n,D,W})(uv)^{(n+1)dns},$$
where now we denote the variable of the power series ring over $`[u,v]_{uv}`$ by $`(uv)^s`$. Using the notation of 1.4 we have the formula
$`Z_W(D,s)`$ $`=(uv)^d{\displaystyle \underset{IT}{}}H(E_I^{}h^1W){\displaystyle \underset{iI}{}}{\displaystyle \frac{uv1}{(uv)^{\nu _i+sN_i}1}}`$
$`[u,v]_{uv}[(uv)^s]_{(1(uv)^{nNs})_{n,N\{0\}}}(u,v)\left((uv)^s\right).`$
($`ii`$) To specialize further to the level of Euler characteristics one takes heuristically the limit of the expression above for $`u,v1`$. We briefly explain the exact argument; see \[DL2, (2.3)\] for the argument starting from $`𝒵_W(D,s)`$. Let $`R`$ denote the subring of $`[u,v]_{uv}[[(uv)^s]]`$ generated by $`[u,v]_{uv}[(uv)^s]`$ and the elements $`\frac{uv1}{1(uv)^{nNs}}`$, where $`n,N\{0\}`$. $`(Z_W(D,s)`$ lives in $`R`$.) By expanding $`(uv)^s`$ and $`\frac{uv1}{1(uv)^{nNs}}`$ formally into series in $`uv1`$, one constructs a canonical algebra morphism
$$R[u,v]_{uv}[s][(n+sN)^1]_{n,N\{0\}}[[uv1]],$$
where $`[[uv1]]`$ denotes completion with respect to the ideal $`(uv1)`$. Composing this morphism with the quotient map given by dividing out $`(uv1)`$ in this last algebra yields a morphism
$$\phi :R\frac{[u,v]_{uv}}{(uv1)}[s][(n+sN)^1]_{n,N\{0\}}.$$
In this last ring the evaluation $`u=v=1`$ is well defined; we put
$`z_W(D,s)=z_W(X,D,s)`$ $`:=\underset{u,v1}{lim}\phi (Z_W(D,s))`$
$`={\displaystyle \underset{IT}{}}\chi (E_I^{}h^1W){\displaystyle \underset{iI}{}}{\displaystyle \frac{1}{\nu _i+sN_i}}(s).`$
When $`X=𝔸^n`$ and $`D`$ is given by a polynomial $`f`$ these invariants are just the topological zeta functions $`Z_{\text{top}}(f,s)`$ and $`Z_{\text{top},0}(f,s)`$ of \[DL1\] if we take $`W=X`$ and $`W=\{0\}`$, respectively.
($`iii`$) As in 1.5 we can consider $`Z_W(D,s)`$ and $`z_W(D,s)`$ also for $``$–divisors $`D`$.
1.7. Now we recall the original motivic integral, introduced by Kontsevich in \[Kon\], using the notation of 1.3. We refer to \[DL3\] for a detailed exposition in a much more general setting; see also the appendix. A nice introduction is \[C\].
We say that dim $`Mn`$ for $`M`$ if $`M`$ can be expressed as a $``$–linear combination of classes of algebraic varieties of dimension at most $`n`$. We consider the decreasing filtration $`(F^m)_m`$ on $`_L`$, where $`F^m`$ is the subgroup of $`_L`$ generated by $`\{[V]L^i|\mathrm{dim}Vim\}`$, and we denote by $`\widehat{}`$ the completion of $`_L`$ with respect to this filtration.
Let again $`D`$ be an effective divisor on $`X`$. We set
$$_W(D)=_W(X,D):=\underset{n}{}\frac{[X_{n,D,W}]}{L^{(n+1)d}}L^n\widehat{};$$
this expression converges in $`\widehat{}`$ since $`\mathrm{dim}[X_{n,D,W}](n+1)d`$. This invariant was denoted as $`[_Xe^D]`$ by Kontsevich (for $`W=X`$) and as $`_{\pi _0^1W}L^{\mathrm{ord}_t𝒪(D)}𝑑\mu `$ in \[DL3\]. In this last paper Denef and Loeser develop an integration theory for semi–algebraic subsets of $`(X)`$ with values in $`\widehat{}`$ such that $`[X_{n,D,W}]/L^{(n+1)d}`$ is just the volume of $`Y_{n,D,W}`$. See also §5 and the appendix.
1.8. Remark. As far as we know it is not clear whether or not the natural morphism $`_L\widehat{}`$ is injective; its kernel is $`_mF_m`$. However for an algebraic variety $`V`$ we have that $`H(V)`$ and $`\chi (V)`$ only depend on the image of $`[V]`$ in $`\widehat{}`$, see 1.12.
###### 1.9. Theorem {\rm\[Kon\]\[DL3, (6.5)\]}
Using the notation of 1.4 we have the following formula for $`_W(D)`$ in terms of a log resolution $`h:YX`$ of supp $`D`$ :
$$_W(D)=L^d\underset{IT}{}[E_I^{}h^1W]\underset{iI}{}\frac{L1}{L^{\nu _i+N_i}1}\text{ in }\widehat{}.$$
In particular $`_W(D)`$ belongs to the image of $`_L[(L^n1)^1]_{n\{0\}}`$ in $`\widehat{}`$.
So by Theorem 1.4 we obtain that $`_W(D)=𝒵_W(D)|_{s=1}`$ in $`\widehat{}`$, where the evaluation ‘$`s=1`$’ means substituting $`L^1`$ for the variable $`L^s`$.
1.10. The following important change of variables formula is a special case of \[DL3, Lemma 3.3\], and was also mentioned in \[Kon\].
###### Theorem
Let also $`X^{}`$ be a smooth irreducible variety and $`\rho :X^{}X`$ a proper birational morphism. Let $`D`$ be an effective divisor on $`X`$. Then
$$_W(X,D)=_{\rho ^1W}(X^{},\rho ^{}D+K_{X^{}|X})$$
where $`K_{X^{}|X}=K_X^{}\rho ^{}K_X`$ is the relative canonical divisor or discrepancy divisor.
1.11. It is possible to generalize the set–up in 1.7 – 1.10 to effective $``$–divisors. We treat this in the appendix. In particular we obtain for an effective $``$–divisor $`D`$ on $`X`$, such that $`rD`$ is a divisor for an $`r\{0\}`$, an analogous invariant $`_W(D)\widehat{}[L^{1/r}]`$. It is given in terms of a log resolution $`h:YX`$ (as in 1.4) by the same formula as in 1.9, where now the $`N_i`$ belong to $`\frac{1}{r}(\{0\}`$). So $`_W(D)`$ belongs to the image of $`[L^{1/r}][(L^{n/r}1)^1]_{n\{0\}}`$ in $`\widehat{}[L^{1/r}]`$.
When $`\mathrm{supp}D`$ has strict normal crossings we extend in the appendix the notion of $`_W(D)`$ further to the case that all coefficients of $`D`$ are $`>1`$. Remark that then in Theorem 1.9 (with $`h=Id_X`$) all $`\nu _i=1`$, and our condition on the coefficients of $`D`$ is thus precisely that all $`\nu _i+N_i>0`$.
1.12. One can also specialize the invariant $`_W(D)`$ to the level of Hodge polynomials and Euler characteristics. We only consider expressions in terms of log resolutions (using the notation of 1.4).
The morphism $`H:[u,v]`$ extends canonically to a morphism $`H:_L[(L^n1)^1]_{n\{0\}}[uv]_{uv}[((uv)^n1)^1]_{n\{0\}}(u,v)`$. Since the kernel of the natural map $`_L\widehat{}`$ is killed by $`H`$ we can in fact consider $`H`$ as a morphism from the image of $`_L[(L^n1)^1]_{n\{0\}}`$ in $`\widehat{}`$ into $`(u,v)`$.
We define for an effective divisor $`D`$ on $`X`$ the invariants
$`E_W(D)=E_W(X,D)`$ $`:=H(_W(D))`$
$`=(uv)^d{\displaystyle \underset{IT}{}}H(E_I^{}h^1W){\displaystyle \underset{iI}{}}{\displaystyle \frac{uv1}{(uv)^{\nu _i+N_i}1}}(u,v)`$
and
$$e_W(D)=e_W(X,D):=\underset{u,v1}{lim}E_W(X,D)=\underset{IT}{}\chi (E_I^{}h^1W)\underset{iI}{}\frac{1}{\nu _i+N_i}.$$
The extended notions of $`_W(D)`$ for $``$–divisors of 1.11 can analogously be specialized. We obtain the same expressions where now the $`N_i`$ are rational; then $`E_W(D)`$ is a rational function in $`u,v`$ with ‘fractional powers’. For $`W=X`$ this was already considered by Batyrev \[B3\].
## 2. Immediate generalizations and applications
2.1. We recall some terminology with origins in the Minimal Model Program. See for example \[KM, KMM, Kol\].
On any normal variety $`V`$ there is a well–defined linear equivalence class of canonical Weil divisors, denoted by $`K_V`$. An arbitrary Weil divisor $`D`$ on $`V`$ is called $``$–Cartier if $`rD`$ is Cartier for some $`r\{0\}`$. A normal variety $`V`$ is called ($``$–)Gorenstein if $`K_V`$ is ($``$–)Cartier.
Let $`X`$ be a normal variety and $`D`$ a $``$–divisor on $`X`$ such that $`K_X+D`$ is $``$–Cartier. (In particular we can have $`D=0`$ and then $`X`$ is $``$–Gorenstein.) Let $`\rho :YX`$ be a log resolution of $`\mathrm{supp}D`$ and denote by $`E_i,iT`$, the irreducible components of $`h^1`$ $`(X_{\mathrm{sing}}\mathrm{supp}D)`$. Then we can write
$$K_Y=\rho ^{}(K_X+D)+\underset{iT}{}(a_i1)E_i$$
in $`\mathrm{Pic}Y`$ and $`a_i=a_i(X,D;E_i)`$ is called the log discrepancy (with respect to the pair $`(X,D)`$) of $`E_i`$ for $`iT`$. This number $`a_i`$ does not depend on the chosen resolution (it is determined by the valuation on $`k(X)`$ associated to $`E_i`$). Remark that when $`X`$ is smooth and $`D=0`$ the numbers $`\nu _i`$ defined in 1.4 are just log discrepancies. ($`i`$) Let first $`D=0`$. The variety $`X`$ is called terminal, canonical, log terminal and log canonical if for some (or, equivalently, any) log resolution of $`X`$ we have that $`a_i>1`$, $`a_i1`$, $`a_i>0`$ and $`a_i0`$, respectively, for all $`iT`$.
($`ii`$) When $`D0`$ the pair $`(X,D)`$ is said to be Kawamata log terminal (shortly klt) if for some (or any) log resolution of supp $`D`$ we have that $`a_i>0`$ for all $`iT`$. In particular this implies that, if $`D=_id_iD_i`$ with the $`D_i`$ irreducible, all $`d_i<1`$. (See \[Kol, S\] for a discussion of other log terminality notions for pairs.)
($`iii`$) A closed subvariety $`CX`$ is called a log canonical centre of $`X`$ if for some log resolution $`\rho :YX`$ there exists $`iT`$ such that $`\rho (E_i)=C`$ and $`a_i0`$. The locus of log canonical singularities of $`X`$, denoted by $`LCS(X)`$, is the union of all log canonical centres of $`X`$. In particular $`LCS(X)=\mathrm{}X`$ is log terminal. (Hence a more appropriate notation for this locus, proposed by Kollár, would be Nlt(X), indicating the locus where $`X`$ is not log terminal.)
2.2. A natural idea, inspired by Theorem 1.10, to generalize the invariant $`_W(X,D)`$ to a ($``$–)divisor $`D`$ on a singular variety $`X`$ is as follows. Take a resolution $`h:YX`$ of $`X`$ and define $`_W(X,D)`$ as $`_{h^1W}(Y,h^{}D+K_{Y|X})`$, whenever this makes sense, and verify independency of the chosen resolution. So we want $`X`$ to be $``$–Gorenstein and $`h^{}D+K_{Y|X}`$ to be effective, or at least that its coefficients are $`>1`$ if its support has normal crossings.
Below we treat the ‘instructional’ case that $`X`$ is ($``$–)Gorenstein and canonical and $`D`$ is an effective ($``$–)Cartier divisor.
###### 2.3. Definition – Proposition
(i) Let $`X`$ be a Gorenstein and canonical variety and $`W`$ a subvariety of $`X`$; let $`D`$ be an effective Cartier divisor on $`X`$. Take a resolution $`h:YX`$ of $`X`$. Then we define
$$_W(X,D):=_{h^1W}(Y,h^{}D+K_{Y|X})\widehat{}.$$
(ii) More generally let $`X`$ be $``$–Gorenstein and canonical and $`W`$ a subvariety of $`X`$; let $`D`$ be an effective $``$–Cartier divisor on $`X`$. Say $`rK_X`$ and $`rD`$ are Cartier for an $`r\{0\}`$. Take a resolution $`h:YX`$ of $`X`$. Then we define $`_W(X,D)\widehat{}[L^{1/r}]`$ as above.
###### Demonstration Proof
(i) The divisor $`h^{}D+K_{Y|X}`$ is effective since $`K_{Y|X}`$ is effective, which is equivalent to $`X`$ being canonical. Let now $`h^{}:Y^{}X`$ be another log resolution of $`X`$. Since two such resolutions are always dominated by a third it is sufficient to consider the case that $`h^{}`$ factors through $`h`$ as $`h^{}:Y^{}\stackrel{𝜋}{}Y\stackrel{}{}X`$. Then by Theorem 1.10 we have
$`_{h^1W}(Y,h^{}D+K_{Y|X})`$ $`=_{\pi ^1h^1W}(Y^{},\pi ^{}(h^{}D+K_{Y|X})+K_{Y^{}|Y})`$
$`=_{h^^1W}(Y^{},h^{}D+K_{Y^{}|X}).`$
(ii) Completely analogous, using the extended theory for $``$–divisors mentioned in 1.11. ∎
When $`h:YX`$ is a log resolution of $`\mathrm{supp}D`$ we have the same formula as in Theorem 1.9, where the $`\nu _i`$ must be generalized according to their meaning as log discrepancies. More precisely, denoting the irreducible components of $`h^1(X_{\mathrm{sing}}\mathrm{supp}D)`$ by $`E_i,iT`$, we set $`h^{}D=_{iT}N_iE_i`$ and $`K_Y=h^{}K_X+_{iT}(\nu _i1)E_i`$. Then $`h^{}D+K_{Y|X}=_{iT}(\nu _i+N_i1)E_i`$ and so
$$_W(X,D)=L^d\underset{IT}{}[E_I^{}h^1W]\underset{iI}{}\frac{L1}{L^{\nu _i+N_i}1},$$
where $`d`$ is the dimension of $`X`$.
2.4. With essentially the same arguments, but needing more material from the appendix, we could introduce $`_W(X,D)`$ for a $``$–Gorenstein variety $`X`$ and a $``$–Cartier divisor $`D`$ on $`X`$ such that the pair $`(X,D)`$ is klt. (Check that this is more general than the case in 2.3 !). On the level of Hodge polynomials this would be possible using \[B2, Theorems 6.27 and 6.28\]. We do not pursue this here; our invariants $`E_W(X,D)`$ in §3 and $`_W(X,D)`$ in §6 cover this case anyhow.
2.5. In the rest of this section we present an application on minimal models, taking $`k=`$.
Recall that an irreducible projective variety $`V`$ is called a minimal model if $`V`$ is terminal and $`K_V`$ is numerically effective (shortly nef ), i.e. the intersection number $`K_VC0`$ for any irreducible curve $`C`$ on $`V`$. The Minimal Model Program predicts the existence of a minimal model in every birational equivalence class $`𝒞`$ of nonnegative Kodaira dimension; furthermore one should be able to transform every smooth irreducible projective variety in $`𝒞`$ by a finite number of divisorial contractions and flips to a minimal model.
In dimension 2 it is well known that each such class has a unique minimal model, which is moreover smooth (then divisorial contractions are just blowing–downs and flips do not occur). In dimension 3 the existence and desired property of minimal models were proved by Mori; here it is crucial to allow terminal singularities, and minimal models are not unique in a given birational equivalence class of nonnegative Kodaira dimension. In dimension $`4`$ the Minimal Model Program is still a major conjecture in algebraic geometry and is becoming a working hypothesis.
It is natural and important in this context to look for invariants which are shared by birationally equivalent minimal models. In \[Wa\] Wang proved that birationally equivalent smooth minimal models have the same Betti numbers, using the following result \[Wa, Corollary 1.10\].
###### 2.6. Proposition
Let $`f:VV^{}`$ be a birational map between two minimal models. Then there exist a smooth projective variety $`Y`$ and birational morphisms $`\phi :YV,\phi ^{}:YV^{}`$ such that $`\phi ^{}K_V=\phi ^{}K_V^{}`$.
(In fact Wang only needs $`V`$ and $`V^{}`$ to be terminal varieties for which $`K_V`$ and $`K_V^{}`$ are nef along the exceptional loci of $`f`$ in $`V`$ and $`V^{}`$, respectively, to conclude.) This result has more interesting consequences.
###### 2.7. Theorem
Let $`V`$ and $`V^{}`$ be birationally equivalent minimal models. Then
(i) $`_V(V,0)=_V^{}(V^{},0)`$, and
(ii) if $`V`$ and $`V^{}`$ are smooth, then $`[V]=[V^{}]`$.
###### Demonstration Proof
(i) Take $`V\stackrel{𝜑}{}Y\stackrel{\phi ^{}}{}V^{}`$ as in Proposition 2.6. Then by Theorem 1.10 (and its generalization in 1.11) we have
$$_V(V,0)=_Y(Y,K_Y\phi ^{}K_V)=_Y(Y,K_Y\phi {}_{}{}^{}K_V^{})=_V^{}(V^{},0).$$
(ii) For any smooth variety $`X`$ we have that $`_X(X,0)=[X]`$. ∎
As a corollary birationally equivalent smooth minimal models have the same Hodge numbers and a fortiori the same Betti numbers. In particular this is true for smooth Calabi–Yau varieties. See also \[B1, Theorems 1.1 and 4.2\].
2.8. Assuming the Minimal Model Program in some dimension $`d`$ we can use Theorem 2.7 to define a birational invariant. For any birational equivalence class $`𝒞`$ of nonnegative Kodaira dimension the expression $`:=_X(X,0)`$ is independent of a chosen minimal model $`X`$. Looking at 2.3 it is given by the following formula in terms of any log resolution $`h:YX`$ of any minimal model $`X`$ of $`𝒞`$. Denote by $`E_i,iT`$, the irreducible components of $`h^1(X_{\text{sing}})`$ and set $`K_Y=h^{}K_X+_{iT}(\nu _i1)E_i`$. Then
$$=L^d\underset{IT}{}[E_I^{}]\underset{iI}{}\frac{L1}{L^{\nu _i}1}.$$
One could extract ‘minimal stringy Hodge numbers’ from (the Hodge polynomial version of) this invariant, see \[B2\]; and maybe it is related to a ‘minimal cohomology theory’ as explained in \[Wa\].
## 3. Singular varieties; on the level of Hodge polynomials and Euler characteristics
3.1. Our aim in this paper is to associate zeta functions and ‘Kontsevich’ invariants to effective $``$–Cartier divisors $`D`$ on arbitrary $``$–Gorenstein varieties $`X`$ for which $`X_{\mathrm{sing}}\mathrm{supp}D`$, generalizing the notions in §1. In this section we realize this on the level of Hodge polynomials and Euler characteristics in a fairly elementary way. The more general case on the level of the Grothendieck ring will be treated in §5.
3.2. We fix notation for this section. Let $`X`$ be a $``$–Gorenstein variety and $`D`$ an effective $``$–Cartier divisor on $`X`$. (When $`\mathrm{dim}X=2`$ we only need that $`X`$ is normal and $`D`$ can be any effective Weil divisor with rational coefficients, see \[V3\].) For a log resolution $`h:YX`$ of $`\mathrm{supp}D`$ we denote by $`E_i,iT`$, the irreducible components of $`h^1(X_{\mathrm{sing}}\mathrm{supp}D)`$ and we put $`E_I^{}:=(_{iI}E_i)(_\mathrm{}IE_{\mathrm{}})`$ for $`IT`$. We also set $`h^{}D=_{iT}N_iE_i`$ and $`K_Y=h^{}K_X+_{iT}(\nu _i1)E_i`$. Remember that now the $`\nu _i`$ and they can be negative or zero.
In the sequel we will again consider arbitrary subvarieties $`W`$ of $`X`$. One can think mainly about $`W`$ being for example $`X`$, $`\mathrm{supp}D`$, $`X_{\mathrm{sing}}`$ or a point of $`X_{\mathrm{sing}}`$.
###### \bf3.3. Definition – Proposition
Let $`X`$ be a $``$–Gorenstein variety of dimension $`d`$ and $`W`$ a subvariety of $`X`$. Let $`D`$ be an effective $``$–Cartier divisor on $`X`$ such that $`X_{\mathrm{sing}}\mathrm{supp}D`$. Take $`r\{0\}`$ with $`rK_X`$ and $`rD`$ Cartier.
(i) The zeta function $`Z_W(D,s)=Z_W(X,D,s)`$ is the unique rational function in the variable $`(uv)^{s/r}`$ and with coefficients in (the fraction field of) $`[u,v][(uv)^{1/r}]`$ such that for $`n>>0`$
$$Z_W(D,n)=E_{h^1W}(Y,nh^{}D+K_{Y|X}),$$
where $`h:YX`$ is a resolution of $`X`$.
(ii) Let $`h:YX`$ be a log resolution of $`\mathrm{supp}D`$. With the notation of 3.2 we have that
$$Z_W(D,s)=\frac{1}{(uv)^d}\underset{IT}{}H(E_I^{}h^1W)\underset{iI}{}\frac{uv1}{(uv)^{\nu _i+sN_i}1}.$$
###### Demonstration Proof
Let $`h_i:Y_iX`$ be resolutions of $`X`$ for $`i=1,2`$. We first show that the defining expressions for $`Z_W(D,n)`$ using $`Y_1`$ and $`Y_2`$ are equal when $`n>>0`$. Take $`n`$ such that $`nh_i^{}D+K_{Y_i|X}`$ is effective for $`i=1,2`$ (here we need that $`X_{\mathrm{sing}}\mathrm{supp}D)`$, and take a resolution $`h:YX`$ of $`X`$ dominating both $`Y_1`$ and $`Y_2`$.
$``$$``$$``$$``$$``$$`\phi _1`$$`\phi _2`$$`h`$$`h_2`$$`h_1`$$`Y_1`$$`Y`$$`Y_2`$$`X`$
Then by Theorem 1.10 (for $``$–divisors and on the level of Hodge polynomials) we have for $`i=1,2`$ that
$`E_{h_i^1W}(Y_i,nh_i^{}D+K_{Y_i|X})`$ $`=E_{\phi _i^1h_i^1W}(Y,\phi _i^{}(nh_i^{}D+K_{Y_i|X})+K_{Y|Y_i})`$
$`=E_{h^1W}(Y,nh^{}D+K_{Y|X}).`$
Choosing now $`h:YX`$ as a log resolution for $`\mathrm{supp}D`$ we have that
$$E_{h^1W}(Y,nh^{}D+K_{Y|X})=\frac{1}{(uv)^d}\underset{IT}{}H(E_I^{}h^1W)\underset{iI}{}\frac{uv1}{(uv)^{\nu _i+nN_i}1}.$$
Hence for $`n>>0`$ the stated rational function in (ii) indeed yields $`E_{h^1W}(Y,nh^{}D+K_{Y|X})`$ when evaluating in $`s=n`$ (i.e. in $`(uv)^{s/r}=(uv)^{n/r})`$.
Finally this rational function must be unique since a polynomial over the domain $`[u,v][(uv)^{1/r}]`$ can have at most finitely many zeroes. ∎
###### 3.4. Definition
With the same notation as in 3.3 we define the topological zeta function of $`D`$ as
$$z_W(D,s)=z_W(X,D,s):=\underset{IT}{}\chi (E_I^{}h^1W)\underset{iI}{}\frac{1}{\nu _i+sN_i}(s).$$
We can justify this definition either by an analogous proof or by obtaining $`z_W(D,s)`$ from $`Z_W(D,s)`$ by a limit argument as in 1.6.
3.5. In the following we extend Kontsevich’s construction $`E_W(X,D)`$ to $``$–Gorenstein varieties $`X`$. We should remark here that in \[DL3\] Denef and Loeser also generalized in a different way $`_W(X,D)`$ to (arbitrary) singular varieties $`X`$. We consider their point of view as more ‘integrational’ and ours as more ‘geometrical’. Our idea is simply to substitute $`s=1`$ in $`Z_W(D,s)`$ when this makes sense or, more generally, to take the limit for $`s1`$.
###### 3.6. Definition
Let $`X`$ be a $``$–Gorenstein variety and $`W`$ a subvariety of $`X`$. Let $`D`$ be an effective $``$–Cartier divisor on $`X`$ such that $`X_{\mathrm{sing}}\mathrm{supp}D`$. Take $`r\{0\}`$ with $`rK_X`$ and $`rD`$ Cartier. Then we put
$$E_W(D)=E_W(X,D):=\underset{s1}{lim}Z_W(X,D,s)(u^{1/r},v^{1/r})\{\mathrm{}\}.$$
Remarks. (1) By $`lim_{s1}`$ we mean taking the limit $`(uv)^{s/r}(uv)^{1/r}`$. This is well defined since $`Z_W(D,s)`$ is a rational function in the variable $`(uv)^{s/r}`$ over a field.
(2) If there exists a log resolution $`h:YX`$ of $`\mathrm{supp}D`$ for which $`\nu _i+N_i0`$ for all $`iT`$, then, because of the formula in 3.3(ii), we obtain $`E_W(D)`$ from $`Z_W(D,s)`$ simply by substituting $`(uv)^{1/r}`$ for $`(uv)^{s/r}`$. (We formulate this below as Proposition 3.7.) If on the other hand there does not exist such a log resolution, then in general we will have $`E_W(D)=\mathrm{}`$. However there are cases where our definition then yields an element in $`(u^{1/r},v^{1/r})`$, see example 4.1.
###### 3.7. Proposition
Let $`WX`$ and $`D`$ be as in 3.6. Let $`h:YX`$ be a log resolution of $`\mathrm{supp}D`$ for which $`\nu _i+N_i0`$ for all $`iT`$ (using the notation of 3.2). Then
$$E_W(D)=\frac{1}{(uv)^d}\underset{IT}{}H(E_I^{}h^1W)\underset{iI}{}\frac{uv1}{(uv)^{\nu _i+N_i}1}.$$
So indeed we extended Kontsevich’s invariant for smooth $`X`$ on the level of Hodge polynomials (1.12).
###### 3.8. Definition – Proposition
Let $`WX`$ and $`D`$ be as in 3.6. We define
$$e_W(D)=e_W(X,D):=\underset{s1}{lim}z_W(X,D,s)\{\mathrm{}\}.$$
Let $`h:YX`$ be a log resolution of $`\mathrm{supp}D`$ for which $`\nu _i+N_i0`$ for all $`iT`$. Then
$$e_W(D)=\underset{IT}{}\chi (E_I^{}h^1W)\underset{iI}{}\frac{1}{\nu _i+N_i}.$$
3.9. Next we introduce analogous invariants for pairs $`(X,D)`$, which will coincide with Batyrev’s stringy $`E`$–function and stringy Euler number for klt pairs \[B3\].
###### 3.10. Definition – Proposition
Let $`X`$ be a $``$–Gorenstein variety and $`W`$ a subvariety of $`X`$. Let $`D`$ be an effective $``$–Cartier divisor on $`X`$ such that $`X_{\mathrm{sing}}\mathrm{supp}D`$. Take $`r\{0\}`$ with $`rK_X`$ and $`rD`$ Cartier.
(i) We put
$$E_W\left((X,D)\right):=\underset{s1}{lim}Z_W(X,D,s)(u^{1/r},v^{1/r})\{\mathrm{}\}.$$
(ii) Let $`h:YX`$ be a log resolution of $`\mathrm{supp}D`$. Using the notation of 3.2, let $`a_i,iT`$, denote the log discrepancy of $`E_i`$ with respect to the pair $`(X,D)`$. Then, if $`a_i0`$ for all $`iT`$, we have
$$E_W\left((X,D)\right)=\frac{1}{(uv)^d}\underset{IT}{}H(E_I^{}h^1W)\underset{iI}{}\frac{uv1}{(uv)^{a_i}1}.$$
Remark. By $`lim_{s1}`$ we mean taking the limit $`(uv)^{s/r}(uv)^{1/r}`$.
###### Demonstration Proof
If $`\nu _iN_i0`$ for all $`iT`$, then, because of the formula for $`Z_W(X,D,s)`$ in 3.3(ii) this limit procedure just means substituting $`(uv)^{1/r}`$ for the variable $`(uv)^{s/r}`$. Clearly we obtain the stated formula for $`E_W\left((X,D)\right)`$ since $`a_i=\nu _iN_i`$ for $`iT`$. ∎
3.11. When the pair $`(X,D)`$ is klt and for $`W=X`$ Batyrev introduced in \[B3\] the same invariant as the stringy $`E`$–function of $`(X,D)`$, denoted by $`E_{st}(X,D)`$. (We do not recover his invariant completely as a special case of $`E_W\left((X,D)\right)`$ because Batyrev only requires $`K_X+D`$ to be $``$–Cartier.) Analogously the invariant $`e_X\left((X,D)\right)`$ below was baptized stringy Euler number by Batyrev and denoted by $`e_{st}(X,D)`$.
###### 3.12. Definition – Proposition
Let $`WX`$ and $`D`$ be as in 3.10. We define
$$e_W\left((X,D)\right):=\underset{s1}{lim}z_W(X,D,s)\{\mathrm{}\}.$$
Let $`h:YX`$ be a log resolution of $`\mathrm{supp}D`$ for which $`\nu _iN_i0`$ for all $`iT`$. Then, denoting by $`a_i`$ the log discrepancy of $`E_i`$ with respect to the pair $`(X,D)`$, we have
$$e_W\left((X,D)\right)=\underset{IT}{}\chi (E_I^{}h^1W)\underset{iI}{}\frac{1}{a_i}.$$
3.13. In Definition–Proposition 3.3, and hence in all subsequent constructions, we required the effective divisor $`D`$ to satisfy $`X_{\mathrm{sing}}\mathrm{supp}D`$. We needed this to assure that for a resolution $`h:YX`$ the divisor $`nh^{}D+K_{Y|X}`$ would be effective for $`n>>0`$. However, using $`A5`$ in the appendix, it is in fact sufficient to require that $`\mathrm{supp}D`$ contains the locus of log canonical singularities $`LCS(X)`$ of $`X`$.
###### 3.14. Definition – Theorem
Let $`X`$ be a $``$–Gorenstein variety of dimension $`d`$ and $`W`$ a subvariety of $`X`$. Let $`D`$ be an effective $``$–Cartier divisor on $`X`$ such that $`LCS(X)\mathrm{supp}D`$. Take $`r\{0\}`$ with $`rK_X`$ and $`rD`$ Cartier.
(i) The zeta function $`Z_W(D,s)=Z_W(X,D,s)`$ is the unique rational function in the variable $`(uv)^{s/r}`$ and with coefficients in (the fraction field of) $`[u,v][(uv)^{1/r}]`$ such that for $`n>>0`$
$$Z_W(D,n)=E_{h^1W}(Y,nh^{}D+K_{Y|X}),$$
where $`h:YX`$ is a log resolution for $`\mathrm{supp}D`$.
(ii) With the notation of 3.2 for $`h`$ we have that
$$Z_W(D,s)=\frac{1}{(uv)^d}\underset{IT}{}H(E_I^{}h^1W)\underset{iT}{}\frac{uv1}{(uv)^{\nu _i+sN_i}1}.$$
###### Demonstration Proof
We proceed analogously as in the proof of 3.3, but now working only with log resolutions $`h:YX`$ of $`\mathrm{supp}D`$. Then for $`n>>0`$ the coefficients $`d_i=nN_i+\nu _i1`$ of $`nh^{}D+K_{Y|X}`$ all satisfy $`d_i>1`$. Indeed any exceptional component $`E_i`$ of $`h`$ for which $`\nu _i0`$ satisfies $`h(E_i)LCS(X)\mathrm{supp}D`$, and hence $`N_i>0`$ for such an $`E_i`$. So in this case the invariant $`E_{h^1W}(Y,nh^{}D+K_{Y|X})`$ is well defined by $`A5`$ and Theorem A6. ∎
Remark. In the formula above the ‘denominators’ $`\nu _i+sN_i`$ are thus always nonzero since either $`\nu _i>0`$ or $`N_i>0`$.
3.15. We can also extend all invariants which we considered in 3.4 – 3.12, i.e. $`z_W(D,s)`$, $`E_W(D)`$, $`e_W(D)`$, $`E_W\left((X,D)\right)`$ and $`e_W\left((X,D)\right)`$, to the case that only $`LCS(X)\mathrm{supp}D`$.
In particular when $`X`$ is log terminal and $`D=0`$, then our invariants $`E_X(0)`$ and $`e_X(0)`$ are precisely the stringy E–function $`E_{st}(X;u,v)`$ and stringy Euler number $`e_{st}(X)`$ of Batyrev \[B2\].
## 4. Examples
In this section we present a number of examples, first in dimension two and then in higher dimension, for which we compute the invariants introduced above. Recall (see (0.4(a)) that in dimension two we can consider more generally Weil divisors instead of Cartier divisors.
4.1. Let $`0X`$ be a normal surface germ with minimal resolution $`h:YX`$ such that $`h^1\{0\}=E_0E_g`$, where $`E_0`$ and $`E_g`$ are nonsingular curves of genus $`0`$ and $`g2`$, respectively, intersecting transversely. So $`h`$ is already a log resolution of $`X`$. (This singularity is quasihomogeneous.) Let $`E`$ be a nonsingular curve (germ) in $`Y`$ intersecting $`E_g`$ transversely in one point and disjoint from $`E_0`$. Denote $`D=h(E)`$; so $`D`$ is a prime Weil divisor on $`X`$ through $`0`$ and $`h`$ is also a log resolution of $`D`$. See Figure 1.
....................................................................................................................................................................................................................................................................................................................................$``$$`D`$$`0`$$`X`$$``$$`h`$ $`E_0`$$`E_g`$$`E`$$`Y`$Figure 1
Let $`\kappa _0`$ and $`\kappa _g`$ denote the self–intersection number of $`E_0`$ and $`E_g`$ on $`Y`$, respectively; we have that $`\kappa _02`$ and $`\kappa _g1`$. We will treat the germs $`0X`$ for which $`N:=2g\kappa _g1>0`$ in order to compute $`z_0(ND,s)`$ and $`e_0(ND)`$ for the effective Weil divisor $`ND`$ on $`X`$.
We denote as usual $`h^{}ND=NE+N_0E_0+N_gE_g`$ and $`K_Y=h^{}K_X+(\nu _01)E_0+`$ $`(\nu _g1)E_g`$. The following relations are well known (see for example \[V3, Lemma 2.3\]) :
$$\{\begin{array}{cc}\hfill \kappa _0N_0& =N_g\hfill \\ \hfill \kappa _0\nu _0& =\nu _g+1\hfill \end{array}\text{ and }\{\begin{array}{cc}\hfill \kappa _gN_g& =N_0+N\hfill \\ \hfill \kappa _g\nu _g& =(\nu _01)+22g.\hfill \end{array}$$
A short computation yields the expression for $`N_0`$, $`\nu _0`$, $`N_g`$ and $`\nu _g`$ in terms of our data $`\kappa _0`$, $`\kappa _g`$ and $`g`$ :
$$\{\begin{array}{cc}\hfill N_0& =\frac{N}{\kappa _0\kappa _g1}=\frac{2g\kappa _g1}{\kappa _0\kappa _g1}\hfill \\ \hfill \nu _0& =\frac{2g+\kappa _g+1}{\kappa _0\kappa _g1}\hfill \end{array}\text{ and }\{\begin{array}{cc}\hfill N_g& =\frac{\kappa _0N}{\kappa _0\kappa _g1}=\frac{\kappa _0(2g\kappa _g1)}{\kappa _0\kappa _g1}\hfill \\ \hfill \nu _g& =\frac{\kappa _0(12g)+1}{\kappa _0\kappa _g1}\hfill \end{array}$$
Remark that $`\nu _0+N_0=0`$ (which, as you can guess, is forced by our choice of $`N`$); nevertheless $`e_0(ND)`$ will be a rational number. We have by definition that
$`z_0(ND,s)`$ $`={\displaystyle \frac{1}{\nu _0+sN_0}}+{\displaystyle \frac{1}{(\nu _0+sN_0)(\nu _g+sN_g)}}+{\displaystyle \frac{2g}{\nu _g+sN_g}}+{\displaystyle \frac{1}{(\nu _g+sN_g)(1+sN)}}`$
$`={\displaystyle \frac{1+(\kappa _02g)(1+sN)}{(\nu _g+sN_g)(1+sN)}}.`$
The fact that $`\nu _0+sN_0`$ cancels in the denominator is a general fact; see \[V3, 2.2\]. Plugging in the expression for $`\nu _g`$ and $`N_g`$ yields
$$z_0(ND,s)=\frac{(\kappa _0\kappa _g1)[1+(\kappa _02g)(1+sN)]}{(12\kappa _0g+\kappa _0(1+sN))(1+sN)}\text{ (with }N=2g\kappa _g1)$$
and
$$e_0(ND)=\underset{s1}{lim}z_0(ND,s)=\frac{(2g\kappa _0)(2g\kappa _g)1}{2g\kappa _g}.$$
One can analogously compute $`Z_0(ND,s)`$ and $`E_0(ND)`$.
4.2. Let $`0X`$ and $`h:YX`$ be as above with $`g=1`$ (instead of $`g2`$). Now let $`E^{}`$ be a nonsingular curve germ in $`Y`$ intersecting $`E_0`$ transversely in one point and disjoint from $`E_1`$, and denote $`D^{}=h(E^{})`$. See Figure 2.
....................................................................................................................................................................................................................................................................................................................................$``$$`D^{}`$$`0`$$`X`$$``$$`h`$Figure 2 $`E_0`$$`E_1`$$`E^{}`$$`Y`$
One easily computes (see \[V3, 2.5\]) that
$$z_0(ND^{},s)=\frac{\kappa _0\kappa _11}{1+sN}\text{and thus}e_0(ND^{})=\frac{\kappa _0\kappa _11}{1+N}.$$
Now choose $`N=\kappa _01`$. It is easy to verify that then $`\nu _1+N_1=0`$; so as in 4.1 we could not have defined $`e_0((\kappa _01)D^{})`$ by the usual formula. However in this example our definition on the level of Hodge polynomials yields $`E_0((\kappa _01)D^{})=\mathrm{}`$.
4.2.1. Remark. One could argue whether in Definition–Proposition 3.8 (and analogously in 3.12) it is more appropriate to introduce $`e_W(D)`$ as $`lim_{u,v1}E_W(D)`$. When $`E_W(D)\mathrm{}`$ this amounts to the same, but when $`E_W(D)=\mathrm{}`$ we then would miss some interesting values of $`e_W(D)`$ as in 4.2 above.
4.3. Let $`X`$ be the quadric hypersurface $`\{xyzw=0\}`$ in $`𝔸^4`$. The origin $`0`$ is the only singular point of $`X`$. Blowing up $`0`$ yields a log resolution $`h_1:Y_1X`$ of $`X`$, which is an isomorphism outside $`h^1\{0\}`$ and with $`E_1=h^1\{0\}(\{xyzw=0\}^3)^1\times ^1`$.
(a) Consider the divisor $`D=E+E^{}`$ on $`X`$, where $`E`$ and $`E^{}`$ are the zero sets of the functions $`zw`$ and $`y`$ on $`X`$, respectively. Remark that $`E`$ is irreducible and that $`E^{}`$ consists of two irreducible components. We want to compute $`z_0(D,s)`$. In this example we will use the same notation for divisors and their strict transforms by blowing–ups.
................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................. ........................................................................................... $`E^{}`$$`E^{}`$$`E`$$`E_1`$Figure 3
In Figure 3 we present the intersection configuration of $`E_1,E`$ and $`E^{}`$ on $`Y_1`$. The variety $`Y_1`$ is naturally covered by 4 affine charts, each isomorphic to $`𝔸^3`$. In the ‘main chart’ the exceptional surface $`E_1`$ and the strict transforms $`E`$ and $`E^{}`$ are given in affine coordinates $`x,z,w`$ by
$`E_1:x=0,`$
$`E:zw=0,`$
$`E^{}:zw=0`$
(in the other charts $`E`$ and $`E^{}`$ do not intersect).
We obtain a log resolution $`h`$ of $`D`$ by composing $`h_1`$ with the blowing–up $`h_2:Y_2Y_1`$ of the curve $`EE^{}(𝔸^1)`$ in $`Y_1`$. The exceptional variety $`E_2`$ of $`h_2`$ is isomorphic to $`𝔸^1\times ^1`$; the intersection configuration of $`E_2,E_1,E`$ and $`E^{}`$ is presented in Figure 4.
............................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................... .............................................................................................................................................................................................................$`E^{}`$$`E^{}`$$`E`$$`E_1`$$`E_2`$Figure 4
Denoting as usual $`h^{}D=E+E^{}+N_1E_1+N_2E_2`$ and $`K_Y=h^{}K_X+(\nu _11)E_1+(\nu _21)E_2`$, one easily verifies that $`(\nu _1,N_1)=(2,2)`$ and $`(\nu _2,N_2)=(2,3)`$. The contributors to $`z_0(D,s)`$ are $`E_1^{},(E_1E_2)^{},(E_1E)^{},(E_1E^{})^{},E_1E_2E`$ and $`E_1E_2E^{}`$. Now $`\chi (E_1^{})=0`$ and the other Euler characteristics are obvious; then
$`z_0(D,s)`$ $`={\displaystyle \frac{1}{\nu _1+sN_1}}\left({\displaystyle \frac{1}{\nu _2+sN_2}}+{\displaystyle \frac{3}{1+s}}+{\displaystyle \frac{3}{(\nu _2+sN_2)(1+s)}}\right)`$
$`={\displaystyle \frac{4}{(2+3s)(1+s)}}.`$
Also $`e_0(D)=lim_{s1}z_0(D,s)=\frac{2}{5}`$ and $`e_0\left((X,D)\right)=lim_{s1}z_0(D,s)=\mathrm{}.`$
(b) Now consider the $``$–divisor $`D=NE+N^{}E^{}`$ with $`N>0,N^{}>0,N1,N^{}1`$ and $`N+N^{}=2`$. The morphism $`h:Y_2X`$ in (a) is of course still a log resolution of $`D`$. The only difference with the data in (a) is that here $`h^{}D=NE+NE^{}+N_1E_1+N_2E_2`$ with $`N_1=N+N^{}=2`$ and $`N_2=N+2N^{}=2+N^{}`$. So
$`z_0(D,s)=`$ $`{\displaystyle \frac{1}{\nu _1+sN_1}}({\displaystyle \frac{1}{\nu _2+sN_2}}+{\displaystyle \frac{1}{1+sN}}+{\displaystyle \frac{2}{1+sN^{}}}+{\displaystyle \frac{1}{(\nu _2+sN_2)(1+sN)}}`$
$`+{\displaystyle \frac{2}{(\nu _2+sN_2)(1+sN^{})}})`$
$`=`$ $`{\displaystyle \frac{1}{2+2s}}{\displaystyle \frac{8+16s+8s^2}{(2+s(2+N^{}))(1+sN)(1+sN^{})}}`$
$`=`$ $`{\displaystyle \frac{4(1+s)}{(2+s(2+N^{}))(1+sN)(1+sN^{})}}.`$
And then $`e_0(D,s)=\frac{8}{(4+N^{})(1+N)(1+N^{})}`$ and $`e_0\left((X,D)\right)=0`$.
4.4. Fix $`d,d3`$. Take a homogeneous polynomial $`F`$ in $`d+1`$ variables of degree $`a2`$ such that $`\{F=0\}^d`$ is nonsingular.
Let $`X`$ be the hypersurface in $`𝔸^{d+1}`$ given by the zero set of $`F`$; so $`X`$ is the affine cone over $`\{F=0\}^d`$ and the origin is the only singular point of $`X`$. Let $`D`$ be the intersection of $`X`$ with a general hyperplane through the origin in $`𝔸^{d+1}`$. The blowing–up $`h:YX`$ of the origin yields a log resolution of $`X`$, which is moreover a log resolution of $`D`$. We denote the strict transform of $`D`$ by $`E`$, and the exceptional variety of $`h`$ by $`E_1`$. Notice that $`E_1`$ is isomorphic to $`\{F=0\}^d`$. We try to give an impression of this situation in Figure 5.
.........................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................$``$$`D`$$`0`$$`X`$..........................................................................................................................................................................$``$$`h`$Figure 5.............................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................. $`E_1`$$`E`$$`Y`$..........................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................
As usual we denote $`K_Y=h^{}K_X+(\nu _11)E_1`$ and $`h^{}(ND)=NE+N_1E_1`$ for $`N`$, $`N>0`$. One can verify that $`\nu _1=d+1a`$ and $`N_1=N`$.
To compute $`z_X(ND,s)`$ we need the Euler characteristics of the varieties $`Y^{},E^{},E_1^{}`$ and $`EE_1`$ (which stratify $`Y`$). Since $`X`$ and $`D`$ are affine cones we have that
$$\chi (E^{})=\chi (D\{0\})=0\text{ and }\chi (Y^{})=\chi (XD)=0.$$
Now $`E_1`$ is a nonsingular hypersurface of degree $`a`$ in $`^d`$, yielding
$$\chi (E_1)=(1a)\left(\frac{(1a)^d1}{a}\right)+d$$
(see for example \[Hirz\]). And because $`D`$ was chosen to be general we have moreover that $`EE_1`$ is a nonsingular hypersurface of degree $`a`$ in $`^{d1}`$; so
$$\chi (EE_1)=(1a)\left(\frac{(1a)^{d1}1}{a}\right)+d1.$$
Then finally $`\chi (E_1^{})=\chi (E_1)\chi (EE_1)=(1a)^d+1`$ and
$`z_X(ND,s)=`$ $`z_0(ND,s)={\displaystyle \frac{\chi (E_1^{})}{\nu _1+sN_1}}+{\displaystyle \frac{\chi (EE_1)}{(\nu _1+sN_1)(1+sN)}}`$
$`=`$ $`{\displaystyle \frac{(1a)^d+1}{d+1a+sN}}+{\displaystyle \frac{(1a)\left(\frac{(1a)^{d1}1}{a}\right)+d1}{(d+1a+sN)(1+sN)}}`$
$`=`$ $`{\displaystyle \frac{(1a)\left(\frac{(1a)^d1}{a}\right)+d+s\left(1(1a)^d\right)N}{(d+1a+sN)(1+sN)}}.`$
A (not very exciting) calculation shows that there is no cancellation in this expression, except when $`d=3`$ and $`a=2`$ or $`3`$, in which case $`z_X(ND,s)`$ is
$$\frac{2}{1+sN}\text{and }\frac{9}{1+sN},$$
respectively. Taking limits we obtain
$$e_X(ND)=\frac{(1a)\left(\frac{(1a)^d1}{a}\right)+d+\left(1(1a)^d\right)N}{(d+1a+N)(1+N)}\text{ if }d+1+Na$$
and
$$e_X\left((X,ND)\right)=\frac{(1a)\left(\frac{(1a)^d1}{a}\right)+d+\left((1a)^d1\right)N}{(d+1aN)(1N)}\text{ if }\{\begin{array}{cc}d+1a+N\hfill & \\ N1.\hfill & \end{array}$$
## 5. Zeta functions associated to divisors and differential forms
5.1. In the $`p`$–adic theory of Igusa’s local zeta functions one also associates this invariant to both polynomials and differential forms, see e.g. \[L, III3.5\]. Let $`f_p[x]=_p[x_1,\mathrm{},x_d]`$ and $`w\mathrm{\Omega }_{𝔸^d}^d`$, i.e. $`w=gdx`$ where $`g_p[x]`$ and $`dx=dx_1\mathrm{}dx_d`$. Then, with the notation of 1.3.1,
$$Z_p(f,w,s):=_{_p^d}|f(x)|^s|g(x)||dx|.$$
With the notation of 1.4 let $`\nu _i^{}1`$ be the multiplicity of $`E_i`$ in the divisor of $`h^{}w`$. Then (for $`f[x]`$) the same formula as in 1.4.1 is valid when we replace $`\nu _i`$ by $`\nu _i^{}`$.
We also want to introduce this notion on the level of the Grothendieck ring of algebraic varieties as in 1.3. Our motivation in this paper is that we will use it to construct on a $``$–Gorenstein variety $`X`$ an invariant $`𝒵_W(X,D,s)`$, generalizing $`Z_W(X,D,s)`$ in 3.3, on the level of the Grothendieck ring. Furthermore we will need this notion in future work.
5.2. We fix notations for this section. Let $`X`$ be an irreducible nonsingular variety of dimension $`d`$ and $`W`$ a subvariety of $`X`$. Let $`D`$ be an effective divisor on $`X`$ and $`J\mathrm{\Omega }_X^d`$ an invertible subsheaf of the sheaf of regular differential $`d`$–forms $`\mathrm{\Omega }_X^d`$ on $`X`$.
We will only consider the situation where supp $`J\mathrm{supp}D`$; we motivate this below.
5.3. First we rephrase the definition of $`𝒵_W(D,s)`$ in terms of the motivic volume $`\mu `$ of \[DL3, 3.2\] or \[DL4\]. Denote by $`𝒞`$ the family of subsets of $`(X)`$ of the form $`\pi _n^1A_n`$ for some $`n`$ and constructible subset $`A_n`$ of $`_n(X)`$. We call these cylindrical subsets as in \[B2\] or \[DL4\]. There exists a unique additive measure $`\mu :𝒞_L`$ satisfying $`\mu (\pi _n^1A_n)=\frac{[A_n]}{L^{(n+1)d}}`$ for $`A_n`$ as above. (In fact this map is denoted by $`\stackrel{~}{\mu }`$ in \[DL3\] and there $`\mu `$ is a map from the more complicated family of semi–algebraic subsets of $`(X)`$ to $`\widehat{M}`$.) For $`A`$ in $`𝒞`$ and $`\alpha :A`$ a bounded function with cylindrical fibres one defines the integral
$$_AL^\alpha 𝑑\mu :=\underset{n}{}L^n\mu (\alpha ^1\{n\})_L.$$
$`\mathrm{5.3.1}`$
Now re–examining the definition of $`𝒵_W(D,s)`$ in 1.3 we have, with the notation introduced there, that $`\mu (Y_{n,D,W})=[X_{n,D,W}]L^{(n+1)d}`$ and hence
$$𝒵_W(D,s)=\underset{n}{}\mu (Y_{n,D,W})L^{ns}_L[[L^s]].$$
5.4. The following construction is a special case of \[DL3, 3.5\]. To the sheaf $`J`$ is associated as follows a measure $`\mu _J`$ on $`𝒞`$, such that $`\mu _{\mathrm{\Omega }_X^d}=\mu `$.
For $`PX`$ let $`dx`$ and $`g_Pdx`$ be local generators of $`\mathrm{\Omega }_X^d`$ and $`J`$, respectively, around $`P`$. Denote then by $`\mathrm{ord}_tJ:(X)\{\mathrm{}\}`$ the function assigning to $`\phi `$ in $`(X)`$ the order of the power series given by $`g_{\pi _0(\phi )}\phi `$. For $`A`$ in $`𝒞`$ we define
$$\mu _J(A):=_AL^{\mathrm{ord}_tJ}𝑑\mu =\underset{\mathrm{}}{}L^{\mathrm{}}\mu (A\{\mathrm{ord}_tJ=\mathrm{}\}).$$
Indeed the sets $`\{\mathrm{ord}_tJ=\mathrm{}\}`$ are cylindrical. For arbitary $`A`$ the right hand side above is only defined as an element in $`\widehat{}`$; however we will only consider sets $`A`$ for which the sum over $`\mathrm{}`$ is finite and then $`\mu _J(A)_L`$. Replacing $`\mu `$ by $`\mu _J`$ we can consider analogous integrals as in (5.3.1).
The following change of variables formula is a special case of \[DL3, 3.5.2\]. (It follows immediately from \[DL3, 3.3\] of which Theorem 1.10 is a special case.)
###### 5.4.1. Proposition
Let $`X^{}`$ be another irreducible smooth variety and $`\rho :X^{}X`$ a proper birational morphism. For $`A`$ in $`𝒞`$ and $`\alpha :A`$ a bounded function with cylindrical fibres we have that
$$_AL^\alpha 𝑑\mu _J=_{\rho ^1A}L^{\alpha \rho }𝑑\mu _{\rho ^{}J}.$$
###### 5.5. Definition
To the data of 5.2 we associate the motivic zeta function
$`𝒵_W(D,J,s)`$ $`=𝒵_W(X,D,J,s):={\displaystyle \underset{n}{}}\mu _J(Y_{n,D,W})L^{ns}`$
$`={\displaystyle \underset{n}{}}\left({\displaystyle \underset{\mathrm{}}{}}L^{\mathrm{}}\mu (Y_{n,D,W}\{\mathrm{ord}_tJ=\mathrm{}\})\right)L^{ns}_L[[L^s]].`$
We explain why the sum over $`\mathrm{}`$ is finite. For any fixed $`n`$ we have that $`Y_{n,D,W}=_{\mathrm{}\{\mathrm{}\}}(Y_{n,D,W}\{\mathrm{ord}_tJ=\mathrm{}\})`$. But our condition $`\mathrm{supp}J\mathrm{supp}D`$ implies that
$$\{\mathrm{ord}_tJ=\mathrm{}\}=(\mathrm{supp}J)(\mathrm{supp}D)=\{\mathrm{ord}_tD=\mathrm{}\},$$
hence we have that $`Y_{n,D,W}\{\mathrm{ord}_tJ=\mathrm{}\}=\mathrm{}`$ and so $`Y_{n,D,W}`$ is the countable union of the cylindrical sets $`Y_{n,D,W}\{\mathrm{ord}_tJ=\mathrm{}\},\mathrm{}`$. Then this union is finite by \[B2, Theorem 6.6\].
###### 5.6. Theorem
Let $`X^{}`$ be another irreducible smooth variety and $`\rho :X^{}X`$ a proper birational morphism. Then
$$𝒵_W(X,D,J,s)=𝒵_{\rho ^1W}(X^{},\rho ^{}D,\rho ^{}J,s).$$
###### Demonstration Proof
This is a consequence of Proposition 5.4.1. ∎
###### 5.7. Theorem
Let $`h:YX`$ be a log resolution of $`\mathrm{supp}D`$. Denote as usual the irreducible components of $`h^1(\mathrm{supp}D)`$ by $`E_i,iT`$. We set $`h^{}D=_{iT}N_iE_i`$ and $`\mathrm{div}(h^{}w)=_{iT}(\nu _i^{}1)E_i`$, where $`w`$ is a local generator of $`J`$. Then
$$𝒵_W(D,J,s)=L^d\underset{IT}{}[E_I^{}h^1W]\underset{iI}{}\frac{L1}{L^{\nu _i^{}+sN_i}1}.$$
Remark. Let as in 1.4 $`dx`$ be a local generator of $`\mathrm{\Omega }_X^d`$ and $`\mathrm{div}(h^{}dx)=_{iT}(\nu _i1)E_i`$. Say $`w=gdx`$ and $`\mathrm{div}(h^{}g)=_{iT}M_iE_i`$. Then $`\nu _i^{}=\nu _i+M_i`$ for $`iT`$.
###### Demonstration Proof
One can adapt the proof of \[DL2, Theorem 2.2.1\] completely to this more general setting with the sheaf $`J`$. ∎
5.8. The notion introduced above is sufficient to introduce zeta functions on the level of the Grothendieck ring for Gorenstein varieties. To cover the case of $``$–Gorenstein varieties we need ‘sheaves of multivalued differential forms’. We briefly describe this generalization.
Now let $`J(\mathrm{\Omega }_X^d)^m`$ be an invertible subsheaf of the $`m`$–fold tensor product of $`\mathrm{\Omega }_X^d`$, still satisfying $`\mathrm{supp}J\mathrm{supp}D`$. We define
$$\mu _{J^{1/m}}(A):=_AL^{\frac{\mathrm{ord}_tJ}{m}}𝑑\mu =\underset{\mathrm{}}{}L^{\mathrm{}/m}\mu (A\{\mathrm{ord}_tJ=\mathrm{}\})_L[L^{1/m}]$$
for the sets $`A`$ in $`𝒞`$ for which the last sum is finite. Then the motivic zeta function is
$`𝒵_W(D,J^{1/m},s)`$ $`=𝒵_W(X,D,J^{1/m},s)`$
$`:={\displaystyle \underset{n}{}}\mu _{J^{1/m}}(Y_{n,D,W})L^{ns}_L[L^{1/m}][[L^s]].`$
Theorem 5.7 easily generalizes to this setting, but now the $`\nu _i^{}\frac{1}{m}(\{0\})`$.
5.9. Finally as in 1.5 we can generalize further to $``$–divisors. Now if $`D`$ is an effective $``$–divisor on $`X`$, such that $`rD`$ is a divisor for an $`r\{0\}`$, we define $`𝒵_W(D,J^{1/m},s):=𝒵_W(rD,J^{1/m},s/r)`$. Again Theorem 5.7 generalizes, with now the $`N_i\frac{1}{r}(\{0\})`$.
## 6. Singular varieties; on the level of the Grothendieck ring
6.1. In this section we generalize the zeta function of 3.3 to the level of the Grothendieck ring. In order to focus on the main idea we first treat the essential case, being an effective Cartier divisor $`D`$ on a Gorenstein variety $`X`$. For a normal variety $`V`$ we denote its canonical sheaf (corresponding to $`K_V`$) by $`\omega _V`$; we have that $`\omega _V`$ is invertible or $`\omega _V^m`$ is invertible for some $`m\{0\}`$ precisely when $`V`$ is Gorenstein or $``$–Gorenstein, respectively.
Also in the sequel $`(F)`$ denotes the sheaf of ideals associated to an effective divisor $`F`$ on a nonsingular variety.
###### 6.2. Definition – Proposition
Let $`X`$ be a Gorenstein variety of dimension $`d`$ and $`W`$ a subvariety of $`X`$. Let $`D`$ be an effective Cartier divisor on $`X`$ such that $`X_{\mathrm{sing}}\mathrm{supp}D`$.
(i) The motivic zeta function
$$𝒵_W(D,s)=𝒵_W(X,D,s):=𝒵_{h^1W}(Y,h^{}D,h^{}\omega _X(ah^{}D),sa)$$
where $`h:YX`$ is a log resolution of $`\mathrm{supp}D`$ and $`a,a>>0`$.
(ii) Let $`h:YX`$ be a log resolution of $`\mathrm{supp}D`$. With the notation of 3.2 we have that
$$𝒵_W(D,s)=L^d\underset{IT}{}[E_I^{}h^1W]\underset{iI}{}\frac{L1}{L^{\nu _i+sN_i}1}.$$
###### Demonstration Proof
We first explain the right hand side of our definition. Since $`X_{\mathrm{sing}}\mathrm{supp}D`$ we have that $`\mathrm{supp}(h^{}\omega _X)\mathrm{supp}(h^{}D)`$, yielding for $`a>>0`$ that $`h^{}\omega _X(ah^{}D)`$ is an invertible subsheaf of $`\mathrm{\Omega }_Y^d`$. So to this sheaf and the effective divisor $`h^{}D`$ we can associate the motivic zeta function of 5.5. The substitution ‘$`sa`$ instead of $`s`$’ means replacing the variable $`L^s`$ by $`L^a(L^s)`$.
Now we show independency of the chosen resolution; it is sufficient to consider another log resolution $`h^{}:Y^{}X`$ that factors as $`h^{}:Y^{}\stackrel{𝜑}{}Y\stackrel{}{}X`$. By Theorem 5.6 we indeed have that
$`𝒵_{h^1W}`$ $`(Y,h^{}D,h^{}\omega _X(ah^{}D),sa)`$
$`=𝒵_{\phi ^1(h^1W)}(Y^{},\phi ^{}h^{}D,\phi ^{}(h^{}\omega _X)\phi ^{}((ah^{}D)),sa)`$
$`=𝒵_{h^1W}(Y^{},h^{}D,h^{}\omega _X(ah^{}D),sa).`$
Let $`\eta `$ and $`f`$ be local generators of $`\omega _X`$ and $`(D)`$, respectively. Then $`(h^{}f)^a(h^{}\eta )`$ is a local generator of $`h^{}\omega _X(ah^{}D)`$ and its divisor of zeroes is $`_{iT}((\nu _i1)+aN_i)E_i`$. Hence Theorem 5.7 (with $`h=Id_Y`$) yields the stated formula for $`𝒵_W(D,s)`$, which also proves independency of the number $`a`$. ∎
6.3. Now let $`X`$ be $``$–Gorenstein and say that $`mK_X`$ is Cartier for some $`m\{0\}`$. We define $`𝒵_W(X,D,s)`$ just as in 6.2, interpreting the expression $`h^{}\omega _X(ah^{}D)`$ as an abbreviation of $`(h^{}(\omega _X^m)(mah^{}D))^{1/m}`$ (see 5.8). Now $`𝒵_W(X,D,s)`$ lives in a localization of $`_L[L^{1/m}][L^s]`$ and is given by the same formula as in 6.2 (with now the $`\nu _i`$).
When $`D`$ is an effective $``$–Cartier divisor we set as usual $`𝒵_W(D,s):=𝒵_W(rD,s/r)`$ if $`rD`$ is Cartier for an $`r\{0\}`$. Then in full generality we have the following.
###### 6.4. Definition – Proposition
Let $`X`$ be a $``$–Gorenstein variety of dimension $`d`$ and $`W`$ a subvariety of $`X`$. Let $`D`$ be an effective $``$–Cartier divisor on $`X`$ (with $`rD`$ Cartier for an $`r\{0\})`$ such that $`X_{\mathrm{sing}}\mathrm{supp}D`$. The motivic zeta function
$$𝒵_W(D,s)=𝒵_W(X,D,s):=𝒵_{h^1W}(Y,h^{}(rD),h^{}\omega _X(arh^{}D),s/ra)$$
where $`h:YX`$ is a log resolution of $`\mathrm{supp}D`$ and $`a,a>>0`$. We have the same formula as in 6.2.
Of course $`𝒵_W(D,s)`$ specializes to the zeta function $`Z_W(D,s)`$ of 3.3.
6.5. Finally we consider for arbitrary $``$–Gorenstein varieties $`X`$ ‘Kontsevich’ invariants $`_W(D)`$ and $`_W\left((X,D)\right)`$ on the level of the Grothendieck ring, which specialize to $`E_W(D)`$ and $`E_W\left((X,D)\right)`$ of 3.6 and 3.10, respectively. Notice first that in 6.4 we have, by the formula for $`𝒵_W(D,s)`$ in terms of a log resolution, that it already belongs to the localization of a polynomial ring $`_L[L^{1/r}][L^{s/r}]`$ with respect to $`(1L^{\alpha \beta s})_{\alpha ,\beta _{>0}}`$. Morally we again take limits for $`s1`$ and $`s1`$ to define $`_W(D)`$ and $`_W\left((X,D)\right)`$, respectively.
###### 6.6. Definition
Let $`X`$ be a $``$–Gorenstein variety of dimension $`d`$ and $`W`$ a subvariety of $`X`$. Let $`D`$ be an effective $``$–Cartier divisor on $`X`$ such that $`X_{\mathrm{sing}}\mathrm{supp}D`$. Take $`r\{0\}`$ with $`rK_X`$ and $`rD`$ Cartier.
(i) If $`𝒵_W(D,s)`$ belongs to the localization of $`_L[L^{1/r}][L^{s/r}]`$ with respect to $`(1L^{\alpha \beta s})_{\alpha ,\beta _{>0},\alpha +\beta 0}`$, then we put
$$_W(D)=_W(X,D):=𝒵_W(D,s)|_{s=1}.$$
Otherwise we put $`_W(D)=_W(X,D):=\mathrm{}`$.
(ii) If $`𝒵_W(D,s)`$ belongs to the localization of $`_L[L^{1/r}][L^{s/r}]`$ with respect to $`(1L^{\alpha \beta s})_{\alpha ,\beta _{>0},\alpha \beta }`$, then we put
$$_W\left((X,D)\right):=𝒵_W(D,s)|_{s=1}.$$
Otherwise we put $`_W\left((X,D)\right):=\mathrm{}`$.
Here the evaluations $`s=1`$ and $`s=1`$ mean substituting the variable $`L^{s/r}`$ by $`L^{1/r}`$ and $`L^{1/r}`$, respectively, yielding a well defined element in $`\widehat{}[L^{1/r}]`$.
###### 6.7. Proposition
Consider the same data as is 6.6.
(i) Suppose there is a log resolution $`h:YX`$ of $`\mathrm{supp}D`$ for which $`\nu _i+N_i0`$ for all $`iT`$ (using the notation of 3.2). Then
$$_W(D)=L^d\underset{IT}{}[E_I^{}h^1W]\underset{iI}{}\frac{L1}{L^{\nu _i+N_i}1}.$$
(ii) Suppose there is a log resolution $`h:YX`$ of $`\mathrm{supp}D`$ for which all log discrepancies $`a_i,iT`$, with respect to the pair $`(X,D)`$ satisfy $`a_i0`$ (using the notation of 3.2). Then
$$_W\left((X,D)\right)=L^d\underset{IT}{}[E_I^{}h^1W]\underset{iI}{}\frac{L1}{L^{a_i}1}.$$
## Appendix
A1. Let in this appendix $`X`$ be a smooth irreducible variety of dimension $`d`$ and $`W`$ a subvariety of $`X`$.
In 1.7 – 1.10 we described the Kontsevich invariant $`_W(D)\widehat{}`$, associated to an effective divisor $`D`$ on $`X`$, and we mentioned its important properties. Here we will generalize this notion to effective $``$–divisors; if $`rD`$ is a divisor for an $`r\{0\}`$ we obtain an invariant $`_W(D)`$ in a finite extension $`\widehat{}[L^{1/r}]`$ of $`\widehat{}`$, and we treat analogous properties. We also introduce this invariant for a $``$–divisor $`D=_id_iD_i`$ (with the $`D_i`$ irreducible) such that all $`d_i>1`$ and $`\mathrm{supp}D=_iD_i`$ is a divisor with strict normal crossings. This is used in 3.14.
A2. First we describe the ring $`\widehat{}[L^{1/r}]`$. Consider the integral ring extension $`_L_L[L^{1/r}]:=\frac{_L[X]}{(X^rL)}`$, where $`L^{1/r}`$ is the class of $`X`$ in this quotient. Each element $`a_L[L^{1/r}]`$ has a unique expression of the form $`a=_{i=0}^{r1}a_iL^{i/r}`$ or $`a=_{i=0}^{r1}a_i^{}L^{i/r}`$ with $`a_i,a_i^{}_L`$.
We extend the decreasing filtration $`(F_m)_m`$ on $`_L`$, introduced in 1.7, to the ring $`_L[L^{1/r}]`$. Let $`F_m^{},m`$, be the subgroup of $`_L[L^{1/r}]`$ generated by
$$\left\{\underset{i=0}{\overset{r1}{}}\frac{[A_i]}{L^{n_i}}L^{i/r}\right|\mathrm{dim}A_in_im\text{ for }i=0,\mathrm{},r1\}.$$
(So indeed $`F_m=_LF_m^{}`$.) We take the completion $`\widehat{}^{}`$ of $`_L[L^{1/r}]`$ with respect to this filtration $`(F_m^{})_m`$; then we have an injection $`\widehat{}\widehat{}^{}`$.
One can verify that $`\widehat{}^{}\widehat{}[L^{1/r}]`$, where the right hand side can be interpreted either as the subring of $`\widehat{}^{}`$ generated by $`\widehat{}`$ and $`L^{1/r}`$, or as $`\frac{\widehat{}[X]}{(X^rL)}`$.
A3. We will use the following notation. Let $`D`$ be a prime divisor on $`X`$. Then $`\mathrm{ord}_tD:(X)\{\mathrm{}\}`$ assigns to $`\phi (X)`$ the order of the power series in $`t`$ given by $`f\phi `$, where $`f`$ is a local equation of $`D`$ at $`\pi _0(\phi )`$. For a $``$–divisor $`D=_id_iD_i`$ (with the $`D_i`$ prime divisors) we then define $`\mathrm{ord}_tD:(X)\{\mathrm{}\}`$ by $`\mathrm{ord}_tD:=_id_i\mathrm{ord}_tD_i`$.
A4. Definition. Let $`D`$ be a $``$–divisor on $`X`$ and $`r\{0\}`$ such that $`rD`$ is a divisor.
($`i`$) If $`D`$ is effective we define for $`n`$ the subscheme $`Y_{n,D,w}`$ of $`(X)`$ and the subscheme $`X_{n,D,w}`$ of $`_n(X)`$ as in 1.3 with only the following adaptation : now $`f`$ is a local equation of the divisor $`rD`$ (instead of $`D`$). Then we set
$$_W(D)=_W(X,D):=\underset{n}{}\frac{[X_{n,D,W}]}{L^{(n+1)d}}L^{n/r}\widehat{}[L^{1/r}].$$
In terms of the motivic volume $`\mu `$ of 5.3 we can describe $`_W(D)`$ as
$$_W(D)=_{\pi _0^1W}L^{\mathrm{ord}_tD}𝑑\mu :=\underset{n}{}\mu (\pi _0^1W\{\mathrm{ord}_tD=\frac{n}{r}\})L^{n/r}.$$
($`ii`$) In general we say that $`\mathrm{ord}_tD:(X)\frac{1}{r}\{\mathrm{}\}`$ is integrable on $`\pi _0^1W`$ if
$$_{\pi _0^1W}L^{\mathrm{ord}_tD}𝑑\mu :=\underset{n}{}\mu (\pi _0^1W\{\mathrm{ord}_tD=\frac{n}{r}\})L^{n/r}$$
converges in $`\widehat{}[L^{1/r}]`$; we then denote this invariant again by $`_W(D)`$.
A5. An important case of this last definition occurs when $`D=_{i=1}^kd_iD_i`$ with the $`D_i`$ irreducible, all $`d_i>1`$, and supp $`D=_{i=1}^kD_i`$ a divisor with strict normal crossings. For $`J\{1,\mathrm{},k\}`$ denote $`D_J^{}:=(_{jJ}D_j)(_\mathrm{}JD_{\mathrm{}})`$ and $`M_J:=\{(m_1,\mathrm{},m_k)^k`$ $`m_j>0jJ\}`$. Then one can compute that
$$_{\pi _0^1W}L^{\mathrm{ord}_tD}𝑑\mu =L^d\underset{J\{1,\mathrm{},k\}}{}(L1)^{|J|}[D_J^{}W]\underset{(m_1,\mathrm{},m_k)M_J}{}L^{_{jJ}(d_j+1)m_j},$$
which converges in $`\widehat{}[L^{1/r}]`$ since all $`d_j+1>0`$. See \[B2, Theorem 6.28\] and \[C, Theorem 1.17\].
###### \bfA6. Theorem
Let also $`X^{}`$ be a smooth irreducible variety and $`\rho :X^{}X`$ a proper birational morphism. Let $`D`$ be a $``$–divisor on $`X`$. Then $`\mathrm{ord}_tD`$ is integrable on $`\pi _0^1W`$ if and only if $`\mathrm{ord}_t(\rho ^{}D+K_{X^{}|X})`$ is integrable on $`\pi _0^1(\rho ^1W)`$; and in this case
$$_W(X,D)=_{\rho ^1W}(X^{},\rho ^{}D+K_{X^{}|X}).$$
###### Demonstration Proof
The proof of \[DL3, Lemma 3.3\], based on the crucial and difficult \[DL3, Lemma 3.4\], can be adapted to this setting. See \[B2, Theorem 6.27\] for an analogous statement and proof when $`W=X`$. We also remark that when $`D`$ is an effective $``$–divisor (implying that both functions are integrable), then one can prove the stated equality using the equality in \[DL3, Lemma 3.3\]. ∎
###### A7. Theorem
Let $`D`$ be a $``$–divisor on $`X`$ (with $`rD`$ a divisor for an $`r\{0\}`$) such that $`\mathrm{ord}_tD`$ is integrable on $`\pi _0^1W`$. Using the notation of 1.4 we have the following formula for $`_W(D)`$ in terms of a log resolution $`h:YX`$ of supp $`D`$ :
$$_W(D)=L^d\underset{IT}{}[E_I^{}h^1W]\underset{iI}{}\frac{L1}{L^{\nu _i+N_i}1}\text{ in }\widehat{}[L^{1/r}].$$
In particular $`_W(D)`$ belongs to the image of $`_L[(1L^{n/r})^1]_{n\{0\}}`$ in $`\widehat{}[L^{1/r}]`$.
###### Demonstration Proof
This follows from A5 and Theorem A6. One can also adapt \[DL3, (6.5)\]. ∎
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# Partially Quenched Chiral Perturbation Theory and the Replica Method
## 1 Introduction
The question of non-perturbative analytical predictions for quenched or partially quenched lattice gauge theory computations has been thoroughly studied in the context of effective chiral Lagrangians . So far the most systematic framework has been the supersymmetric formulation of Bernard and Golterman , which builds on an idea first introduced in the context of staggered lattice fermions . Here one introduces $`k`$ additional quark species (of conventional statistics) on top of the $`N_f`$ physical “sea” quarks, and $`k`$ “ghost” quarks of opposite statistics to cancel the effects of the additional quarks. When $`N_f`$ is taken to vanish this gives the fully quenched theory, while for $`N_f`$ non-zero it gives the partially quenched theory. Both are accessible to a study by Monte Carlo techniques in lattice gauge theory. The chiral flavor symmetry group is in that formulation extended to a super Lie group which in perturbation theory can be taken as SU($`N_f+k|k`$). (For this reason it is commonly known as the supersymmetric method although it, as applied, has nothing to do with space-time supersymmetry, but rather is a graded symmetry.) Based on the usual assumption of spontaneous chiral symmetry breaking (here extended to the super group case) the effective low-energy theory of the lowest-lying hadronic excitations is that of a chiral Lagrangian, now with fields living on the coset of super Lie groups. This effective Lagrangian can be studied by the conventional methods of chiral perturbation theory. In what follows we denote fully and partially quenched chiral perturbation theory by QChPT and PQChPT, respectively.
The supersymmetric framework has also proven to be an efficient means of deriving analytical results for the soft part of the Dirac operator spectrum in finite volume, by taking an appropriate discontinuity of the partially quenched chiral condensate . This has brought earlier results derived entirely from universal Random Matrix Theory (for a very recent comprehensive review, see ref. ) in direct contact with the effective Lagrangian of QCD. In particular, a series of very compact relations that described general $`k`$-point spectral correlation functions of low-lying Dirac operator eigenvalues in terms of effective partition functions with additional quark species can now be understood as due to the cancelling pairs of fermionic and bosonic valence quarks. When taking the same discontinuity near the origin in PQChPT it has also been shown that one recovers among other terms the analytical prediction for the slope of the spectral density of the Dirac operator at the origin , a formula first derived in the QCD case by Smilga and Stern . The same analysis has recently been extended to the two other major chiral symmetry breaking classes by Toublan and Verbaarschot . There are thus also plenty of physical applications of PQChPT that have nothing to do with the artifacts of the quenched approximation at all.
While the supersymmetric approach to QChPT and PQChPT has been well tested, and is by now quite well understood, it is still of interest to find alternative means of formulating the same problem. In particular, the supersymmetry itself is not fundamental and not an inherent property of QChPT and PQChPT. Indeed, it has recently been shown in the context of the finite-volume effective chiral Lagrangian related to the soft part of the Dirac operator spectrum that the so-called replica method can provide a useful alternative technique . Here full or partial quenching is instead achieved by adding $`N_v`$ valence quarks (of usual statistics), and then taking the limit $`N_v0`$ at the end of the calculation. In ordinary QCD perturbation theory this procedure trivially kills all valence quark loops. In the framework of the effective Lagrangian of Goldstone bosons it is far from obvious that such a procedure can be carried out explicitly. It entails an extension of the chiral symmetry group U($`N`$) to non-integer $`N`$, and integrals over such a group are not known in closed form. Nevertheless, it turns out that in series expansions the required analytical continuation can be carried out explicitly , and results agree with what was earlier established by the supersymmetric method . This suggests that also conventional QChPT and PQChPT can be performed by simply taking the limit $`N_v0`$. In this paper we shall show that this is indeed the case. We shall give the very simple Feynman rules, and explain the intimate relationship to QChPT and PQChPT in the supersymmetric formulation. As a simple illustration we show how to derive the partially quenched chiral condensate to one-loop order using this replica method. This fully or partially quenched chiral condensate is a particularly convenient observable on which to test the non-perturbative finite-volume scaling results discussed above . The way partially quenched chiral perturbation theory smoothly matches on to this regime has been explained in ref. .
After providing the Feynman rules, it becomes quite obvious how the replica method in perturbation theory is equivalent to the supersymmetric method. We illustrate a few of the counting rules by considering a chiral $`k`$-point function below. Mainly out of curiosity, we also show how a variant of the replica method that is supersymmetric can be used to provide identical results. This supersymmetric variant is however slightly more cumbersome than the conventional $`N_v0`$ replica method, and we do not propose to use that particular variant for practical calculations.
## 2 The replica method
As explained above, with the replica method one adds $`N_v`$ valence quarks to the QCD Lagrangian, which here can be taken as any SU($`N_c3`$) gauge theory with $`N_f`$ physical (sea) quark flavors. Depending on the applications, it can be convenient to introduce $`k`$ sets of such valence quarks with $`k`$ different masses $`m_{v_j}`$, each set containing $`N_v`$ new quark flavors. The physical quark masses are denoted by $`m_f`$. The QCD partition function with these $`kN_v`$ additional quark species reads
$$𝒵^{(N_f+kN_v)}=[dA]\underset{j=1}{\overset{k}{}}det(i\text{/}Dm_{v_j})^{N_v}\underset{f=1}{\overset{N_f}{}}det(i\text{/}Dm_f)e^{S_{\mathrm{YM}}[A]}.$$
(1)
This partition function can be viewed as an unnormalized average of $`k`$ sets of $`N_v`$ identical replicas of the following partition functions of quarks in a fixed gauge field background $`A_\mu `$:
$$𝒵_{v_j}[d\overline{\psi }_jd\psi _j]\mathrm{exp}\left[d^4x\overline{\psi }_j(i\text{/}Dm_{v_j})\psi _j\right]$$
(2)
in the sense that
$$𝒵^{(N_f+kN_v)}=[dA]\underset{j=1}{\overset{k}{}}\left[𝒵_{v_j}\right]^{N_v}\underset{f=1}{\overset{N_f}{}}det(i\text{/}Dm_f)e^{S_{\mathrm{YM}}}.$$
(3)
Clearly, if we set $`N_v=0`$ this just reproduces the original QCD partition function. But the theory extended with $`kN_v`$ additional quark species in this way is a generating functional for partially quenched averages of $`\overline{\psi }_j\psi _j`$ and mixed averages also involving physical quark fields. One simply sets $`N_v`$ to zero after having performed the functional differentiations
$$\chi (m_{v_1},\mathrm{},m_{v_k},m_{f_1},\mathrm{},m_{f_l},\{m_f\})\underset{N_v0}{lim}\frac{1}{N_v^k}\frac{1}{N_f^l}\frac{}{m_{v_1}}\mathrm{}\frac{}{m_{v_k}}\frac{}{m_{f_1}}\mathrm{}\frac{}{m_{f_l}}\mathrm{ln}𝒵^{(N_f+N_v)}.$$
(4)
Technically, it can be convenient to add local sources for both scalar and pseudoscalar quark bilinears $`\overline{\psi }_j(x)\psi _j(x)`$ and $`\overline{\psi }_j(x)\gamma _5\psi _j(x)`$ and similarly for the vector and axial vector currents (for simplicity taken flavor diagonal). If needed, one can of course introduce corresponding sources in the physical quark sector. Because such terms have no bearing on our arguments presented below, we shall for simplicity omit them here.
### 2.1 Adapting the replica method to the chiral Lagrangian
For $`N_v`$ integer, and $`N_f+kN_v`$ small enough, chiral symmetry is assumed to be spontaneously broken according to the standard pattern of SU$`{}_{L}{}^{}(N_f+kN_v)\times `$SU$`{}_{R}{}^{}(N_f+kN_v)`$ SU($`N_f+kN_v`$). The effective low-energy theory can therefore be described in the entirely conventional framework of a chiral Lagrangian based on SU($`N_f+kN_v`$), with no new assumptions about the pattern of chiral symmetry breaking.<sup>1</sup><sup>1</sup>1The reader might worry about the assumption that $`N_f+kN_v`$ should be taken small enough for the theory to support spontaneous chiral symmetry breaking. Actually, there will be no new constraint from this. We simply analyze the chiral Lagrangian for arbitrary $`N_f+kN_v`$ even though this Lagrangian is only the low-energy theory of QCD for $`N_f+kN_v`$ sufficiently small. However, we take the limit $`N_v0`$ in the end. Then we must meet only the usual constraint that the number of physical light quarks $`N_f`$ should be small enough to lead to spontaneous chiral symmetry breaking. The cases $`N_f=0`$ and $`N_f=1`$ are obviously very special here. For $`N_f=1`$ there is not any spontaneous breaking of chiral symmetry in the theory after taking $`N_v`$ to zero, and the case $`N_f=0`$ (which would correspond to full quenching) is so unusual that we shall discuss it separately.
Having in mind a possibly non-trivial rôle played by the flavor singlet meson, the lowest-order effective chiral Lagrangian is taken to be the usual $`𝒪(p^2)`$ expression
$$=\frac{F^2}{4}\text{Tr}(_\mu U^\mu U^{})\frac{\mathrm{\Sigma }}{2}\text{Tr}(U+U^{})+\frac{\mu ^2}{2N_c}\mathrm{\Phi }_0^2+\frac{\alpha }{2N_c}_\mu \mathrm{\Phi }_0^\mu \mathrm{\Phi }_0.$$
(5)
Here the field $`U\mathrm{exp}[i\sqrt{2}\mathrm{\Phi }/F]`$ is an element of SU($`N_f+kN_v`$), and we have kept the flavor-singlet field $`\mathrm{\Phi }_0\text{Tr}\mathrm{\Phi }`$. As in the supersymmetric method , it proves convenient to work in a “quark basis” where $`\mathrm{\Phi }_{ij}`$ corresponds to $`\overline{\psi }_i\psi _j`$. With all external sources set to zero, this gives a simple propagator for the “off-diagonal” mesons corresponding to $`\mathrm{\Phi }_{ij}\overline{\psi }_i\psi _j,ij`$:
$$D_{ij}(p^2)=\frac{1}{p^2+M_{ij}^2},$$
(6)
while for the “diagonal” mesons $`\mathrm{\Phi }_{ii}\overline{\psi }_i\psi _i`$ the propagator can be written in the form
$$G_{ij}(p^2)=\frac{\delta _{ij}}{(p^2+M_{ii}^2)}\frac{(\mu ^2+\alpha p^2)/N_c}{(p^2+M_{ii}^2)(p^2+M_{jj}^2)(p^2)}.$$
(7)
Here $`M_{ij}^2(m_i+m_j)\mathrm{\Sigma }/F^2`$ and
$$(p^2)1+\frac{\mu ^2+\alpha p^2}{N_c}\left(\underset{j=1}{\overset{k}{}}\frac{N_v}{p^2+M_{v_jv_j}^2}+\underset{f=1}{\overset{N_f}{}}\frac{1}{p^2+M_{ff}^2}\right).$$
(8)
Note that $`N_v`$ enters as a parameter due to the mass degeneracy of the valence quarks in each of the $`k`$ sets. This is exactly what is required in order to apply the replica method. We remark that the unusual form of the propagator (7) just stems from using the quark basis and including the flavor singlet field $`\mathrm{\Phi }_0=\text{Tr}\mathrm{\Phi }`$, and not from any peculiarities of partial quenching. Although we borrow the result (7) from ref. , it is also unrelated to the supersymmetry of the method discussed there.
By including the $`\mathrm{\Phi }_0`$ field in the Lagrangian we have kept open the possibility of studing various expansion schemes (see $`e.g.`$ the second reference of ). The $`\mathrm{\Phi }_0`$ terms affect only $`G_{ij}`$. For $`G_{ii}`$ the flavor-singlet $`\mathrm{\Phi }_0`$ can give rise to double poles in the partially quenched limit, but the appearance of such double poles is not special to the replica method. Indeed such double poles are also present in the supersymmetric formulation where a thorough study has been done . As we prove in the next section the two formulations have equivalent perturbative expansions. The appearance of the double pole in the replica method is therefore completely analogous to the case of the supersymmetric formulation. In particular, we note that also in the replica formalism the case $`N_f=0`$ is quite special since in that case $`(p^2)`$ simply becomes unity, and the double pole in $`G_{ii}`$ is unavoidable. Moreover, in just that case there is no decoupling as the scale $`\mu `$ is sent to infinity.
In Table 1 we give the explicit relation between the Feynman rules based on the replica method, and those based on the supersymmetric formulation. The supersymmetry Feynman rules are supplemented by the standard relative minus sign between boson and fermion loops. Despite the additional minus signs in the Feynmann rules of the supersymmetric formulation, the Green functions are identical in the two formulations. As we show below, the signs due to combinatorics in the replica method match those arising from statistics and the supertrace in the supersymmetric formulation.
## 3 The equivalence between replica and supersymmetric PQChPT
In this section we formulate the equivalence between the generating functional of PQChPT in the replica and supersymmetric formulations. The equivalence proof is by default restricted to perturbation theory (expressed in terms of the Feynman rules), and we can in principle not make any statements at the non-perturbative level. But this is as it should be, as our whole framework in any case is restricted to chiral perturbation theory. The Lagrangian itself contains an infinitely long string of interactions that become relevant with increasing loop order, and we shall only demonstrate the equivalence at the one-loop level. However, seeing how the equivalence proof proceeds, it is pretty obvious how to generalize this to arbitrarily high order.
Our claim is: The generating functional of replica PQChPT for $`N_f+kN_v`$ flavors with $`k`$ sets of $`N_v`$ mass-degenerate quarks is in perturbation theory equivalent to the generating functional of supersymmetric PQChPT for $`N_f+k`$ fermionic and $`k`$ bosonic quarks.
By equivalence between the SU$`(N_f+kN_v)`$ and the SU$`(N_f+k|k)`$ generating functionals is meant that the chiral expansions are equivalent order by order. Of course, the respective limits, $`N_v0`$ and mass degeneracy between the $`k`$ bosons and $`k`$ of the fermions, are to be introduced at the end of the calculations. While we believe that this statement is true we will as mentioned above only address the equivalence at the one-loop level. At this one-loop level the contributions from the $`𝒪(p^4)`$ chiral Lagrangian act as counter terms and we can base the discussion on the Lagrangian of (5).
Let us first consider the sea sector. (The term sea sector is used when only sea quark masses are involved in differentiations of the generating functional.) For this sector both methods are equivalent to SU$`(N_f)`$ ChPT. In the replica formulation the contributions from the valence quarks at one-loop to any of the correlators
$$\chi (m_{f_1},\mathrm{},m_{f_l},\{m_f\})\frac{1}{N_f^k}\frac{}{m_{f_1}}\mathrm{}\frac{}{m_{f_k}}\mathrm{ln}𝒵^{(N_f+N_v)},$$
(9)
are necessarily proportional to positive powers of $`N_v`$. Hence the dependence on the valence quarks vanishes as $`N_v0`$, leaving the sea sector equivalent to SU$`(N_f)`$ ChPT. The analogous statement in supersymmetric PQChPT was proven in ref. . This equivalence was formulated as three theorems in that reference. At the risk of making some oversimplifications we state them compactly as follows:
I) The sea sector of SU$`(N_f+k|k)`$ PQChPT is equivalent to SU$`(N_f)`$ ChPT.
II) The super-$`\eta ^{}`$ is equivalent to the conventional $`\eta ^{}`$ of SU$`(N_f)`$ ChPT.
III) The double pole of $`G_{ii}`$ arise at a given fermionic quark mass if and only if all fermionic quarks with this mass are paired up by bosonic quarks.
In the supersymmetric formalism theorem I is established by noting that $`k`$ of the fermionic quarks and the $`k`$ bosonic quarks only appear as virtual loops in the sea sector. Since these 2$`k`$ quarks are paired up in masses the virtual loops cancel explicitly. This cancellation is also responsible for establishing theorem II in the supersymmetric formalism, only now it takes place in the quark loop corrections to the $`\eta ^{}`$-propagator. Finally theorem III follows directly from the structure of the last term in $`G_{ii}`$. We emphasize here that the obvious analogs of both theorems I and II are completely trivial in the present replica formalism. Theorem III, when re-stated in the language of the replica formalism, stipulate under what circumstances the potential double pole of $`G_{ii}`$ is cancelled: By inspection this occurs when $`M_{ii}=M_{ff}`$ for at at least one physical meson labelled by $`ff`$. The proof of theorem III is then almost identical in the replica and supersymmetric formulations. In the phrasing of refs. the double poles can only occur at mass scales that are completely quenched.
In the remaining quark sectors the equivalence is far less trivial. However, the supersymmetric bosonic Green functions equal the fermionic ones up to a well defined sign. So we can focus on the sectors involving fermionic valence quarks.
The equivalence in these sectors is not just of academic interest. As mentioned in the introduction, differentiations with respect to valence quark masses may be related to physical quantities. For instance the partially quenched chiral condensate for the valence quarks,
$$\mathrm{\Sigma }(m_v,\{m_f\})\underset{N_v0}{lim}\frac{1}{N_v}\frac{}{m_v}\mathrm{ln}𝒵^{(N_f+N_v)},$$
(10)
can be used to determine the Dirac spectral density. This density is given by the discontinuity of the partially quenched chiral condensate across a cut on the imaginary axis :
$$\rho (\lambda ;\{m_f\})=\frac{1}{2\pi }\mathrm{Disc}|_{m_v=i\lambda }\mathrm{\Sigma }(m_v,\{m_f\})=\frac{1}{2\pi }\underset{ϵ0}{lim}[\mathrm{\Sigma }(i\lambda +ϵ,\{m_f\})\mathrm{\Sigma }(i\lambda ϵ,\{m_f\})].$$
(11)
(This identification holds when one considers $`\mathrm{\Sigma }(m_v,\{m_f\})`$ as a function of a real mass $`m_v`$, and then replaces $`m_vi\lambda \pm ϵ`$.)
In the valence sector and the mixed sector the equivalence is established in two steps. First, notice that the propagator (7) of replica PQChPT for $`N_v=0`$ is identical to the one for the fermionic sector of the corresponding supersymmetric PQChPT in the limit where each of the boson masses is paired up with a fermion mass, see Table 1. (This equivalence holds trivially for the off-diagonal quark anti-quark propagators.) Second, the signs arising from combinatorics in the replica method is exactly matched by the opposite signs of boson and fermion loops occurring in the supersymmetric formulation.
In order to see exactly how the signs come to match in the two approaches, we explicitly give the derivation of the $`k`$-point function in the valence sector. The generalization to the mixed sector follows in complete analogy.
### 3.1 The one-point function in the valence sector
In this first example we give the contributions to the valence quark mass dependent chiral condensate defined in (10). We show how the cancellations that occur exactly match those of the supersymmetric formulation. It turns out that this simple 1-point function actually is ideally suited for illustrating the equivalence between the replica method and the supersymmetric method, as all essential properties of the propagators and of the combinatorics come into play.
To evaluate the one-point function we need to introduce just one set of replica fermions. Explicitly performing the differentiation of the generating functional, see Eq. (10), or alternatively counting the number of realizations of quark flow diagrams we have to one-loop
$$\frac{\mathrm{\Sigma }(m_v,\{m_f\})}{\mathrm{\Sigma }}=\underset{N_v0}{lim}\frac{1}{N_v}(N_v\frac{1}{F^2}\left(N_v\underset{f=1}{\overset{N_f}{}}\mathrm{\Delta }(M_{vf}^2)+N_v(N_v1)\mathrm{\Delta }(M_{vv}^2)+N_v\frac{1}{V}\underset{p}{}G_{vv}(p^2)\right))$$
(12)
where
$$\mathrm{\Delta }(M_{ij}^2)\frac{1}{V}\underset{p}{}\frac{1}{p^2+M_{ij}^2}\frac{1}{V}\underset{p}{}D_{ij}(p^2)$$
(13)
is a one-loop integral of the standard diagonal propagator for the off-diagonal mesons, $`\mathrm{\Phi }_{ij}\overline{\psi }_i\psi _j`$, $`ij`$. (We write everything in finite-volume notation, having also in mind applications of the kind discussed in refs. .) The first term in $`\frac{1}{V}_pG_{vv}(p^2)`$ is simply $`\mathrm{\Delta }(M_{vv}^2)`$. For arbitrary $`N_v`$ this term is seen to cancel against the term just before $`G_{vv}`$. In the $`N_v0`$ limit we also get rid of the term proportional to $`N_v`$, leaving simply
$$\frac{\mathrm{\Sigma }(m_v,\{m_f\})}{\mathrm{\Sigma }}=1\frac{1}{F^2}\left(\underset{f=1}{\overset{N_f}{}}\mathrm{\Delta }(M_{vf}^2)\frac{1}{V}\underset{p}{}\frac{(\mu ^2+\alpha p^2)/N_c}{(p^2+M_{vv}^2)(p^2+M_{vv}^2)(p^2)}\right).$$
(14)
This is completely analogous to the result obtained in the supersymmetric formulation. In that case a similar cancellation takes place between the first term in $`\frac{1}{V}_pG_{vv}(p)`$ and the loop of the meson built up by the fermionic and bosonic valance quark. It is also instructive to trace the cancellation of valence quark loops. In the supersymmetric formulation this cancellation occurs because of a matching boson loop, while in the present formulation it is due to the lack of a replica fermion. Pictorially speaking, this lack of a replica fermion acts like a boson.
### 3.2 The $`k`$-point function in the valence sector
As for the condensate, the $`k`$-fold derivative, $`k2`$, of $`\mathrm{ln}𝒵^{(N_f+kN_v)}`$ with respect to each of the valence quark masses is related to the spectral $`k`$-point function. The evaluation of the $`k`$-fold derivative is quite simple but we need to treat the case $`k=2`$ separately. The reason is simple: The product
$$\underset{j,k=1}{\overset{N_f+2N_v}{}}\mathrm{\Phi }_{ij}(x_1)\mathrm{\Phi }_{ji}(x_1)\mathrm{\Phi }_{lk}(x_2)\mathrm{\Phi }_{kl}(x_2),il$$
(15)
occurring in the two point function includes two connected terms, namely
$$\mathrm{\Phi }_{v_1v_1}(x_1)\mathrm{\Phi }_{v_2v_2}(x_2)\mathrm{\Phi }_{v_2v_2}(x_2)\mathrm{\Phi }_{v_1v_1}(x_1)$$
and
$$\mathrm{\Phi }_{v_1v_2}(x_1)\mathrm{\Phi }_{v_2v_1}(x_2)\mathrm{\Phi }_{v_1v_2}(x_2)\mathrm{\Phi }_{v_2v_1}(x_1).$$
For $`k>2`$ there is no connected analogue of the latter “crossed diagram”, since $`k`$ of the indices must be different (we differentiate with respect to different masses). The 2-point function is thus different from higher $`k`$-point functions because meson loops correspond to just quark-antiquark lines. In terms of the propagators the two-point function is<sup>2</sup><sup>2</sup>2This chiral 2-point function has been analyzed in the supersymmetric formulation by Osborn, Toublan, and Verbaarschot (private communication).
$$\frac{\chi (m_{v_1},m_{v_2},\{m_f\})}{\mathrm{\Sigma }^2}=\underset{N_v0}{lim}\frac{1}{N_v^2}\frac{1}{F^4}\left(N_v^2\frac{1}{V}\underset{p}{}D_{v_1v_2}(p^2)D_{v_2v_1}(p^2)+N_v^2\frac{1}{V}\underset{p}{}G_{v_1v_2}(p^2)G_{v_2v_1}(p^2)\right).$$
(16)
Whereas for $`k>2`$ there is no crossed diagram, and we are left with
$`{\displaystyle \frac{\chi (m_{v_1},\mathrm{},m_{v_k},\{m_f\})}{\mathrm{\Sigma }^k}}`$ $`=`$ $`\underset{N_v0}{lim}{\displaystyle \frac{1}{N_v^k}}(1)^k{\displaystyle \frac{1}{F^{2k}}}N_v^k{\displaystyle \frac{1}{V}}{\displaystyle \underset{p}{}}G_{v_1v_2}(p^2)\mathrm{}G_{v_kv_1}(p^2).`$ (17)
We observe that in both cases the $`N_v`$-dependence is such that the limit $`N_v0`$ becomes trivial. The corresponding expressions in the supersymmetric formalism are identical. Note that sea fermion and “ghost”-loops only appear in the one-point function.
## 4 From replicas to supersymmetry
Interestingly, in perturbation theory it is possible to use a peculiar variant of the replica method that is supersymmetric. This is because all $`N_v`$-dependence in the propagators and vertices is entirely parametric. We can thus make replicas of an arbitrary real number of valence quarks. Moreover, partial quenching can be achieved not only by taking $`N_v0`$, but also by taking $`N_v`$ to any fixed number of quarks $`N_v^{}`$, and re-interpreting the remaining $`N_f+N_v^{}`$ as physical quarks (of which it just happens that at least $`N_v^{}`$ are degenerate in mass). Because the $`N_v`$-dependence is parametric in perturbation theory, we can trivially go one step further and consider a partially quenched theory of $`N_f`$ physical fermions as the limit $`N_v\stackrel{~}{N}_v`$ of a theory based on $`N_f+\stackrel{~}{N}_v+N_v`$ quarks, out of which the $`\stackrel{~}{N}_v+N_v`$ quarks are degenerate in mass $`\stackrel{~}{m}_v=m_v`$. This corresponds to considering the effective theory of a fundamental partition function that is partially supersymmetric (for simplicity considering only one such set of replica quarks):
$`𝒵^{(N_f+\stackrel{~}{N}_v+N_v)}|_{N_v=\stackrel{~}{N}_v}`$ $`=`$ $`{\displaystyle [dA]det(i\text{/}Dm_v)^{N_v}\underset{f=1}{\overset{N_f+\stackrel{~}{N}_v}{}}det(i\text{/}Dm_f)e^{S_{\mathrm{YM}}[A]}}|_{N_v=\stackrel{~}{N}_v}`$ (18)
$`=`$ $`{\displaystyle [dA]\frac{det(i\text{/}D\stackrel{~}{m}_v)^{\stackrel{~}{N}_v}}{det(i\text{/}Dm_v)^{\stackrel{~}{N}_v}}\underset{f=1}{\overset{N_f}{}}det(i\text{/}Dm_f)e^{S_{\mathrm{YM}}[A]}}.`$ (19)
At this level the partition function is exactly as the starting point of the supersymmetric method. However, when we consider the effective partition function in terms of the Goldstone bosons, the working rules are entirely different. We keep our Feynman rules of Table 1, and just remember to take the limit $`N_v\stackrel{~}{N}_v`$ in the end. The fact that this procedure works is of course a direct consequence of the fact that in perturbation theory we can get bosons from fermions by letting the number of (degenerate) species go from positive to negative (also the “statistics” sign of closed fermion loops relative to closed boson loops comes out right in this way).
It is instructive to see how this supersymmetric variant of the replica method works in detail. Consider again our prototype of a Green function, that of the partially quenched chiral condensate. Using the notation of above, we find
$`{\displaystyle \frac{\mathrm{\Sigma }(m_v,\{m_f\})}{\mathrm{\Sigma }}}`$ $``$ $`\underset{\genfrac{}{}{0pt}{}{N_v\stackrel{~}{N}_v}{m_v\stackrel{~}{m}_v}}{lim}{\displaystyle \frac{}{m_v}}\mathrm{ln}𝒵^{(N_f+\stackrel{~}{N}_v+N_v)}`$ (21)
$`=`$ $`\underset{\genfrac{}{}{0pt}{}{N_v\stackrel{~}{N}_v}{m_v\stackrel{~}{m}_v}}{lim}{\displaystyle \frac{1}{N_v}}(N_v{\displaystyle \frac{1}{F^2}}(N_v[{\displaystyle \underset{f=1}{\overset{N_f}{}}}\mathrm{\Delta }(M_{vf}^2)+\stackrel{~}{N}_v\mathrm{\Delta }(M_{v\stackrel{~}{v}}^2)]+N_v(N_v1)\mathrm{\Delta }(M_{vv}^2)`$
$`+N_v{\displaystyle \frac{1}{V}}{\displaystyle \underset{p}{}}G_{vv}(p^2)))`$
where $`M_{v\stackrel{~}{v}}^2(m_v+\stackrel{~}{m}_v)\mathrm{\Sigma }/F^2`$, and $`G_{vv}(p^2)`$ is as in Table 1, except for the obvious change that now
$$(p^2)1+\frac{\mu ^2+\alpha p^2}{N_c}\left(\frac{N_v}{p^2+M_{vv}^2}+\frac{\stackrel{~}{N}_v}{p^2+M_{\stackrel{~}{v}\stackrel{~}{v}}^2}+\underset{f=1}{\overset{N_f}{}}\frac{1}{p^2+M_{ff}^2}\right).$$
(22)
Taking the degenerate mass limit $`\stackrel{~}{m}_v=m_v`$ and letting $`N_v\stackrel{~}{N}_v`$ we note that terms cancel out exactly as in the previous $`N_v0`$ replica method. For instance, in $`(p^2)`$ the terms linear in $`N_v`$ and $`\stackrel{~}{N}_v`$ just cancel each other. In eq. (21) the term proportional to $`N_v^2`$, which previously dropped out trivially in the $`N_v0`$ limit, is now precisely cancelled by a similar term proportional to $`N_v\stackrel{~}{N}_v\stackrel{~}{N}_v^2`$. All “unwanted” terms thus exactly cancel as they should, and we are left with the correct one-loop result (14). As we mentioned earlier, this example of the one-point function is actually the most instructive for illustrating the cancellations. The other $`k`$-point functions clearly proceed analogously.
Although it is thus possible to make a supersymmetric variant of the replica method, it is obviously rather pointless to do so. The simplest Feynman rules come from using just the conventional $`N_v0`$ limit. We also note that although the starting partition function (19) is identical to that forming the basis for the supersymmetric chiral Lagrangian , the effective theory one works with in the analogous supersymmetric replica scheme is of a very different nature, and has in fact here only been defined by means of the perturbative expansion.
## 5 Conclusions
We have shown how the replica method can been adapted to chiral perturbation theory. This provides a new and systematic realization of quenched and partially quenched chiral perturbation theory. We have demonstrated how the replica method is equivalent to the supersymmetric formulation in perturbation theory. This equivalence is quite trivial in the sector of physical quarks, and has allowed us to extend the three theorems of to the present replica formulation of PQChPT. The equivalence between the replica and the supersymmetric formalisms also extends outside the sea sector. The complete agreement (at least to one-loop order) of the two approaches offers a non-trivial consistency check. In particular, the assumed extension of the standard symmetry breaking pattern to the supergroup case is avoided in the present context. The fact that results agree can be taken as independent confirmation of the validity of both approaches.
As an equivalent but nevertheless independent formulation of PQChPT the replica method illustrates the fact that supersymmetry is a technical tool for quenching rather than of fundamental nature. For practical purposes the usefulness of the replica method as compared to the supersymmetric formulation is perhaps a matter of taste. The advantage of having fewer sign-rules using the replica method is to some extent traded for the marginally simpler combinatorics in the supersymmetric formulation.
Finally, the replica method presented here gives the background and the justification for the rules observed by Colangelo and Pallante in . Within the supersymmetric formulation they studied fully quenched chiral perturbation theory to one loop. Based on an explicit calculation of the divergent parts of the generating functional for both SU$`(k|k)`$ (and the additional U(1) of the $`\mathrm{\Phi }_0`$) and standard SU$`(N_f)`$ chiral perturbation theory (without the $`\mathrm{\Phi }_0`$), they proposed a set of rules for writing down large parts of the SU$`(k|k)`$ generating functional from that of SU$`(N_v)`$. The equivalence between the SU$`(k|k)`$ and SU$`(N_v0)`$ theories (when the $`\mathrm{\Phi }_0`$ is included in both), is a special case of the general equivalence established here. This formally establishes the rules suggested in and furthermore shows that the terms missing in SU$`(N_f0)`$ chiral perturbation are just those produced by including the $`\mathrm{\Phi }_0`$. The procedure to compute in partially quenched chiral perturbation theory to any order is now extremely simple. One must take a usual chiral SU($`N_f+N_v`$) chiral Lagrangian and add the contributions from $`\mathrm{\Phi }_0`$. For example, to order $`p^6`$ the whole list of divergent contributions in the case of a degenerate SU($`N_v`$) theory is provided in ref. . This can form the basis for a fully quenched calculation once the contributions from the flavor singlet have been included (for a discussion of the large-$`N_c`$ limit, see $`e.g.`$ ref. ).
Acknowledgement:
This work was supported in part by EU TMR grant no. ERBFMRXCT97-0122.
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# Inhomogeneity-Induced Superconductivity?
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## Abstract
A $`t`$-$`J`$-like model for inhomogeneous superconductivity of cuprate oxides is presented, in which local anisotropic magnetic terms are essential. We show that this model predicts pairing, consistent with experiments, and argue how the macroscopic phase-coherent state gradually grows upon lowering of the temperature. We show that appropriate inhomogeneities are essential in order to have significant pair binding in the thermodynamic limit. Particularly, local breaking of $`SU(2)`$ spin symmetry is an efficient mechanism for inducing pairing of two holes, as well as explaining the magnetic scattering properties. We also discuss the connection of the resulting inhomogeneity-induced superconductivity to recent experimental evidence for a linear relation between magnetic incommensurability and the superconducting transition temperature, as a function of doping.
There is a growing body of experimental evidence suggesting that the superconducting state in cuprate oxides is “inhomogeneous,” such that the locally defined charge density varies across the sample in the ground state. Spatially inhomogeneous features in the spin and charge channels have been indicated in a number of experiments on high-$`T_c`$ materials . The simplest realization of this state is the so-called “stripe” phase where charges cluster in nanoscale linear patterns and the remainder of the sample is essentially an antiferromagnetically correlated insulator. This represents a nanoscale distribution of charge and spin, rather than a global phase separation. These experiments lead us to a central question: Is the superconducting state found in high-$`T_c`$ cuprates inhomogeneous as a result of spin/charge inhomogeneities? We believe that the answer to this question is yes. Moreover, we argue that spatial spin/charge inhomogeneities are in fact necessary for pairing and subsequent formation of the superconducting state in these compounds. This situation should be contrasted with the case of conventional superconductors, resolved by the BCS (Bardeen-Cooper-Schrieffer) theory, that starts with a homogeneous metallic state and describes the formation of a homogeneous superconducting state. It is commonly believed that magnetic correlations, characterized by the spin exchange energy $`J1500K`$ are responsible for the pairing interactions in the cuprates and are, therefore, crucial for our inhomogeneous exchange approach. Moreover, the existence of a spin gap has been experimentally proven . A model that naturally incorporates these features (inhomogeneities, magnetism and spin gap) is a $`t`$-$`J`$-like model with explicitly broken spatial and magnetic symmetries.
Of central importance for the present work is a microscopic model which captures the main low energy physics of doped antiferromagnetic (AF) Mott insulators. In particular, we show that our minimal model properly describes the magnetic properties observed in a wide variety of doped cuprate oxide materials. The key ingredient is the existence of magnetic perturbations which explicitly break local spin-rotational invariance (e.g., due to local spin-orbit coupling ) and thereby induce substantial hole pair binding. We then develop a mean-field theory of superconductivity based upon a phenomenology from our microscopic model. We emphasize that in our approach there are two, in principle different, energy scales; one associated to the pairing of holes and another related to the phase coherence of the pairs (that establishes $`T_c`$). Basically the inhomogeneities induce a strong hole pairing, which in turn Josephson-tunnel coherently between stripes, separated by insulating AF regions, phase-locking into a macroscopic supercurrent superfluid stiffness. Recently, a simple linear relation between the superconducting transition temperature $`T_c`$ and the AF incommensuration $`\delta `$ has been observed for the LSCO and YBCO high-$`T_c`$ compounds: $`k_BT_c\delta `$, where $`k_B`$ is the Boltzman constant. We find that for this relation to hold we need a power law Josephson tunneling.
Microscopic Model. Our model Hamiltonian describing the low energy dynamics of CuO<sub>2</sub> planes is $`H=H_{tJ}+H_{\mathrm{inh}}`$, where the background Hamiltonian $`H_{tJ}`$ is the standard $`t`$-$`J`$ model
$$H_{tJ}=t\underset{i,j,\sigma }{}c_{i\sigma }^{}c_{j\sigma }^{}+J\underset{i,j}{}(𝐒_i𝐒_j\frac{1}{4}\overline{n}_i\overline{n}_j).$$
(1)
For the inhomogeneous component, we take
$$H_{\mathrm{inh}}=\underset{\alpha ,\beta }{}\delta J_zS_\alpha ^zS_\beta ^z+\frac{\delta J_{}}{2}\left(S_\alpha ^+S_\beta ^{}+S_\alpha ^{}S_\beta ^+\right),$$
(2)
with $`\delta J_{}\delta J_z`$, representing the magnetic perturbation of a static local Ising anisotropy, locally lowering spin symmetry. Here $`i,j`$ are near-neighbor sites, while $`\alpha ,\beta `$ are two near-neighbor sites characterizing the bonds that are perturbed and where $`SU(2)`$ spin-rotational invariance is explicitly broken. The network of perturbed bonds form mesoscopic patterns determined by the distribution of stripe segments. The spin-$`\frac{1}{2}`$ operator $`𝐒_i=\frac{1}{2}c_{i\sigma }^{}𝝉_{\sigma \sigma ^{}}^{}c_{i\sigma ^{}}^{}`$, the electron occupation number $`\overline{n}_i=c_i^{}c_i^{}+c_i^{}c_i^{}`$, and $`c_{i\sigma }^{}(c_{i\sigma }^{})`$ creates(annihilates) an electron of spin $`\sigma `$ in a Wannier orbital centered at site $`i`$; $`𝝉`$ are the 3 Pauli matrices. This is a three-state model with the hopping constrained to the subspace with no doubly occupied sites. In the following, all energies will be measured in units of $`J`$.
Our modeling strategy consists in assuming the existence of an inhomogeneous mesoscopic skeleton of stripe segments, and then exploring its consequences, mainly the competition between magnetism and superconductivity. We do not address here the important problem of the formation and stability of this skeleton morphology. The origin(s) of “stripe segment” formation in high temperature superconductors is as yet unclear and several physical mechanisms could act cooperatively and be responsible for the generation of multiple length scales, among them: spin-orbit coupling, local Jahn-Teller distortions induced by the hole, effective interactions coming from a multi-band Hubbard Model (HM) (including explicitly the oxygen and copper bands), oxygen buckling at the stripe, and other local magnetoelastic effects . Competitions between attractive short range forces and repulsive long range ones can certainly spontaneously break translational and/or rotational invariance in the CuO<sub>2</sub> planes , but this is not necessarily the only mechanism. However, we show below that the mere existence of appropriate local magnetic anisotropies is crucial for pair-formation.
We start by showing that, as far as we could numerically determine, only by including a local Ising perturbation such as $`H_{\mathrm{inh}}`$ in Eq. 1 can a strong pairing of holes be obtained (see Fig. 1). All the calculations were made using exact diagonalization in small clusters with periodic boundary conditions in all spatial directions. We studied one-dimensional (1D) chains up to 16 sites and $`8\times 2`$ clusters. Hereafter, we will view our clusters as simulating systems where the longer direction is perpendicular to the stripes. We investigated the system size scaling for the binding energy of two holes defined as $`E_b=(E_{2\mathrm{holes}}E_{0\mathrm{hole}})2(E_{1\mathrm{hole}}E_{0\mathrm{hole}})`$ for several models in 1D and 2D. Although for small enough systems the binding energies could be very large, they all seem to extrapolate towards no (or extremely small) binding in the thermodynamic limit, with the clear exception of the inhomogeneous $`t`$-$`J`$$`J_z`$ case. We have also studied several one-band HMs, but we could again not find definite binding. The $`t`$-$`J`$$`J_z`$ model, the only one unambiguously giving binding in the thermodynamic limit, is obtained by breaking spin-rotation symmetry in $`d`$ near-neighbor bonds $`\alpha ,\beta `$, repeated with period $`P`$, by an amount $`\delta J_z=0`$, $`\delta J_{}<0`$ in Eq. 1. This $`t`$-$`J`$$`J_z`$ model is a most natural way to induce a spin-gap. We have checked that the spin-gap is present for our $`t`$-$`J`$$`J_z`$ model. The inhomogeneities forming the superstructure, which we impose by hand in the Hamiltonian, we term stripes. In Fig. 2 we show the hole correlation function $`g|n_0.n_i|g`$, where $`|g`$ is the ground state of the system ($`g|g=1`$). This correlation function gives information about the structure of the pair. It can be seen that as the hopping strength $`t`$ is increased beyond a characteristic value the second hole jumps from one stripe to the neighboring one, starting from an initial configuration where both holes are in the same stripe for small $`t`$. This can be understood as a result of a length(time)-scale competition: the pair size exceeds the stripe width.
To explore the nature of the binding, we have examined the canonical transformation of a $`t`$-$`J`$ model from a one-band HM and traced what kind of perturbations would produce a $`t`$-$`J`$$`J_z`$ term. This corresponds to a term like $`V_\sigma \overline{n}_{i,\sigma }\overline{n}_{j,\overline{\sigma }}`$ ($`i,j`$ first neighbors) in an extended HM, which may in turn arise, for example, from local magnetoelastic (e.g. oxygen (un)buckling) or spin-orbit couplings. Note, again, that here this is a perturbation only at the stripes. This kind of anisotropy manifests itself in two different ways in the $`t`$-$`J`$$`J_z`$ model, enhancing both the $`\frac{1}{4}\overline{n}_i\overline{n}_j`$ and the easy axis (Ising) terms of the Hamiltonian. The first one is an explicit pairing term for electrons. To see the relative importance of each term we have calculated the binding energy of a Hamiltonian like (1) but excluding the $`\frac{1}{4}\overline{n}_i\overline{n}_j`$ term and including bonds with broken spin-rotational symmetry. This model corresponds to holes (with no spin) propagating in an antiferromagnet, but not derived from a canonical transformation of a one-band HM. We find that it still has binding, as should be expected. Thus, the easy axis exchange term is partially responsible for the binding energy.
In order to understand this exchanged-based pairing mechanism, it is useful to explore some limiting cases. When the magnetic energy scales are the most relevant ones: $`(J_z=J+\delta J_z,J_{}=J+\delta J_{})t`$, it is easy to realize that, depending upon $`J_z,J_{}`$ being smaller or larger than $`J`$, the holes will prefer to be in the stripes or between stripes (with no binding), respectively. The opposite limit, i.e. purely kinetic energy, leads to delocalized holes and no binding. The situation where $`J_z<J`$ and $`t`$ is relevant corresponds to the intermediate regime where pairing is observed. Notice that pairing of holes does not necessarily imply that holes should share the same stripe, they can occupy neighboring ones (see Fig. 2, upper panel), thus avoiding phase separation. Details of the charge confinement and pairing potentials from the (dynamic) spin-field profiles in the superlattice skeletons will be given elsewhere.
Having demonstrated a minimal model for hole binding, we have computed spin correlation functions in clusters of size $`N_x\times N_y=N`$ ($`N_x=8,N_y=2`$). Here, we simulate the stripes by including an anisotropic $`\delta J_{}<0`$ in one $`y`$-bond with $`P=4`$; the rest of the bonds, including all the $`x`$-bonds, were not changed from the background $`t`$-$`J`$ model (see inset, Fig. 3). We cannot perform scaling on this size of inhomogeneous system, but the binding energy is still considerable. We have included up to 6 holes. In the case of four holes (the one more relevant to the stripes in the underdoped regime for cuprate oxides) and small $`t(J)`$ the holes bind in pairs on each site of the inhomogeneous bond (see Fig. 2, lower panel).
In Fig. 3 we show the spin-structure factor function $`S(𝕜=(k_x,k_y))`$ defined as: $`S_𝕜S_𝕜=(1/N)_{𝕚,𝕛}\mathrm{exp}(i𝕜.𝕣_𝕚)g|𝕊_𝕛.𝕊_{𝕛+𝕚}|g`$. This function corresponds to the observable in the elastic neutron scattering experiments. For $`t`$ small, only one peak occurs with $`k_y=\pi `$, corresponding to two essentially uncorrelated AF domains, isolated from each other by the pinned hole wall. As $`t`$ increases ($`t/J2`$, near the accepted set of values of the 2D $`t`$-$`J`$ model for cuprates ), the holes gain kinetic energy by visiting the first neighbor sites around the anisotropy region, but still bind together. The effective width of the pair thus increases to two sites. Magnetic energy is then gained if the two domains shift their staggered magnetization by $`\pi `$. We suggest that these $`𝒪(t^2)`$ processes are responsible for the incommensuration ($`\delta `$) in $`S(𝕜)`$ observed in the experiments. This $`\delta `$ is the inverse of twice the period $`P`$ of the stripes. In this picture the incommensuration is a consequence of the holes and their kinetic energy, and is a property of the ground state. Basically, it results from the competition between hole delocalization and magnetic fluctuations. This contrasts with some other proposed explanations, where $`\delta `$ is a magnetic thermodynamic property . It is interesting to note that this incommensuration arises even in the homogeneous $`t`$-$`J`$ model although for different values of $`t`$. This suggests that the experimentally observed magnetic properties are already present in a homogeneous $`t`$-$`J`$ model, but in order to obtain binding of holes appropiate inhomogeneous terms must be included.
When more than four holes are added to the system, but only two bonds are perturbed, $`S(𝕜)`$ changes qualitatively. Instead of showing an incommensurability around $`𝕜=𝐐=(\pi ,\pi )`$, it has a broad peak at $`𝕜=(0,\pi )`$. In this case the extra holes are delocalized in the middle of the AF space between stripes. This suggests that when the stripes reach their minimum separation, extra holes are responsible for the experimental increase and ultimate disappearance of the incommensurability.
Model of Josephson Spaghetti. It is important to relate the above discussion to the experimental evidence for the incommensurate neutron scattering peak, seen in LSCO (e.g., ) and YBCO compounds . In both of these cases a simple linear relation between $`T_c`$ and the peak incommensuration $`\delta `$ near $`𝐐`$ (or peak width in YBCO) is obeyed . Namely,
$`k_BT_c=\mathrm{}v^{}\delta .`$ (3)
The anomalously low velocity values for $`v^{}`$ depend on the compound . These velocities are independent of the carrier concentration and the only doping ($`x`$) dependence entering Eq. 3 is through $`\delta (x)`$.
An interpretation of this relation is to connect possible superconductivity mechanisms to the existence of the fluctuating stripes. Here we focus on the simple proportionality between $`T_c(x)`$ and a doping dependent length $`\mathrm{}(x)`$, determined from the neutron scattering: $`T_c(x)1/\mathrm{}(x),\mathrm{}(x)=1/\delta (x)`$. We consider how the Josephson tunneling of pairs between stripe segments can produce the relation between the phase ordering transition temperature $`T_c`$ and the typical length $`\mathrm{}(x)`$. The stripe-stripe distance $`r`$ is a random quantity due to intrinsic mechanisms as well as disorder and/or crystal imperfections . Therefore, we will assume that the mean-field transition temperature depends upon the Josephson coupling $`J(r)`$, averaged with some probability distribution of stripe separations.
Our model Hamiltonian of random stripe separation and associated inter- and intra-stripe random Josephson coupling (see Fig. 4) is
$`={\displaystyle \underset{ij}{}}J_{ij}\mathrm{exp}[i(\varphi _i\varphi _j)],`$ (4)
$`J_{ij}=J(r_{ij})=t_0/r_{ij}^\alpha ,`$ (5)
where the summation is taken over the coarse-grained regions $`i=1,\mathrm{},𝒩`$ with well-defined phases, labeled $`\varphi _i=\varphi (r_i)`$ and $`J_{ij}`$ becomes zero eventually at large distances. Next, we will assume some probability distribution $`P(r)`$ for the stripe-stripe distance. For simplicity we will take the “box” distribution $`P(r)`$ centered around $`\mathrm{}=1/\delta `$ and with finite width $`a=\nu \mathrm{}`$, where $`\nu =𝒪(1)`$ is a parameter. $`P(r)=C`$, for $`\mathrm{}ar\mathrm{}+a`$, and zero otherwise. Here we have simplified to one length scale for both $`a`$ and $`\mathrm{}`$. The normalization constant in 2D is $`C=[4\pi \mathrm{}a]^1`$. In this model one easily finds
$`J(r)`$ $`=`$ $`{\displaystyle d^2rP(r)J(r)}={\displaystyle \frac{2\pi t_0C}{2\alpha }}a_1\mathrm{}^{2\alpha },`$ (6)
$`r`$ $`=`$ $`{\displaystyle \frac{2\pi C}{3}}a_2\mathrm{}^3,`$ (7)
where the constants $`a_1,a_2`$ are $`𝒪(1)`$. Thus, for $`\alpha =1`$, we obtain the experimentally observed relation
$`T_c(x)J(r)[r]^1=\delta (x).`$ (8)
We have examined a variety of distributions $`P(r)`$ and functional dependences for $`J_{ij}`$; Eq. 5 with $`\alpha =1`$ is the only one reproducing the experimental data (at our mean-field random Josephson coupling level. Implicit in $`J(r)`$ is the exponential cutoff at lengths much larger than the stripe-stripe distance. This cutoff is necessary to have a well defined thermodynamic limit but is not important for short length scales). The screening mechanism (magnetic, elastic fluctuations, etc.) responsible for this form requires detailed microscopic modeling . The present model does not allow us to determine the magnitude of $`v^{}`$ without making specific assumptions about parameters such as $`t_0`$.
In conclusion, we have presented a microscopic model that captures the essential magnetic and pairing properties of high-temperature cuprate superconductors. Pairing of holes is a consequence of the existence of an AF background. (Analogous scenarios in other broken symmetry backgrounds, e.g. doped charge-density-wave bismuthates, are likely.) Crucially, however, the glue is provided by magnetic inhomogeneities whose precise origin remains to be unraveled, although it seems fundamental that these perturbations should locally break spin-rotational invariance. This pairing mechanism is kinetic exchange-interaction based and involves a competition between Ising and $`XY`$ symmetries. We emphasize that the pair-binding occurs only for intermediate strengths of $`t`$ and (local) $`J_z`$. We also introduced a phenomenological model and scenario for the macroscopic superconductivity based upon coherent Josephson-tunneling of pairs of holes between these magnetic inhomogeneities in a mesoscopic liquid-crystal-like skeleton. We have shown that this approach is able to recover the magnetic incommensuration $`\delta `$ and its experimentally observed relation to $`T_c(x)`$. Finally, we note that we have assumed static magnetic inhomogeneities. The case where the broken spin-symmetry follows the hole is also interesting. Elsewhere, we will discuss this generalization of coupling the inhomogeneity selfconsistently to dynamic holes.
Work at Los Alamos is sponsored by the US DOE under contract W-7405-ENG-36.
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# 1 Predicted ratios 𝑅=Γ(𝑋^∗→𝑋𝑒⁺𝑒⁻)/Γ(𝑋^∗→𝑋𝛾), for some mesonic systems 𝑋. The branching ratios ℬ(𝑋^∗→𝑋𝛾) and mass differences 𝑚_𝑋^∗-𝑚_𝑋 are from []. We have assumed 𝑚_𝐵_𝑠^∗-𝑚_𝐵_𝑠=𝑚_𝐵_𝑑^∗-𝑚_𝐵_𝑑. The number for the 𝐾^∗(892) decay is included for completeness only, since here form factors might become important, and of course, the 𝐾^∗(892) is a 1⁻–particle. Note however that 𝐾^∗(892)→𝐾𝑒⁺𝑒⁻ has not been observed. The only ratio for a heavy system 𝑋 observed so far is 𝐵^∗ [], where 𝑅=(4.7±1.1±0.9)×10⁻³.
Determining the quantum numbers of excited heavy mesons
G. Eilam and F. Krauss
Technion–Israel Institute of Technology
32000 Haifa, Israel
Abstract
We discuss the decays $`X^{}Xe^+e^{}`$ (“Dalitz decays”) of excited heavy mesons into their ground states and an electron–positron pair. We argue that the measurement of the invariant mass spectrum of the lepton pair gives clear indication on the quantum numbers of the excited meson and thus provides an experimental test of the quark model predictions.
PACS numbers: 14.40.Lb, 14.40.Nd, 12.39.Jh, 13.40Hq, 13.20Fc, 13.20.He, 13.25.Ft, 13.25.Hw
We investigate radiative decays of excited heavy mesons with charm and beauty, i.e. $`D^0`$, $`D_s^+`$, $`B^0`$ and $`B_s^0`$, into their ground state plus an electron–positron pair. There are several reasons that led us to investigate these decays:
1. The quantum numbers $`J`$ and $`P`$ of the excited mesons given in the Review of Particle Physics are either merely predictions of the quark model, which gives for the particles discussed here $`J^P=1^{}`$, or at best the quantum numbers need confirmation <sup>1</sup><sup>1</sup>1We would like to mention here, that for the $`B`$–meson even the quantum numbers of the ground state are not known experimentally; nevertheless we assume $`J^P(B)=0^{}`$..
2. The quantum numbers of some even higher heavy meson states, depend on the correct assignments for the $`X^{}`$’s which are the focus of the present work.
3. The central values of the branching ratios for $`D^0D^0\gamma `$ and $`D^0D^0\pi `$ sum up to exactly 100 $`\%`$ . Only closer inspection of both relevant experimental papers reveals that this was one of the assumptions underlying the analyses. A similar assumption was made for $`D_s^\pm `$, where again it was assumed that branching ratios for $`D_s^+D_s^+\gamma `$ and $`D_s^+D_s^+\pi ^0`$ sum to 100 $`\%`$ . Although the current errors on the individual branching ratios are clearly higher than the branching ratios for the decay into a Dalitz pair which, as we will see below, are of the order of $`0.5\%`$ of their real photon emission counterparts, in the long term this channel cannot be neglected. Especially in light of future precision measurements in the $`B`$–sector, a precise tracking of all outgoing $`D`$’s is required. In other words, it is important to know which decay of $`X^{}`$ has the closest branching ratio to the decays that have been already measured. Let us remark, that for the $`B^{}`$ decay, for which the dominant mode is $`B^{}B\gamma `$, its decay into $`Be^+e^{}`$ presented here has a much larger rate (about $`0.5\%`$, see below) than the decay $`B^{}B\gamma \gamma `$, recently discussed in .
4. The only recent theoretical analysis of an $`X^{}Xe^+e^{}`$ decay was performed for $`D^{}De^+e^{}`$ and yielded $`R0.001`$, where the ratio $`R`$ is defined as
$`R={\displaystyle \frac{\mathrm{\Gamma }(X^{}Xe^+e^{})}{\mathrm{\Gamma }(X^{}X\gamma )}}={\displaystyle \frac{\mathrm{\Gamma }_{ee}}{\mathrm{\Gamma }_\gamma }}`$ (1)
with $`X=D`$. As found below, we predict $`R`$ to be about 5 times larger than the result of . Let us note here that the suggestion to use the ratio between the Dalitz decay and the real photon emission–and in particular its $`q^2m_{12}^2`$ dependence–to determine the quantum numbers of the decaying particle, was made many years ago . Our results are consistent with theirs.
Since we will be mainly concerned with ratios $`R`$, where $`X=D,B`$, it is sufficient to parametrize the $`X^{}X\gamma `$ transition with some effective coupling constant $`g_{X^{}X\gamma }`$ and a suitable form factor $``$. We replace the various form factors possible for off–shell photons by a single one, and furthermore assume that it is independent of $`q^2`$. This is justified by the observation that for all the $`X^{}`$$`X`$ combinations we consider, the mass difference $`\mathrm{\Delta }_{X^{}X}m_X^{}m_X`$ is of the order of up to $`150`$ MeV and thus much smaller than the $`\rho `$–mass which is the most relevant one in the vector dominance model for the form factors. Therefore the influence of the form factors on the results is indeed very small. We have confirmed that our numerical results are practically unaffected by the assumption $`(q^2)(0)`$.
The matrix elements we consider read
$`_{1^{}0^{}\gamma }`$ $`=`$ $`g_{X^{}X\gamma }(q^2)ϵ^{\alpha \beta \mu \nu }ϵ_\alpha (\gamma )ϵ_\beta ^{}(X^{})P_\mu (X^{})q_\nu (\gamma )`$
$`_{2^+0^{}\gamma }`$ $`=`$ $`g_{X^{}X\gamma }(q^2)ϵ^{\alpha \beta \mu \nu }ϵ_\alpha (\gamma )P_\beta (X^{})ϵ_{\mu \rho }^{}(X^{})q_\nu (\gamma )q^\rho (\gamma ).`$ (2)
We have concentrated here on the $`1^{}`$ and $`2^+`$ quantum numbers for the excited and $`0^{}`$ for the ground state, since in the $`D`$–system the quantum numbers of the ground states are known to be $`0^{}`$ and because the $`D^{}`$ decays into a $`D`$ and both a pion or a photon. Therefore, the $`1^{}`$ and $`2^+`$ are the lowest lying quantum numbers allowed for the $`D^{}`$. In fact, for the $`B`$–system the situation is slightly different. Here the small mass difference of the $`B^{}`$ and the $`B`$ of roughly 45 MeV does not provide enough phase space for the $`B^{}`$ to decay into a $`B`$ and a $`\pi `$, hence the $`B^{}`$ could in principle be a $`1^+`$. This would cause the relevant strong coupling constant $`g_{B^{}B\pi }`$ to vanish. However, we consider this idea to be too far–out and will not discuss it here.
We use the completeness relations (for the $`2^+`$ see e.g: )
$`{\displaystyle ϵ^\mu (p)ϵ^\nu (p)}`$ $`=`$ $`\left(g^{\mu \nu }{\displaystyle \frac{p^\mu p^\nu }{p^2}}\right)`$
$`{\displaystyle ϵ^{\mu \nu }(p)ϵ^{\rho \sigma }(p)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}[(g^{\mu \rho }{\displaystyle \frac{p^\mu p^\rho }{p^2}})(g^{\sigma \nu }{\displaystyle \frac{p^\sigma p^\nu }{p^2}})`$ (3)
$`+\left(g^{\sigma \mu }{\displaystyle \frac{p^\sigma p^\mu }{p^2}}\right)\left(g^{\nu \rho }{\displaystyle \frac{p^\nu p^\rho }{p^2}}\right)`$
$`{\displaystyle \frac{2}{3}}(g^{\mu \nu }{\displaystyle \frac{p^\mu p^\nu }{p^2}})(g^{\rho \sigma }{\displaystyle \frac{p^\rho p^\sigma }{p^2}})].`$
After squaring, summing and averaging we obtain for the the decay into a real photon
$`3{\displaystyle \overline{\left|_{1^{}0^{}\gamma }\right|^2}}`$ $`=`$ $`2g_{X^{}X\gamma }^2^2(0)(Pq)^2`$
$`5{\displaystyle \overline{\left|_{2^+0^{}\gamma }\right|^2}}`$ $`=`$ $`g_{X^{}X\gamma }^2^2(0){\displaystyle \frac{(Pq)^4}{P^2}},`$ (4)
where $`P`$ and $`q`$ are the momenta of the excited meson and the photon, respectively. The resulting branching ratios are given by
$`\mathrm{\Gamma }_{1^{}0^{}\gamma }`$ $`=`$ $`g_{X^{}X\gamma }^2^2(0){\displaystyle \frac{(m_X^{}^2m_X^2)^3}{96\pi m_X^{}^3}}`$
$`\mathrm{\Gamma }_{2^+0^{}\gamma }`$ $`=`$ $`g_{X^{}X\gamma }^2^2(0){\displaystyle \frac{(m_X^{}^2m_X^2)^5}{1280\pi m_X^{}^5}}.`$ (5)
With suitable replacements we recover the known result for the width of the decay $`a_2\pi \gamma `$ .
For the decay into $`e^+e^{}`$ the polarization vector $`ϵ_\mu `$ of the photon has to be replaced by the lepton–current $`e\overline{u}(e^{})\gamma _\mu u(e^+)`$. Squaring the matrix elements and summing and averaging over polarizations yields
$`3{\displaystyle }\overline{\left|_{1^{}0^{}ee}\right|^2}=2g_{X^{}X\gamma }^2{\displaystyle \frac{^2(q^2)}{q^4}}`$ (6)
$`\left[4m_e^2\left((Pq)^2P^2q^2\right)+q^2\left(2(Pp_e)^2+2(Pp_{\overline{e}})^2P^2q^2\right)\right]`$
$`5{\displaystyle }\overline{\left|_{2^+0^{}ee}\right|^2}=g_{X^{}X\gamma }^2{\displaystyle \frac{^2(q^2)}{q^4}}{\displaystyle \frac{(Pq)^2P^2q^2}{P^2}}`$ (7)
$`\left[4m_e^2\left((Pq)^2P^2q^2\right)+q^2\left(2(Pp_e)^2+2(Pp_{\overline{e}})^2P^2q^2\right)\right].`$
The resulting widths $`\mathrm{\Gamma }(X^{}Xe^+e^{})`$ and their respective ratios to the real photon widths agree with the results given in . Our results for the ratio $`R`$, defined in Eq. 1, in some meson systems are displayed in Table 1. All of them are of the order of $`510^3`$.
As can be observed from the table, the ratio $`R`$ may serve as an indicator for the $`J^P`$ quantum numbers of $`X^{}`$ mesons, which are believed to have $`J^P=1^{}`$. At present, the only ratio for which an experimental number exists is $`R=(4.7\pm 1.1\pm 0.9)\times 10^3`$, for $`B^{}`$ decay . Although the central value agrees well with the quark model quantum numbers $`J^P(B^{})=1^{}`$, it is premature to claim, in view of the large error, a clear–cut rejection of the $`2^+`$ possibility.
A better indicator for the quantum numbers of $`X^{}`$ is the distribution of the invariant mass squared ($`q^2=m_{12}^2`$) of the $`e^+e^{}`$ pair. To illustrate the effect of different quantum numbers on $`d\mathrm{\Gamma }(X^0X^0e^+e^{})/dq^2`$, scaled by $`\mathrm{\Gamma }(B^0B^0\gamma )`$, we display in Fig. 1 the results for the $`B`$–system. Similar results are obtained for $`X=D`$. Clearly, the change in the quantum numbers affects the tail of the $`q^2`$–distribution of the electron–positron pairs by factors about 2 and larger. Modifications of $`q^2`$ distribution depicted in Fig. 1 by Vector Meson Dominance form factors would be practically invisible. This provides an excellent tool to determine the quantum numbers of $`X`$.
To summarize, we have presented results for the decay widths of $`D^{}De^+e^{}`$, $`D_s^\pm D_s^\pm e^+e^{}`$ and $`B^0B^0e^+e^{}`$. All of them are of the order of roughly $`0.5\%`$ of the corresponding decays into real photons, as expected from a QED–like calculation. The difference between the two $`J^P`$ assignments $`1^{}`$ and $`2^+`$ for the the decaying particle, is about $`5\%`$. A significant difference arises for the $`q^2`$–distribution, especially at high $`q^2`$. In the light of the recent experimental results by on the invariant mass spectrum in the Dalitz decay of $`B^{}`$’s we are very optimistic that the ambitious measurements we advertise are feasible.
Acknowledgments
We thank H. Landsman, Y. Rozen and P. Singer for helpful discussions. The research of G.E. was supported in part by the Israel Science Foundation founded by the Israel Academy of Sciences and Humanities and by the Harry Werksman Research Fund. F.K. would like to acknowledge financial support by the Minerva–foundation.
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# Untitled Document
Eclatement de réseaux de quadriques
et bord des instantons de degré $`3`$
Blowing-up of nets of quadrics and boundary
of mathematical instantons of degree $`3`$
Nicolas PERRIN
Ecole Normale Supérieure
6 rue Einstein
92 160 Antony
e-mail : nperrin@clipper.ens.fr
Résumé:
Dans leur article , L. Gruson et M. Skiti ont décrit une application birationnelle de la variété $`𝐈_\mathrm{𝟑}`$ des instantons de degré $`3`$ vers les réseaux de quadriques de $`\stackrel{ˇ}{𝐏^\mathrm{𝟑}}`$. Ils font ainsi apparaitre deux composantes du bord de $`𝐈_\mathrm{𝟑}`$ associées l’une au diviseur des réseaux contenant une quadrique de $`\stackrel{ˇ}{𝐏^\mathrm{𝟑}}`$ dégénérée en deux plans et l’autre au diviseur des réseaux de Lüroth. Dans cet article, on décrit une composante irréductible du bord de $`𝐈_\mathrm{𝟑}`$ comme le diviseur exceptionnel de l’éclatement du fermé des réseaux de quadriques formé par les quadriques de rang $`3`$.
Abstract :
In their article , L. Gruson and M. Skiti have constructed a birationnal map from the variety $`𝐈_\mathrm{𝟑}`$ of mathematical instantons of degree $`3`$ to the variety of nets of quadrics in $`\stackrel{ˇ}{𝐏^\mathrm{𝟑}}`$. They describe by this way two irreducible componants of the boundary of $`𝐈_\mathrm{𝟑}`$ associated to the divisor of nets which contain a two-plane degenerated quadric and the divisor of Lüroth nets. In this article we describe an irreducible componante of the boundary of $`𝐈_\mathrm{𝟑}`$ as the exceptionnal divisor of the blowing-up of the closed set of nets of quadrics of rank $`3`$.
Remerciements :
Je tiens à remercier ici mon directeur de thèse Laurent Gruson pour toute l’aide qu’il m’a apportée durant la préparation de ce travail.
Introduction :
Soit $`V`$ un espace vectoriel de dimension $`4`$ sur $`𝐂`$, on note $`𝐏^3`$ l’espace projectif $`𝐏(V)`$. Un instanton de degré $`n`$, est un élément de $`𝐌_{𝐏^\mathrm{𝟑}}(0,n,0)`$ (l’espace des modules des faisceaux semi-stables sans torsion de classes de chern $`(0,n,0)`$) qui est localement libre, stable et qui vérifie la propriété cohomologique $`h^1E(2)=0`$. On connait toutes les composantes irréductibles du bord des instantons de degré $`1`$ et $`2`$ (voir par exemple pour le degré $`1`$ et pour le degré $`2`$). On se propose ici d’étudier une composante irréductible du bord de la variété $`𝐈_\mathrm{𝟑}`$ des instantons de degré $`3`$.
L. Gruson et M. Skiti ont montré dans que la variété $`𝐈_\mathrm{𝟑}`$ des instantons de degré $`3`$ est birationnelle à la variété $`𝐑=\mathrm{G}(3,S^2\stackrel{ˇ}{V})`$ (qui paramétrise les sous espace vectoriels de dimension $`3`$ de $`S^2\stackrel{ˇ}{V}`$) des réseaux de quadriques de $`\stackrel{ˇ}{𝐏^\mathrm{𝟑}}`$. L’ouvert de $`𝐈_\mathrm{𝟑}`$ formé par les instantons sans droite trisauteuse s’envoie birationnellement sur un ouvert $`𝐑_4`$ de $`𝐑`$ formé par les réseaux $`R`$ de quadriques tels que l’application $`RV\stackrel{ˇ}{V}`$ soit de rang $`4`$.
Cette description leur a permis d’identifier deux composantes irréductibles du bord de $`𝐈_\mathrm{𝟑}`$ données par les hypersurfaces $`𝐑^{}`$ (respectivement $`𝐑^{\prime \prime }`$) des réseaux contenant une quadrique décomposée en deux plans (respectivement des réseaux de Lüroth). On va ici décrire une troisième composante du bord. L’idée consiste à éclater certaines sous variétés de $`𝐑`$ afin de pouvoir prolonger le morphisme des réseaux vers les faisceaux. On considère ainsi la sous variété $`𝐅`$ de la Grassmannienne $`\mathrm{Grass}(4,V)`$ des quotients de rang $`4`$ de $`V`$ au dessus de $`𝐑`$ (où $``$ est le sous fibré tautologique de $`𝐑`$) qui vérifie le fait que $`V𝒲`$ factorise l’application $`V𝒪\stackrel{ˇ}{V}`$ (on a ici noté $`𝒲`$ le quotient tautologique de $`\mathrm{Grass}(4,V))`$. Remarquons que $`\pi :𝐅𝐑`$ est un isomorphisme au dessus de $`𝐑_4`$. Notons $`𝐑_i`$ le localement fermé des réseaux tels que la flèche $`V\stackrel{\psi }{}\stackrel{ˇ}{V}`$ est de rang $`i`$ et posons $`𝐅_i=\pi ^1(𝐑_i)`$.
Remarque : Le schéma $`𝐅_3`$ est irréductible et réduit. En effet, on a un morphisme de $`𝐅_3`$ vers $`\mathrm{G}(3,\stackrel{ˇ}{V})`$ qui a un quotient $`W`$ de rang $`4`$ de $`V`$ associe l’image de la composée $`VW\stackrel{ˇ}{V}`$ qui est de dimension $`3`$ car on est dans $`𝐅_3`$. Mais alors la fibre de ce morphisme au dessus d’un sous espace $`K`$ de dimension $`3`$ de $`\stackrel{ˇ}{V}`$ est donnée par $`\mathrm{G}(3,S^2K\stackrel{ˇ}{V})`$ donc on voit que $`𝐅_3`$ est donné par $`\mathrm{G}(3,S^2𝒦\stackrel{ˇ}{V})`$$`𝒦`$ est le sous fibré tautologique de $`\mathrm{G}(3,\stackrel{ˇ}{V})`$. Grace à cette description, on sait que $`𝐅_3`$ est irréductible et réduit.
Proposition: Au voisinage de $`𝐅_3`$, le morphisme $`\pi :𝐅𝐑`$ est l’éclatement de $`𝐑_3`$.
Démonstration: Prenons un réseau $`R`$ de $`𝐑_3`$ et plaçons nous sur un ouvert affine de $`𝐑`$ contenant ce point. On sait alors qu’il existe un mineur $`3\times 3`$ de $`RV\stackrel{ˇ}{V}`$ qui est inversible. Ceci signifie que l’on a un sous fibré $`Q`$ de rang $`3`$ de $`V`$ et un sous fibré $`K`$ de $`\stackrel{ˇ}{V}`$ tels que la restriction de $`\psi `$ à $`Q`$ est un isomorphisme sur $`K`$. On peut alors se placer sur un ouvert affine de $`𝐑`$ d’anneau $`A`$ tel que l’application $`QA\stackrel{\psi }{}KA`$ est inversible en tout point de $`\mathrm{Spec}(A)`$. On a alors le diagramme suivant de $`A`$-modules :
$$\begin{array}{ccc}Q& \stackrel{}{}& K\\ & & \\ V& \stackrel{\psi }{}& 𝒪\stackrel{ˇ}{V}\\ & & \\ A^9& \stackrel{\phi }{}& A\end{array}$$
Le fermé $`𝐑_3`$ dans l’ouvert affine $`\mathrm{Spec}(A)`$ est donné par le $`0^{ieme}`$ idéal de Fitting de $`\phi `$ c’est à dire l’annulation de $`\phi `$. L’éclatement de ce lieu singulier est alors un schéma $`X`$ au dessus de $`\mathrm{Spec}(A)`$ solution du problème universel suivant : si $`X^{}`$ est un schéma muni d’un morphisme $`f`$ vers $`\mathrm{Spec}(A)`$ tel que l’image de $`f^{}\phi :f^{}(A^9)f^{}(A)`$ est un idéal inversible alors le morphisme de $`X`$ vers $`\mathrm{Spec}(A)`$ se factorise par $`f`$.
Il reste à montrer que $`𝐅`$ au dessus de cet ouvert est également solution de ce problème universel. Or au dessus de $`\mathrm{Spec}(A)`$, le schéma $`𝐅`$ vérifie la propriété universelle suivante : Si $`X^{}`$ est un schéma muni d’un morphisme $`f`$ vers $`\mathrm{Spec}(A)`$ et muni d’un module $`W`$ localement libre de rang $`4`$ tel que $`f^{}\psi :f^{}(V)f^{}(𝒪\stackrel{ˇ}{V})`$ se factorise par $`W`$, le morphisme de $`𝐅`$ vers $`\mathrm{Spec}(A)`$ se factorise par $`f`$. Soit un schéma $`X^{}`$ muni d’un morphisme $`f`$ vers $`\mathrm{Spec}(A)`$ et tel que l’image de $`f^{}\phi :f^{}(A^9)f^{}(A)`$ est un idéal inversible $`I`$, on construit alors un module localement libre de rang quatre $`W`$ tel que $`f^{}\psi `$ se factorise par $`W`$. En effet, il suffit de poser $`W=f^{}QI`$ qui est localement libre de rang $`4`$ et on a une application $`f^{}(V)W`$ qui factorise le morphisme $`f^{}\psi `$. La propriété universelle de $`𝐅`$ au dessus de $`\mathrm{Spec}(A)`$ nous permet alors de dire que l’on a une application de $`𝐅`$ vers $`X^{}`$ ce qui nous donne le résultat.
Remarque : M. Skiti, dans un travail en préparation montre qu’au voisinage de $`𝐅_2`$, le morphisme $`\pi `$ est l’éclatement de $`𝐑_2`$. Il en déduit une description de la sous variété (notée $`𝐈_\mathrm{𝟑}^\mathrm{𝟏}`$) de $`𝐈_\mathrm{𝟑}`$ des instantons ayant une droite trisauteuse. L’application birationnelle de $`𝐑`$ dans $`𝐈_\mathrm{𝟑}`$ s’étend à $`𝐅`$ de telle sorte qu’elle devienne un isomorphisme sur un voisinage de $`𝐅_2`$ dans $`𝐅`$. Elle identifie $`𝐅_2`$ et $`𝐈_\mathrm{𝟑}^\mathrm{𝟏}`$. Ceci permet de décrire $`𝐈_\mathrm{𝟑}^\mathrm{𝟏}`$ comme l’éclatement de $`𝐑_2`$ dans $`𝐑`$. On va faire la même construction avec $`𝐅_3`$.
On étudie la sous variété $`𝐈_\mathrm{𝟑}`$ de $`𝐌_{𝐏^\mathrm{𝟑}}(0,3,0)`$ formée par les faisceaux sans torsion non localement libres dont le bidual est un instanton de degré $`1`$ et tels que le conoyau de l’injection canonique dans ce bidual est une théta-caractéristique (décalée de $`2`$) sur une conique lisse. Les faisceaux $`E`$ de cette famille sont donc donnés par les noyaux de surjections $`E^{\prime \prime }\theta (2)`$$`E^{\prime \prime }`$ est un instanton de degré $`1`$ et $`\theta `$ est une théta-caractéristique sur une conique lisse. Cette famille est irréductible de dimension $`20`$. Sur un ouvert de $`𝐈_\mathrm{𝟑}`$, les faisceaux $`E`$ ont une cohomologie naturelle.
Sur l’ouvert $`𝐔`$ de $`𝐌_{𝐏^\mathrm{𝟑}}(0,3,0)`$ formé par les faisceaux à cohomologie naturelle (donc minimale), on sait définir un morphisme vers $`𝐅`$. Si $`E`$ est un tel faisceau, alors on a les égalités $`h^1E(1)=3`$, $`h^1E=4`$ et $`h^1E(1)=1`$. Ainsi, on définit un premier morphisme $`f_0`$ vers $`𝐑`$ qui prolonge celui de en associant à $`E`$ le réseau $`H^1E(1)H^1E(1)S^2\stackrel{ˇ}{V}`$. Au dessus de ce réseau, on associe à $`E`$ un quotient de rang $`4`$ de $`RV`$ donné par $`H^1E(1)VH^1E`$ qui nous donne le morphisme $`f`$ souhaité. Ce morphisme est ainsi défini sur l’ouvert $`𝐔`$ de $`𝐈_\mathrm{𝟑}`$ formé par les faisceaux à cohomologie naturelle.
Il existe sur un ouvert de $`𝐑`$ (et donc sur un ouvert de $`𝐅`$) un morphisme $`g`$ réciproque (voir ). On va le prolonger à un ouvert de $`𝐅_3`$.
Proposition: Le morphisme $`f`$ restreint à $`𝐔`$ est à valeurs dans $`𝐅_3`$ et est dominant sur $`𝐅_3`$. De plus sur son image, on définit une réciproque $`g`$ à $`f`$.
Démonstration: On commence par montrer que sur l’image de $`𝐔`$ on sait définir une réciproque à $`f`$. En effet, la réciproque est donnée de la façon suivante : soit $`R`$ le réseau et $`W`$ le quotient de $`RV`$. Les applications linéaires $`RVW`$ et $`W\stackrel{ˇ}{V}`$ nous donnent un complexe :
$$R\mathrm{\Omega }^2(2)W\mathrm{\Omega }^1(1)𝒪_{𝐏^\mathrm{𝟑}}$$
et la cohomologie au centre (décalée de $`1`$) de ce complexe nous donne le faisceau recherché. En effet, si le réseau $`R`$ et le quotient $`W`$ sont assez généraux, le faisceau ainsi obtenu est dans $`𝐌_{𝐏^\mathrm{𝟑}}(0,3,0)`$. Soit $`E`$ dans $`𝐔`$, la suite spectrale de Beilinson nous dit que le faisceau $`E(1)`$ est la cohomologie du complexe précédent si on prend $`R=H^1E(1)`$, $`W=H^1E`$ et que l’on identifie $`H^1E(1)`$ à $`𝐂`$. Les applications linéaires sont les multiplications du module de Rao de $`E`$. Ainsi, sur l’image de $`𝐔`$ par $`f`$ on a une réciproque $`g`$ qui prolonge le morphisme défini sur $`𝐑`$.
Montrons maintenant que cette image est contenue dans $`𝐅_3`$. Pour ceci, il suffit de voir que le réseau associé, qui est $`H^1E(1)H^1E(1)S^2\stackrel{ˇ}{V}`$ est tel que l’application $`H^1E(1)VH^1E(1)\stackrel{ˇ}{V}`$ est de rang trois. Cette application se décompose en deux applications : $`H^1E(1)VH^1E`$ qui est génériquement surjective et $`H^1EH^1E(1)\stackrel{ˇ}{V}`$. Mais si notre faisceau est donné par la suite exacte
$$0EE^{\prime \prime }\theta (2)0$$
alors $`H^1E`$ s’identifie à $`H^0\theta (2)`$ et $`H^1E(1)`$ est un quotient de rang $`1`$ de $`H^0\theta (3)`$. La multiplication du module de Rao est donc nulle pour l’élément $`H`$ de $`V`$ qui définit le plan de la conique. Ainsi l’application $`H^1EH^1E(1)\stackrel{ˇ}{V}`$ est de rang trois et a pour image $`(V/H\stackrel{ˇ}{)}`$. Par conséquent l’application $`H^1E(1)VH^1E(1)\stackrel{ˇ}{V}`$ est aussi de rang trois car $`H^1E(1)VH^1E`$ est surjective.
On sait maintenant que $`f`$ a un morphisme réciproque sur $`f(𝐔)`$ qui est contenu dans $`𝐅_3`$. On sait donc déjà que l’image de $`𝐔`$ est de dimension $`20`$. Mais alors, comme $`𝐅_3`$ est réduit irréductible et que sa dimension est aussi $`20`$, on sait que $`f(𝐔)`$ contient un ouvert de $`𝐅_3`$ et $`f|_𝐔`$ est donc dominant sur $`𝐅_3`$.
Corollaire : Le prolongement de $`g`$ est birationnel au voisinage de $`𝐅_3`$. Dans la description de $`𝐈_\mathrm{𝟑}`$ avec les réseaux, la variété $`𝐈_\mathrm{𝟑}`$ est donc le diviseur exceptionnel de l’éclatement de $`𝐑_3`$ dans $`𝐑`$. La famille $`𝐈_\mathrm{𝟑}`$ forme une composante irréductible du bord de $`𝐈_\mathrm{𝟑}`$.
Démonstration: On a vu que $`g`$ est défini sur $`f(𝐔)`$ qui contient un ouvert de $`𝐅_3`$. Ainsi, sur un voisinage (dans $`𝐅`$) de cet ouvert, le morphisme $`g`$ est bien un morphisme réciproque à $`f`$. Le morphisme $`g`$ est bien birationnel au voisinage de $`𝐅_3`$.
De plus, on a vu que $`𝐅𝐑`$ est l’éclatement de $`𝐑_3`$ au voisinage de $`𝐅_3`$ donc $`g`$ identifie cette situation à celle de $`𝐈_\mathrm{𝟑}`$ et $`𝐈_\mathrm{𝟑}`$. On voit donc que $`𝐈_\mathrm{𝟑}`$ est adhérente à $`𝐈_\mathrm{𝟑}`$ et que dans la description de $`𝐈_\mathrm{𝟑}`$ avec les réseaux, la variété $`𝐈_\mathrm{𝟑}`$ est le diviseur exceptionnel de l’éclatement de $`𝐑_3`$ dans $`𝐑`$ et forme donc une composante irréductible du bord de $`𝐈_\mathrm{𝟑}`$.
Remarque : On sait décrire les éléments de saut d’un faisceau $`E`$ de $`𝐈_\mathrm{𝟑}`$.
Les plans instables forment une courbe dans $`\stackrel{ˇ}{𝐏^\mathrm{𝟑}}`$ de degré $`6`$ et de genre $`3`$ ayant un point triple au point correspondant au plan de la conique et dont la courbe des trisécantes est tracées sur le complexe de droites associé au bidual.
Les droites bisauteuses forment une courbe de degré $`8`$ et de genre $`3`$ de la Grassmannienne qui est réunion d’une quintique rationnelle tracée sur le complexe de droites défini par le bidual et d’une cubique tracée dans le $`(\beta )`$-plan des droites du plan de la conique.
Il y a un lien entre ces deux courbes de saut. Les droites bisauteuses sont les trisécantes à la courbe des plans instables. Réciproquement, la courbe des plans instables forme le lieu triple de la surface réglée décrite par les droites bisauteuses.
Références :
L. Gruson et M. Skiti : $`3`$-instantons et réseaux de quadriques. Math. Ann. 298 (1994).
M.S. Narashiman et G. Trautmann : Compactification of $`M_{𝐏^\mathrm{𝟑}}(0,2)`$ and Poncelet pairs of conics. Pacific Journal of Mathematics 145 (1990).
C. Okonek, M. Schneider et H. Spindler : Vector bunbles on complex projective spaces. Basel, Boston, Stuttgart : Birkhäuser 1980.
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# Analytic Light-Curves of Gamma-Ray Burst Afterglows: Homogeneous versus Wind External Media
## 1 Introduction
One of the most important issues regarding Gamma-Ray Bursts (GRBs) is the nature of the object that releases the relativistic ejecta generating the high energy emission of the main event and the lower frequency emission during the ensuing afterglow. Some insight about the GRB progenitor can be obtained from the properties of the circum-burst medium, which can be inferred from the features of the afterglow emission. If the ejecta is expelled during the merging of two compact objects (Mészáros & Rees 1997b), it is expected that the medium surrounding the GRB source is homogeneous. However, if a collapsing massive star (Woosley 1993, Paczyński 1998) is the origin of the relativistic fireball, the circum-burst medium is the wind ejected by the star prior to its collapse, whose density decreases outwards. The two models differ in the dependence on radius of the particle density of the circum-burst medium which the GRB remnant interacts with, and in the value of this density at the deceleration length-scale. The former modifies the rate of decline of the afterglow, while the latter determines the overall afterglow brightness. Therefore it is possible to correlate afterglow emission features to a specific type of external medium.
Significant work in this direction has been done by many researchers. The two afterglows that exhibited breaks consistent with the effects arising from strong collimation of ejecta – GRB 990123 (Kulkarni et al. 1999a) and GRB 990510 (Stanek et al. 1999, Harrison et al. 1999) – indicate that the external gas was homogeneous (recent work by Kumar & Panaitescu 2000 shows that jets interacting with winds cannot produce sharp breaks in the afterglow light-curve). The optical emission of three afterglows had a steeper than usual decline. GRB 970228 decayed as $`T^{1.7}`$ after the subtraction of an underlying supernova emission (Reichart 1999, Galama et al. 2000). The light-curve of GRB 980326 fell off as $`T^{2.1}`$ (Groot et al. 1998) and an emission in excess of the early time extrapolation was detected $`20`$ days after the main event, indicating a supernova contribution (Bloom et al. 1999). A $`T^2`$ decay was observed for the afterglow of GRB 980519 (Halpern et al. 1999). Such steep declines can be produced either by a fireball interacting with a pre-ejected wind (Chevalier & Li 1999) and an electron index around 3, or by a narrow jet expanding laterally in a homogeneous external medium and an electron index slightly larger than 2. Chevalier & Li (1999) found that the radio emission of the afterglow of GRB 980519 is consistent with an external wind; however Frail et al. (2000) point out that the interstellar scintillation present in the radio data does not allow ruling out the jet model. Nevertheless, the existence of supernovae associated with GRB 970228 and GRB 980326 points toward a massive star as the origin of these bursts, implying a pre-ejected wind as the external medium. From the analysis of the optical radio emission of the afterglow of GRB 970508, Chevalier & Li (2000) conclude that the surrounding medium was a wind. Frail, Waxman & Kulkarni (2000) argue that the same radio afterglow can be explained by a homogeneous external medium.
In this work we investigate the differences between the light-curves of afterglows arising for each type of external medium, with the aim of finding ways for distinguishing between the two models. This study is done within the usual framework of a relativistic remnant interacting with a cold external gas. As the fireball is decelerated, a shock front sweeps up the external gas, accelerating relativistic electrons and generating a magnetic field in the shocked gas. We ignore the emission from electrons accelerated by the reverse shock which propagates through the ejecta at very early times. At optical wavelengths this emission is short lived, lasting up to few tens of seconds after the main event (Sari & Piran 1999), but it could be important for the radio emission until few days (Kulkarni et al. 1999b).
Analytical afterglow light-curves for spherical remnants interacting with homogeneous external media have been previously published by Sari, Piran & Narayan (1998). Features of afterglows from spherical fireballs, such as peak flux, break frequencies, and time evolution of fluxes at a fixed frequency, have been studied by Mészáros & Rees (1997a), Waxman (1997b), Wijers & Galama (1999), and Dai & Lu (2000) for homogeneous media, and by Chevalier & Li (2000) for pre-ejected winds. In this work we present and compare analytical and numerical light-curves at various observing frequencies, covering all cases of interest, for both types of external media, taking into account the differential arrival-time delay and Doppler boosting due to the spherical shape of the source. We take into account first order IC scattering, calculate its effect on the electron cooling and on the afterglow synchrotron emission, and study briefly the high-energy emission resulting from the up-scattering of synchrotron photons. The possible importance of IC scatterings for the early afterglow emission was pointed out by Waxman (1997a) and Wei & Lu (1998).
## 2 Simple Dynamics of Relativistic Remnants
For the calculation of the afterglow emission it is necessary to know how the remnant Lorentz factor $`\mathrm{\Gamma }`$ evolves with observer time $`T`$, as all other quantities that appear in the expression of the spectral flux are functions of $`\mathrm{\Gamma }`$ and of the remnant radius $`r`$ and external medium density $`n(r)`$. We shall assume that the remnant is adiabatic, i.e. the energy carried away by the emitted photons is a negligible fraction of the total energy of the remnant. This assumption is correct if the energy density of the electrons accelerated at the shock front is a fraction $`\epsilon _e1`$ of the total energy density in post-shock fluid or if most of the electrons are adiabatic, i.e. their radiative cooling timescale exceeds that of the adiabatic losses due to the remnant expansion.
Assuming that the internal energy of the ejecta is negligible compared to its rest-mass energy and that the ratio internal-to-rest mass energy in the energized external medium is $`\mathrm{\Gamma }1`$ (i.e. its “temperature” tracks that of the freshly shocked gas), conservation of energy leads to
$$m(r)\mathrm{\Gamma }^2+M_{fb}\mathrm{\Gamma }[m(r)+M_{fb}\mathrm{\Gamma }_0]=0$$
(1)
where $`M_{fb}`$ and $`\mathrm{\Gamma }_0`$ are the initial mass and Lorentz factor of the fireball (whose energy is $`E=M_{fb}\mathrm{\Gamma }_0`$) and
$$m(r)=\frac{4\pi }{3s}m_pn(r)r^3$$
(2)
is the mass of swept-up material ($`m_p`$ being the proton’s mass). The external medium particle density is
$$n(r)=Ar^s,$$
(3)
with $`s=0`$ for a homogeneous medium and $`s=2`$ for a wind ejected by the GRB progenitor at a constant speed.
Equation (2) is valid if the remnant is spherical, but can also be used for collimated ejecta when the lateral spreading (Rhoads 1999) is insignificant if the quantity $`E`$ above is defined as the energy the fireball would have if it were spherical. Throughout this work we shall assume that the remnant is a jet with an initial half-angle larger than $`>20^\mathrm{o}`$, in which case the sideways expansion is negligible during the relativistic phase. The following analytical calculations of the afterglow emission can be extended to sideways expanding jets and non-relativistic remnants by first determining $`\mathrm{\Gamma }(r)`$. The set of coupled differential equations describing the evolution of the jet Lorentz factor and its opening can be solved analytically for $`s=0`$ (Rhoads 1999). The lack of a good approximation for the jet dynamics in the case of pre-ejected winds is the main motivation for restricting the following analytical calculations to spherical or wide-angle remnants.
The solution of equation (1) is
$$\mathrm{\Gamma }(r)=\frac{1}{2}\left[\sqrt{4x^{3s}+1+(2x^{3s}/\mathrm{\Gamma }_0)^2}1\right]x^{s3}\mathrm{\Gamma }_0,$$
(4)
where $`x`$ is the radial coordinate $`r`$ scaled to
$$r_0=\left(\frac{3s}{4\pi }\frac{E}{m_pc^2A\mathrm{\Gamma }_0^2}\right)^{1/(3s)},$$
(5)
the deceleration length-scale, at which $`m(r_0)=E/(c^2\mathrm{\Gamma }_0^2)=M_{fb}/\mathrm{\Gamma }_0`$. The result given in equation (4) is also valid in the non-relativistic regime. For $`x1`$ $`\mathrm{\Gamma }<\mathrm{\Gamma }_0`$, while for $`1xx_{nr}`$ we find $`\mathrm{\Gamma }=x^{(3s)/2}\mathrm{\Gamma }_0`$. Here $`x_{nr}=(\mathrm{\Gamma }_0^2/3)^{1/(3s)}`$ marks the end of the relativistic regime: $`\mathrm{\Gamma }(x_{nr})=2`$. For the ease of analytical calculations we shall assume that the power-law behavior of $`\mathrm{\Gamma }`$ lasts from $`x=1`$ to $`x=x_{nr}`$.
The Lorentz factor given in equation (4) represents a “dynamical” average of the Lorentz factors at which different regions of the shocked remnant move (the Blandford – McKee solution). The Lorentz factor of the shock front that propagates into the external gas $`\mathrm{\Gamma }_{sh}=\sqrt{2}\mathrm{\Gamma }`$, with $`\mathrm{\Gamma }`$ given by equation (4), matches that given in equation (69) of Blandford & McKee (1976) for the power-law regime $`1xx_{nr}`$ if $`E`$ is multiplied by $`(174s)/(124s)`$. This correction factor ($`17/8`$ for $`s=0`$ and $`9/4`$ for $`s=2`$) will be used in the following results.
The constant $`A`$ in equation (3) is the number density $`n_{}`$ of the external homogeneous medium for $`s=0`$, while for $`s=2`$
$$A=\frac{1}{4\pi }\frac{\stackrel{.}{M}}{m_pv}=3.0\times 10^{35}A_{}\mathrm{cm}^1,$$
(6)
where $`\stackrel{.}{M}`$ is the mass loss rate of the massive star that ejected the wind at constant speed $`v`$, and $`A`$ was scaled to
$$A_{}=\frac{\stackrel{.}{M}/10^5\mathrm{M}_{}\mathrm{yr}^1}{v/10^3\mathrm{km}\mathrm{s}^1},$$
(7)
as in Chevalier & Li (2000) for a Wolf-Rayet star. For further calculations it is convenient to use equation (3) in the form $`n(r)=n_0(r/r_0)^s`$ where $`n_0=n_{}`$ for $`s=0`$, while for $`s=2`$
$$n_0=1.9\times 10^4E_{53}^2\mathrm{\Gamma }_{0,2}^4A_{}^3\mathrm{cm}^3$$
(8)
is the wind particle density at the deceleration radius (eq. ):
$$(s=2)r_0=4.0\times 10^{15}E_{53}\mathrm{\Gamma }_{0,2}^2A_{}^1\mathrm{cm}.$$
(9)
The usual notation $`C_n=10^nC`$ is used throughout this work.
Note that, for the reference values used here, the deceleration radius in the wind model is 1.5 orders of magnitude lower than that for a homogeneous external medium:
$$(s=0)r_0=1.3\times 10^{17}E_{53}^{1/3}\mathrm{\Gamma }_{0,2}^{2/3}n_{,0}^{1/3}\mathrm{cm},$$
(10)
and even smaller for higher initial fireball Lorentz factors or slower winds (i.e. a larger parameter $`A_{}`$). If GRBs are due to internal shocks occurring in unstable relativistic fireballs (Rees & Mészáros 1994, Paczyński & Xu 1994, Piran 1999), then the external shock resulting from the interaction of the fireball with the pre-ejected (non-relativistic) wind may occur before the internal shocks are over. In this case successive internal collisions occur when faster parts of the ejecta catch up with the decelerating leading edge of the fireball, a scenario suggested within the $`s=0`$ model by Fenimore & Ramirez-Ruiz (2000), but which is more likely to happen if the external medium is the gas ejected by a massive star. The GRB itself would then exhibit the erratic variability characteristic of internal shocks until a time of the order
$$\frac{r_0}{c\mathrm{\Gamma }_0^2}10E_{53}\mathrm{\Gamma }_{0,2}^4A_{}^1s,$$
(11)
(which has a strong dependence on $`\mathrm{\Gamma }_0`$), after which there may be significant emission from internal shocks on the outermost part of the fireball and from the external shock that plows through the external gas. The former mechanism generates pulses of increasing duration as the fireball expands, while the later leads to a continuous emission.
The time $`T`$ when the observer receives a photon emitted along the line of sight toward the fireball center can be calculated by integrating
$$\mathrm{d}T=(1\beta )\mathrm{d}t=\frac{1}{2}\frac{\mathrm{d}t}{\mathrm{\Gamma }^2},$$
(12)
where $`\beta `$ is the shocked fluid speed and $`t=r/c`$ is the time measured in the laboratory frame. Approximating the solution given in equation (4) with $`\mathrm{\Gamma }=\mathrm{\Gamma }_0`$ for $`x<1`$ and $`\mathrm{\Gamma }=x^{(3s)/2}\mathrm{\Gamma }_0`$ for $`1<x<x_{nr}`$, one obtains
$$T=T_0(x^{4s}+3s),T_0\frac{1}{2(4s)}\frac{r_0}{c\mathrm{\Gamma }_0^2}.$$
(13)
Note that $`T`$ given in equation (13) is the earliest time a photon emitted by the remnant at time $`t`$ can reach the observer. Photons emitted by the fluid moving at an angle $`\theta =1/\mathrm{\Gamma }`$ off the center–observe axis arrive at $`T=t(1\mathrm{cos}\theta )=t/2\mathrm{\Gamma }^2`$, which is a factor $`4s`$ larger than the time corresponding to $`\theta =0`$: $`T=t/2(4s)\mathrm{\Gamma }^2`$.
From equation (13) $`r(T)`$ can be found and then substituted in the expressions for $`\mathrm{\Gamma }(r)`$ and $`n(r)`$ to obtain these quantities as a function of the observer time. For the power-law phase the results are
$$(s=0)\mathrm{\Gamma }(T)=6.3E_{53}^{1/8}n_{,0}^{1/8}T_d^{3/8},$$
(14)
$$(s=0)r(T)=8.2\times 10^{17}E_{53}^{1/4}n_{,0}^{1/4}T_d^{1/4}\mathrm{cm},$$
(15)
$`n=n_{}`$ (constant) for a homogeneous external medium and
$$(s=2)\mathrm{\Gamma }(T)=7.9E_{53}^{1/4}A_{}^{1/4}T_d^{1/4},$$
(16)
$$(s=2)r(T)=6.4\times 10^{17}E_{53}^{1/2}A_{}^{1/2}T_d^{1/2}\mathrm{cm},$$
(17)
$$(s=2)n(T)=0.73E_{53}^1A_{}^2T_d^1\mathrm{cm}^3,$$
(18)
for an external wind, $`T_d`$ being the observer time measured in days. Note that, at least for the scaling values chosen here, $`\mathrm{\Gamma }`$, $`r`$, and $`n`$ have about the same values at $`T=1`$ day in both models. Also note that the above quantities (and thus the afterglow emission) are independent of the fireball initial Lorentz factor $`\mathrm{\Gamma }_0`$.
## 3 Break Frequencies
Within the synchrotron emission model there are three expected breaks in the afterglow spectrum: $`(i)`$ an injection break, at the synchrotron frequency $`\nu _i`$ at which the bulk of the electrons injected by the shock front radiate, $`(ii)`$ a cooling break, at the synchrotron frequency $`\nu _c`$ of electrons whose radiative cooling time equals the expansion timescale, and $`(iii)`$ an absorption break, at $`\nu _a`$ below which the synchrotron photons are absorbed by electrons in free-free transitions in a magnetic field (synchrotron self-absorption).
The break frequencies can be calculated if the distribution of the injected electrons and the strength of the magnetic field are known. The distribution of the injected electrons is assumed to be $`𝒩_i(\gamma )\gamma ^p`$ starting from a minimum random Lorentz factor given by
$$\gamma _i=\frac{m_p}{m_e}\epsilon _e(\mathrm{\Gamma }1),$$
(19)
where $`m_e`$ is the electron mass. The energy carried by this electron distribution is a fraction $`\frac{p1}{p2}\epsilon _e`$ of the total internal energy. The post-shock magnetic field strength in the co-moving frame is given by
$$\frac{B^2}{8\pi }=\epsilon _Bm_pc^2n_e^{}(\mathrm{\Gamma }1)=4\epsilon _Bm_pc^2n(r)(\mathrm{\Gamma }1)\left(\mathrm{\Gamma }+\frac{3}{4}\right),$$
(20)
where $`\epsilon _B`$ is the fractional energy carried by the magnetic field and $`n_e^{}`$ is the co-moving frame electron density behind the shock front. Equations (19) and (20) are based on that the internal to rest-mass energy density ratio in the shocked fluid is $`\mathrm{\Gamma }1`$; the derivation of the latter equation also used that the co-moving particle density is $`4\mathrm{\Gamma }+3`$ times larger than that ahead of the shock.
### 3.1 Injection Break
Using the relativistic Doppler factor $`2\mathrm{\Gamma }`$ corresponding to the motion of the source toward the observer (i.e. $`\theta =0`$), the synchrotron emission from a power-law distribution of electrons peaks at the observer frame frequency
$$\nu _i=\frac{3x_p}{2\pi }\frac{e}{m_ec}\gamma _i^2B\mathrm{\Gamma }=8.4\times 10^6x_p\gamma _i^2B\mathrm{\Gamma }\mathrm{Hz},$$
(21)
where the factor $`x_p`$ is calculated in Wijers & Galama (1999) for various values of the electron index $`p`$. We shall use $`x_p=0.52`$, which is strictly correct only for $`p=2.5`$. With the aid of equations (14), (16), (18), (19), and (20) one obtains:
$$(s=0)\nu _i=0.92\times 10^{13}E_{53}^{1/2}\epsilon _{e,1}^2\epsilon _{B,2}^{1/2}T_d^{3/2}\mathrm{Hz},$$
(22)
$$(s=2)\nu _i=1.9\times 10^{13}E_{53}^{1/2}\epsilon _{e,1}^2\epsilon _{B,2}^{1/2}T_d^{3/2}\mathrm{Hz}.$$
(23)
Note that $`\nu _i`$ has the same scalings with the model parameters for $`s=0`$ and $`s=2`$. The ratio of the two frequencies is $`\sqrt{\frac{17}{72}}`$.
### 3.2 Cooling Break
The relativistic electrons cool radiatively through synchrotron emission and IC scatterings of the synchrotron photons on a co-moving frame timescale
$$t_{rad}^{}(\gamma )=\frac{t_{sy}^{}(\gamma )}{Y+1}=\frac{6\pi }{Y+1}\frac{m_ec}{\sigma _e}\frac{1}{\gamma B^2},$$
(24)
where $`t_{sy}^{}(\gamma )`$ is the co-moving frame synchrotron cooling timescale of electrons of Lorentz factor $`\gamma `$, $`Y`$ is the Compton parameter, and $`\sigma _e`$ the cross-section for electron scatterings<sup>1</sup><sup>1</sup>1 The up-scattering of the $`\nu _i`$ synchrotron photons on the $`\gamma _c`$\- and $`\gamma _i`$-electrons occurs in the Thomson regime for $`T>10^3`$ day, i.e. the scattering cross-section in equation (24) is not reduced by the Klein-Nishina effect. Using equation (20), the Lorentz factor of the electrons that cool radiatively on a timescale equal to the remnant age
$$t^{}=\frac{1}{c}\frac{\mathrm{d}r}{\mathrm{\Gamma }}=\frac{2}{5s}\frac{r}{c\mathrm{\Gamma }}$$
(25)
can be written as
$$\gamma _c=\frac{3\pi (5s)}{Y+1}\frac{m_ec^2}{\sigma _e}\frac{\mathrm{\Gamma }}{B^2r}=\frac{77(5s)}{(Y+1)\epsilon _B}\frac{1}{n\mathrm{\Gamma }r_{18}},$$
(26)
with $`n`$ in $`\mathrm{cm}^3`$.
The observer-frame frequency $`\nu _c`$ of the cooling break is
$$(s=0)\nu _c=3.7\times 10^{14}E_{53}^{1/2}n_{,0}^1(Y+1)^2\epsilon _{B,2}^{3/2}T_d^{1/2}\mathrm{Hz},$$
(27)
$$(s=2)\nu _c=3.5\times 10^{14}E_{53}^{1/2}A_{}^2(Y+1)^2\epsilon _{B,2}^{3/2}T_d^{1/2}\mathrm{Hz}.$$
(28)
Hereafter we shall use the terminology “radiative electrons” for the case where the $`\gamma _i`$-electrons cool mostly through emission of radiation (i.e. $`t_{rad}^{}(\gamma _i)<t^{}`$ and $`\gamma _c<\gamma _i`$), and we shall refer to “adiabatic electrons” if the $`\gamma _i`$-electrons cool mostly adiabatically (i.e. $`t_{rad}^{}(\gamma _i)>t^{}`$ and $`\gamma _i<\gamma _c`$).
#### 3.2.1 Compton Parameter and Electron Distribution
For the calculation of the Compton parameter $`Y`$, we take into account only one up-scattering of the synchrotron photons. Multiple IC scatterings of the same photon have an important effect on the electron cooling only if the $`Y`$ parameter for single scatterings is above unity. As shown in §3.2.2 and §3.2.3 this occurs if $`1)`$ $`\epsilon _B<10^2\epsilon _{e,1}`$ and if $`2)`$ $`T<T_y`$, where $`T_y`$ is the time when $`Y`$ falls below unity.
Using equations (19) and (20), it can be shown that IC scatterings of order higher than two are suppressed by the Klein-Nishina effect. For adiabatic electrons, a second IC scattering occurs in the Thomson regime if $`3a)`$ $`T>1E_{53}^{1/3}n_{,0}^{1/9}\epsilon _{e,1}^{20/9}\epsilon _{B,4}^{2/9}\mathrm{d}`$ for $`s=0`$, and if $`3b)`$ $`T>1E_{53}^{1/2}A_{}^{1/4}\epsilon _{e,1}^{5/2}\epsilon _{B,4}^{1/4}\mathrm{d}`$ for $`s=2`$. If a second IC scattering is ignored in these cases then a higher energy component peaking around 1 GeV is left out, otherwise the synchrotron and first IC emissions remain unaltered if the electrons are adiabatic.
The effect of second order up-scatterings is more important when electrons are radiative, i.e. for $`4)`$ $`T<T_r`$ with $`T_r`$ calculated in §3.2.2, as in this case it reduces the intensity of the synchrotron and first stage IC components. With the aid of equations (19) and (20) it can be shown that for $`s=0`$ a second up-scattering occurs in the Thomson regime and at $`T>10^2\mathrm{d}`$ if $`5a)`$ $`n>25E_{53}^{7/11}\epsilon _{e,1}^{10/11}\epsilon _{B,3}^{2/11}\mathrm{cm}^3`$, while for $`s=2`$ the second IC emission is not suppressed by the Klein-Nishina effect if $`5b)`$ $`T>0.05E_{53}^{1/3}A_{}^{11/6}\epsilon _{e,1}^{5/6}\epsilon _{B,3}^{2/3}\mathrm{d}`$.
Concluding, second stage up-scatterings can be ignored if the set of conditions $`1),2),3)`$ or $`1),4),5)`$ are not simultaneously satisfied. For a single up-scattering, the Compton parameter is
$$Y=\frac{4}{3}_{𝒩_e}\gamma ^2d\tau _e=\frac{4}{3}\tau _e_{𝒩_e}\gamma ^2𝒩_e(\gamma )d\gamma ,$$
(29)
where $`𝒩_e(\gamma )`$ is the normalized electron distribution and $`\tau _e`$ is the optical thickness to electron scattering, given by
$$\tau _e=\frac{1}{4\pi }\frac{\sigma _em(r)}{m_pr^2}=\frac{1}{3s}\sigma _enr.$$
(30)
If the injected $`\gamma _i`$-electrons cool faster than the timescale of their injection, then $`\gamma _c`$ given by equation (26) is the typical electron Lorentz factor in the remnant, and the electron distribution in the shocked fluid can be approximated by
$$𝒩_e^{(r)}\{\begin{array}{cc}\gamma ^2\hfill & \gamma _c<\gamma <\gamma _i\hfill \\ \gamma ^{(p+1)}\hfill & \gamma _i<\gamma \hfill \end{array},$$
(31)
with $`p>2`$. In the opposite case most electrons have a random Lorentz factor $`\gamma _i`$, and the electron distribution is
$$𝒩_e^{(a)}\{\begin{array}{cc}\gamma ^p\hfill & \gamma _i<\gamma <\gamma _c\hfill \\ \gamma ^{(p+1)}\hfill & \gamma _c<\gamma \hfill \end{array}.$$
(32)
#### 3.2.2 Radiative Electrons
For $`\gamma _c\gamma _i`$ equations (29) and (31) lead to
$$Y_r=\frac{4}{3}\gamma _i\gamma _c\tau _e.$$
(33)
Substituting $`\gamma _c`$ with the aid of equation (26), one obtains
$$Y_r(Y_r+1)=\frac{5s}{8(3s)}\frac{n_e^{}m_ec^2\gamma _i}{B^2/8\pi }=\frac{5s}{8(3s)}\frac{\epsilon _e}{\epsilon _B},$$
(34)
where we used equations (19) and (20). Therefore the Compton parameter during the electron radiative phase is
$$Y_r=\frac{1}{2}\left(\sqrt{\frac{5s}{2(3s)}\frac{\epsilon _e}{\epsilon _B}+1}1\right).$$
(35)
Hence the electron cooling is dominated by IC scatterings (i.e. $`Y_r>1`$) for $`\epsilon _B<\stackrel{~}{\epsilon _B}`$, where
$$\stackrel{~}{\epsilon _B}=\frac{5s}{16(3s)}\epsilon _e.$$
(36)
Note that $`Y_r`$ is time-independent. Therefore the $`\nu _c`$ given by equations (27) and (28) decreases with time for $`s=0`$ and increases in the $`s=2`$ model. Thus, for observations made at a fixed frequency, the electrons emitting at that frequency change their cooling regime from adiabatic to radiative in the case of a homogeneous external gas, and from radiative to adiabatic for an external wind (Chevalier & Li 2000).
With the aid of equations (19), (20), and (24), it can be shown that the electrons are radiative if
$$nr\mathrm{\Gamma }^2>\frac{3(5s)}{32(Y_r+1)}\frac{(m_e/m_p)^2}{\sigma _e\epsilon _e\epsilon _B}=\frac{4.2\times 10^{16}}{\epsilon _e\epsilon _B}\frac{5s}{Y_r+1}\mathrm{cm}^2,$$
(37)
which, with the further use of equations (14) – (18), leads to the conclusion that the electrons are radiative until the observer time $`T_r`$ given by
$$(s=0)T_r=0.025E_{53}n_{,0}(Y_r+1)^2\epsilon _{e,1}^2\epsilon _{B,2}^2\mathrm{day},$$
(38)
$$(s=2)T_r=0.23A_{}(Y_r+1)\epsilon _{e,1}\epsilon _{B,2}\mathrm{day}.$$
(39)
#### 3.2.3 Adiabatic Electrons
For $`\gamma _i\gamma _c`$ equations (29) and (32) give the Compton parameter
$$Y_a=\frac{4}{3}\tau _e\times \{\begin{array}{cc}\gamma _i^{p1}\gamma _c^{3p}\hfill & 2<p<3\hfill \\ \gamma _i^2\hfill & 3<p\hfill \end{array}.$$
(40)
Case 1: $`2<p<3`$. By substituting equations (19) and (26) in equation (40), one obtains
$$Y_a(Y_a+1)^{3p}=_p(T)c_s(p)\epsilon _e^{p1}\epsilon _B^{p3}(n\mathrm{\Gamma }^2r_{18})^{p2},$$
(41)
where $`\mathrm{log}c_s(p)=(3p)\mathrm{log}(5s)\mathrm{log}(3s)+1.4p3.7`$. The Compton parameter can be obtained by solving numerically the above equation. For analytical purposes, one can approximate $`Y_a=_p`$ for $`_p<1`$, in which case the IC losses are less important, and
$$Y_a=_p^{\frac{1}{4p}}\mathrm{for}_p>1.$$
(42)
In the latter case the IC scatterings affect the electron cooling.
Note that the quantity $`n\mathrm{\Gamma }^2r`$ in equation (41) decreases with time. Thus, for $`\epsilon _B<\stackrel{~}{\epsilon _B}`$, the Compton parameter $`Y_a`$ is above unity until a time $`T_y`$ which can be determined by substituting equations (14) – (18) in (42):
$$(s=0)T_y=10^{\frac{83p}{p2}}E_{53}n_{,0}\epsilon _{e,1}^{2\frac{p1}{p2}}\epsilon _{B,3}^{2\frac{3p}{p2}}\mathrm{day},$$
(43)
$$(s=2)T_y=10^{\frac{4.91.6p}{p2}}A_{}\epsilon _{e,1}^{\frac{p1}{p2}}\epsilon _{B,3}^{\frac{3p}{p2}}\mathrm{day}.$$
(44)
Thus, for $`\epsilon _B<\stackrel{~}{\epsilon _B}`$ and $`T_r<T<T_y`$, the Compton parameter determines the evolution of the cooling break frequency (eqs. and ):
$$\nu _c\stackrel{(s=0)}{=}10^{15+\frac{2.5p5.5}{4p}}\left[E_{53}^{\frac{p}{2}}n_{,0}^2\epsilon _{e,1}^{2(p1)}\epsilon _{B,3}^{\frac{p}{2}}T_d^{\frac{3p8}{2}}\right]^{\frac{1}{4p}}\mathrm{Hz},$$
(45)
$$\nu _c\stackrel{(s=2)}{=}10^{15+\frac{2.2p5.5}{4p}}E_{53}^{\frac{1}{2}}\left[A_{}^4\epsilon _{e,1}^{2(p1)}\epsilon _{B,3}^{\frac{p}{2}}T_d^{\frac{3p4}{2}}\right]^{\frac{1}{4p}}\mathrm{Hz}.$$
(46)
Note that for a homogeneous medium ($`s=0`$) and $`8/3<p<3`$, the cooling break frequency increases with time, unlike the decreasing behavior it has for $`T<T_r`$.
Case 2: $`p>3`$. This case is treated here for completeness, as there are no afterglows for which such a steep electron index has been found. Equations (19), (30), and (40) lead to
$$Y_a=\frac{3}{3s}\epsilon _e^2n\mathrm{\Gamma }^2r_{18}.$$
(47)
For $`\epsilon _B<\stackrel{~}{\epsilon _B}`$ the Compton parameter is above unity until
$$(s=0)T_y=0.11E_{53}n_{,0}\epsilon _{e,1}^4\mathrm{day},$$
(48)
$$(s=2)T_y=0.87A_{}\epsilon _{e,1}^2\mathrm{day}.$$
(49)
For $`\epsilon _B<\stackrel{~}{\epsilon _B}`$ and $`T_r<T<T_y`$, the evolution of the cooling break frequency is
$$(s=0)\nu _c=1.1\times 10^{17}E_{53}^{3/2}n_{,0}^2\epsilon _{e,1}^4\epsilon _{B,3}^{3/2}T_d^{1/2}\mathrm{Hz},$$
(50)
$$(s=2)\nu _c=1.5\times 10^{16}E_{53}^{1/2}A_{}^4\epsilon _{e,1}^4\epsilon _{B,3}^{3/2}T_d^{5/2}\mathrm{Hz}.$$
(51)
Note that in this regime $`\nu _c`$ increases with time in both models.
### 3.3 Absorption Break
The synchrotron self-absorption frequency $`\nu _a`$ can be calculated with the aid of equation (6.50) from Rybicki & Lightman (1979). With the notations $`\gamma _p=\mathrm{min}(\gamma _i,\gamma _c)`$, $`\nu _p=\mathrm{min}(\nu _i,\nu _c)`$, and $`\nu _0=\mathrm{max}(\nu _i,\nu _c)`$ it can be shown that optical thickness to synchrotron self-absorption can be approximated by
$$\tau _{ab}(\nu )5\frac{e\mathrm{\Sigma }}{B\gamma _p^5}\times \{\begin{array}{cc}(\nu /\nu _p)^{5/3}\hfill & \nu <\nu _p\hfill \\ (\nu /\nu _p)^{(q+4)/2}\hfill & \nu _p<\nu <\nu _0\hfill \end{array},$$
(52)
where $`\mathrm{\Sigma }=(3s)^1nr`$ is the remnant electron column density, $`q=2`$ for radiative electrons ($`\gamma _c<\gamma _i`$), $`q=p`$ for adiabatic electrons ($`\gamma _i<\gamma _c`$).
#### 3.3.1 Radiative Electrons
Equations (20), (26), and (52) give the optical thickness at the cooling break frequency
$$\tau _c=\frac{5}{3s}\frac{enr}{B\gamma _c^5}=\frac{1.1\times 10^3}{(3s)(5s)^5}(Y_r+1)^5\epsilon _B^{9/2}n^{11/2}\mathrm{\Gamma }^4r_{18}^6.$$
(53)
For $`s=0`$ equations (14), (15), and (53) lead to
$$(s=0)\tau _c=0.11E_{53}^2n_{,0}^{7/2}(Y_r+1)^5\epsilon _B^{9/2}.$$
(54)
For $`\tau _c<1`$ the optical thickness to synchrotron self-absorption is unity at $`\nu _a`$ given by $`\nu _a=\nu _c\tau _c^{3/5}`$:
$$(s=0)\nu _a=6.5\times 10^9E_{53}^{7/10}n_{,0}^{11/10}(Y_r+1)\epsilon _{B,1}^{6/5}T_d^{1/2}\mathrm{Hz}.$$
(55)
For $`s=2`$ equations (16) – (18), and (53) give
$$(s=2)\tau _c=0.44E_{53}^{3/2}A_{}^7(Y_r+1)^5\epsilon _B^{9/2}T_d^{7/2}.$$
(56)
For $`T>T_a`$ we have $`\tau _c<1`$ and $`\nu _a<\nu _c`$:
$$(s=2)\nu _a=1.4\times 10^{12}E_{53}^{2/5}A_{}^{11/5}(Y_r+1)\epsilon _{B,2}^{6/5}T_{d,2}^{8/5}\mathrm{Hz}.$$
(57)
#### 3.3.2 Adiabatic Electrons
The optical thickness to synchrotron self-absorption at the injection break can be found using equations (19), (20), and (52):
$$\tau _i=\frac{5}{3s}\frac{enr}{B\gamma _i^5}=\frac{2.9\times 10^2}{3s}\frac{n^{1/2}r_{18}}{\epsilon _{e,1}^5\epsilon _B^{1/2}\mathrm{\Gamma }^6}.$$
(58)
If the external medium is homogeneous
$$(s=0)\tau _i=1.3\times 10^6E_{53}^{1/2}n_{,0}\epsilon _{e,1}^5\epsilon _{B,2}^{1/2}T_d^{5/2},$$
(59)
$$(s=0)\nu _a=2.6\times 10^9E_{53}^{1/5}n_{,0}^{3/5}\epsilon _{e,1}^1\epsilon _{B,2}^{1/5}\mathrm{Hz},$$
(60)
which is time-independent.
For the wind model
$$(s=2)\tau _i=6.8\times 10^7E_{53}^{3/2}A_{}^2\epsilon _{e,1}^5\epsilon _{B,2}^{1/2}T_d^{3/2},$$
(61)
$$(s=2)\nu _a=3.7\times 10^9E_{53}^{2/5}A_{}^{6/5}\epsilon _{e,1}^1\epsilon _{B,2}^{1/5}T_d^{3/5}\mathrm{Hz}.$$
(62)
Note that, in general, the absorption frequency decreases faster for a remnant interacting with a wind than for one running into a homogeneous external medium.
## 4 Analytical Light-Curves
If the effects arising from the remnant spherical shape (see Appendix) are ignored, than the observed flux peaks at $`\nu _p=\mathrm{min}(\nu _i,\nu _c)`$, where it has a value
$$F_{\nu _p}\frac{\sqrt{3}\varphi _p}{4\pi D^2}\frac{e^3}{m_ec^2}\mathrm{\Gamma }BN_e.$$
(63)
Here $`\varphi _p`$ is a factor calculated by Wijers & Galama (1999), which we shall set $`\varphi _p=0.63`$, $`D=(1+z)^{1/2}D_l(z)`$ with $`D_l`$ the luminosity distance, and $`N_e=m(r)/m_p`$ is the number of electrons in the remnant. Equations (19), (20), and (63) give
$$F_{\nu _p}=\frac{57}{3s}D_{28}^2\epsilon _B^{1/2}\mathrm{\Gamma }^2n^{3/2}r_{18}^3\mathrm{mJy}.$$
(64)
The afterglow emission at any given frequency and time can be calculated using the synchrotron spectrum for the electron distributions given in equations (31) and (32) (e.g. Sari et al. 1998):
$$F_\nu =F_{\nu _c}\{\begin{array}{ccc}(\nu /\nu _a)^2(\nu _a/\nu _c)^{1/3}\hfill & \nu <\nu _a\hfill & (1)\hfill \\ (\nu /\nu _c)^{1/3}\hfill & \nu _a<\nu <\nu _c\hfill & (2)\hfill \\ (\nu /\nu _c)^{1/2}\hfill & \nu _c<\nu <\nu _i\hfill & (3)\hfill \\ (\nu /\nu _i)^{p/2}(\nu _c/\nu _i)^{1/2}\hfill & \nu _i<\nu \hfill & (4)\hfill \end{array}$$
(65)
for $`T<T_r`$, assuming that $`\nu _a<\nu _c`$, and
$$F_\nu =F_{\nu _i}\{\begin{array}{ccc}(\nu /\nu _a)^2(\nu _a/\nu _i)^{1/3}\hfill & \nu <\nu _a\hfill & (5)\hfill \\ (\nu /\nu _i)^{1/3}\hfill & \nu _a<\nu <\nu _i\hfill & (6)\hfill \\ (\nu /\nu _i)^{(p1)/2}\hfill & \nu _i<\nu <\nu _c\hfill & (7)\hfill \\ (\nu /\nu _c)^{p/2}(\nu _i/\nu _c)^{(p1)/2}\hfill & \nu _c<\nu \hfill & (8)\hfill \end{array}$$
(66)
for $`T>T_r`$, assuming that $`\nu _a<\nu _i`$.
In Figures 1 and 2 the plane $`T\epsilon _B`$ is divided into several regions which are labeled as in equations (65) – (66), according to the ordering of the observing frequency $`\nu `$ and of the three break frequencies $`\nu _a`$, $`\nu _i`$, and $`\nu _c`$. The observed fluxes in each case are given in the Appendix, and a set of multi-wavelength light-curves is shown in Figure 3. The results shown in Figures 1–4 have been obtained using equations which are valid in any relativistic regime, such as equations (4), (19), (20), and (63).
### 4.1 Inverse Compton Emission
The IC emission can be easily calculated by using the above equations for the synchrotron spectrum and Compton parameter. The up-scattered spectrum peaks at $`\nu _{ic}\gamma _c^2\nu _c`$ if electrons are radiative and at $`\nu _{ic}\gamma _i^2\nu _i`$ if electrons are adiabatic. It can be shown that, for any electron radiative regime, the flux of the up-scattered emission at this frequency is
$$F_{\nu _{ic}}^{(ic)}=\tau _eF_{\nu _p},$$
(67)
where $`\tau _e`$ and $`F_{\nu _p}`$ are given by equations (30) and (64), respectively. For the up-scattering of synchrotron photons above $`\nu _a`$ and assuming that $`Y_a<1`$, the resulting IC light-curves have the following behaviors:
$$\left(\begin{array}{c}s=0\hfill \\ T<T_r\hfill \end{array}\right)F_\nu ^{(ic)}\{\begin{array}{cc}T^{1/3}\hfill & \nu <\gamma _c^2\nu _c\hfill \\ T^{1/8}\hfill & \gamma _c^2\nu _c<\nu <\gamma _i^2\nu _i\hfill \\ T^{(9p10)/8}\hfill & \gamma _i^2\nu _i<\nu \hfill \end{array},$$
(68)
$$\left(\begin{array}{c}s=0\hfill \\ T>T_r\hfill \end{array}\right)F_\nu ^{(ic)}\{\begin{array}{cc}T^1\hfill & \nu <\gamma _i^2\nu _i\hfill \\ T^{(9p11)/8}\hfill & \gamma _i^2\nu _i<\nu <\gamma _c^2\nu _c\hfill \\ T^{(9p10)/8}\hfill & \gamma _c^2\nu _c<\nu \hfill \end{array},$$
(69)
$$\left(\begin{array}{c}s=2\hfill \\ T<T_r\hfill \end{array}\right)F_\nu ^{(ic)}\{\begin{array}{cc}T^{5/3}\hfill & \nu <\gamma _c^2\nu _c\hfill \\ T^0\hfill & \gamma _c^2\nu _c<\nu <\gamma _i^2\nu _i\hfill \\ T^{(p1)}\hfill & \gamma _i^2\nu _i<\nu \hfill \end{array},$$
(70)
$$\left(\begin{array}{c}s=2\hfill \\ T>T_r\hfill \end{array}\right)F_\nu ^{(ic)}\{\begin{array}{cc}T^{1/3}\hfill & \nu <\gamma _i^2\nu _i\hfill \\ T^p\hfill & \gamma _i^2\nu _i<\nu <\gamma _c^2\nu _c\hfill \\ T^{(p1)}\hfill & \gamma _c^2\nu _c<\nu \hfill \end{array}.$$
(71)
For external media that are not denser than assumed so far, the IC emission is weaker than synchrotron, even in soft X-rays. As shown in the upper left panel of Figure 3, for $`n_{}>10\mathrm{cm}^3`$ and $`A_{}>1`$ the up-scattered radiation can dominate the synchrotron emission at times $`T>10^1`$ day, diminishing the decay rate of the X-ray emission. This is due to that the flux at the IC peak (eq. ) depends strongly on the external medium density: $`F_{\nu _{ic}}^{(ic)}n_{}^{5/4}`$ for $`s=0`$ and $`F_{\nu _{ic}}^{(ic)}A_{}^{5/2}`$ for $`s=2`$. The afterglow flattening is strongly dependent on the observing frequency, being absent in the optical and below.
## 5 Conclusions
Using the analytical results given in equations (B4) – (C21), the afterglow light-curve can be calculated at any frequency and at observing times up to the onset of the non-relativistic phase. As illustrated in Figure 3, the largest differences between the afterglow emission in the two models for the external medium is seen at low frequencies (lower panels). However, the scintillation due to the local interstellar medium (Goodman 1997), may hamper the use of the radio light-curves to identify the type of external medium and geometry of the ejecta (Frail et al. 2000).
Figure 4 shows that, for various model parameters, the rate of change of the afterglow emission at $`\nu 10^{12}`$ Hz and at early times (when the jet effects are negligible, provided that the jet is initially wider than a few degrees) exhibits a strong dependence on the type of external medium. If the external medium is homogeneous the sub-millimeter afterglow should rise slowly at times between $`1`$ hour and $`1`$ day, while for a pre-ejected wind the emission should fall off steeply, followed by a plateau<sup>2</sup><sup>2</sup>2The light-curves presented in the upper left panel of Figure 4 show that this criterion for determining the type of external medium fails only if the particle density of the homogeneous medium exceeds $`10\mathrm{cm}^3`$. In this case the X-ray emission may help to distinguish between the two models of external medium, as the absence of a flattening of the high energy emission is compatible only with a pre-ejected wind with $`A_{}`$ less than a few.. Therefore observations made at sub-millimeter frequencies with the SCUBA (James Clerk Maxwell Telescope) or with the MAMBO (IRAM Telescope) instruments would be very powerful in determining if the medium which the remnant runs into is homogeneous or follows a $`r^2`$ law.
We note that if turbulence in the shocked fluid does not lead to a significant mixing, then the inhomogeneous electron distribution will alter the afterglow spectrum below the absorption frequency $`\nu _a`$ as described by Granot, Piran & Sari (2000). The result is that the afterglow emission at $`\nu <\nu _a`$ rises more slowly than calculated here. For instance, the $`T^2`$ rise exhibited by the $`\nu =10^{12}`$ Hz light-curves shown in Figure 4 for the wind model at early times becomes a $`T^1`$ rise. Nevertheless, the basic difference mentioned above between the temporal behaviors of the sub-millimeter light-curves at $`10^2\mathrm{day}<T<1\mathrm{day}`$ remains unchanged.
The IC losses alter the evolution of the cooling break $`\nu _c`$ if the electrons injected with minimal energy are adiabatic and if the Compton parameter is above unity (i.e. the magnetic field parameter is weaker than that given in eq. ). In the case of a homogeneous external medium, the cooling break frequency decreases as $`T^{1/2}`$ if the electrons are radiative. When the electrons become adiabatic, this break evolves as $`T^{\frac{3p8}{82p}}`$ for $`p<3`$ and increases as $`T^{1/2}`$ for $`p>3`$. For an external wind the change is from $`\nu _cT^{1/2}`$ to $`\nu _cT^{\frac{3p4}{82p}}`$ for $`p<3`$, and to $`\nu _cT^{5/2}`$ for $`p>3`$. Consequently the power-law decay of the afterglow emission at frequencies above the cooling break flattens by up to $`1/2`$ if the external medium is homogeneous and by up to $`1`$ if the medium is a wind. For an electron index $`p<3`$, the flattening is mild and likely to be seen only in the optical emission from a remnant interacting with a pre-ejected wind.
The IC emission itself is generally weaker than the synchrotron emission. Nevertheless, if the external medium is sufficiently dense (i.e. $`n_{}>10\mathrm{cm}^3`$ or $`A_{}>`$ few), a flattening of the soft X-ray light-curve should be seen few hours after the main event, at fluxes well above the threshold of BeppoSAX (see Figure 3, upper left panel). The flattening of the afterglow emission due to the up-scattered radiation is a chromatic feature, appearing only at high frequencies, and its strength is moderately dependent on the remaining model parameters.
Finally, another possible signature of the interaction with the wind ejected by a Wolf-Rayet star should be found during the GRB emission in the form of smooth pulses of increasing duration. Such pulses are generated in internal shocks when the decelerating outermost shell is hit from behind by shells ejected at later times. This phenomenon is more likely to be seen in the wind model, for which the deceleration radius is smaller than for a homogeneous medium.
AP acknowledges the support received from Princeton University through the Lyman Spitzer, Jr. Fellowship. We thank Bohdan Paczyński for useful discussions.
APPENDIX
The observer time calculated with equation (12) is the time when photons emitted by the shocked gas moving precisely toward the observer arrive at detector. Photons emitted by the fluid moving at a non-zero angle relative to the direction toward the observer are less boosted by the relativistic motion of the source and arrive at the observer later. The effect of the remnant spherical curvature on the afterglow emission can be reduced to a correction factor that must be applied to the light-curve obtained analytically using the remnant parameters at laboratory frame time $`t`$ related to the observer time $`T`$ through equation (12). This correction factor is dependent on the observing frequency, as described below.
## A Correction Factors for Analytical Light-Curves
The flux received by the observer at time $`T`$ can be calculated by integrating the remnant emission over the equal arrival-time surface (Panaitescu & Mészáros 1998). This surface is defined by $`cT=ctr\mathrm{cos}\theta `$, where $`\theta `$ is the azimuthal angle measured relative to the observer’s line of sight toward the remnant center. Using the equation for dynamics of relativistic remnants ($`\mathrm{\Gamma }r^{(3s)/2}`$), the integral can be written as
$$F_\nu (T)=\frac{A(cT)^2}{2(3s)D^2}_0^{R(T)}\frac{(P_\nu ^{}^{})_er^{4s}}{\mathrm{\Gamma }^3(x+1)^2}dr,$$
(A1)
where $`x(3s)(r/R)^{4s}`$, and $`R`$ is the radius for which photons emitted at $`\theta =0`$ arrive at $`T`$, i.e. the maximal radius on the equal arrival-time surface. In equation (A1) $`(P_\nu ^{}^{})_e`$ is the co-moving frame power per electron, taking into account the spectrum of the emission. For the synchrotron spectra given in equations (65) and (66), one can write generically $`(P_\nu ^{}^{})_e=\sqrt{3}\varphi _p(e^3B/m_ec^2)(\nu ^{}/\nu _a^{})^{\alpha _a}(\nu ^{}/\nu _i^{})^{\alpha _i}(\nu ^{}/\nu _c^{})^{\alpha _c}`$, where $`\nu ^{}`$ is the co-moving frame observing frequency. The integral in equation (A1) is determined by the the values at of the integrand at $`r<R`$.
If the effect of the remnant geometrical curvature on the photon arrival-time were ignored, then the observed flux would be
$$F_\nu ^{(0)}(R)=\frac{1}{4\pi D^2}\left[(P_\nu ^{}^{})_e\mathrm{\Gamma }N_e\right]_{R(T)},$$
(A2)
which can be calculated based on equations (63), (65) and (66). The $`r`$-dependent quantities in equation (A1) can be expressed in terms of their values at $`r=R`$, so that the flux $`F_\nu `$ can be written as $`F_\nu ^{(0)}`$ times a correction factor
$$K=2(4s)^{2(\alpha _a+\alpha _i+\alpha _c)}_0^1u^{f(s)}\left[1+(3s)u^{4s}\right]^{\alpha _a+\alpha _i+\alpha _c2}du,$$
(A3)
where $`f(s)=72.5s2\alpha _a(q3s)+\alpha _i(20.5s)\alpha _c(2+0.5s)`$, with $`q=1,2`$ for radiative, radiative electrons, respectively. The $`K`$ factor depends on $`\nu `$ and also on $`p`$ if $`\nu >\nu _i`$. Table 1 gives its values for various cases.
## B Homogeneous External Medium ($`s=0`$)
By substituting the equations for the break frequencies and equation (64) in equations (65) and (66), and taking into account the above correction factors for the remnant curvature, the following fluxes are obtained
$$F_\nu \stackrel{(1)}{=}0.3D_{28}^2(Y_r+1)^1n_{,0}^1\epsilon _{B,1}^1\nu _{9.7}^2T_{d,1}\mathrm{mJy},$$
(B4)
$$F_\nu \stackrel{(2)}{=}10D_{28}^2(Y_r+1)^{2/3}E_{53}^{7/6}n_{,0}^{5/6}\epsilon _{B,2}\nu _{14.6}^{1/3}T_{d,2}^{1/6}\mathrm{mJy},$$
(B5)
$$F_\nu \stackrel{(3)}{=}40D_{28}^2(Y_r+1)^1E_{53}^{3/4}\epsilon _{B,1}^{1/4}\nu _{14.6}^{1/2}T_{d,2}^{1/4}\mathrm{mJy},$$
(B6)
$$F_\nu \stackrel{(4)}{=}10^{2.10.6p}D_{28}^2(Y_r+1)^1E_{53}^{\frac{p+2}{4}}\epsilon _{e,1}^{p1}\epsilon _{B,1}^{\frac{p2}{4}}\nu _{14.6}^{\frac{p}{2}}T_d^{\frac{3p2}{4}}\mathrm{mJy},$$
(B7)
$$F_\nu \stackrel{(5)}{=}30D_{28}^2E_{53}^{1/2}n_{,0}^{1/2}\epsilon _{e,1}\nu _{9.7}^2T_{d,1}^{1/2}\mathrm{mJy},$$
(B8)
$$F_\nu \stackrel{(6)}{=}1D_{28}^2E_{53}^{5/6}n_{,0}^{1/2}\epsilon _{e,1}^{2/3}\epsilon _{B,4}^{1/3}\nu _{14.6}^{1/3}T_{d,2}^{1/2}\mathrm{mJy},$$
(B9)
$$F_\nu \stackrel{(7)}{=}10^{2.11.3p}D_{28}^2E_{53}^{\frac{p+3}{4}}n_{,0}^{1/2}\epsilon _{e,1}^{p1}\epsilon _{B,4}^{\frac{p+1}{4}}\nu _{14.6}^{\frac{p1}{2}}T_d^{\frac{3}{4}(p1)}\mathrm{mJy},$$
(B10)
$$F_\nu \stackrel{(8)}{=}10^{2.40.8p}D_{28}^2E_{53}^{\frac{p+2}{4}}\epsilon _{e,1}^{p1}\epsilon _{B,2}^{\frac{p2}{4}}\nu _{14.6}^{\frac{p}{2}}T_d^{\frac{3p2}{4}}\mathrm{mJy},$$
(B11)
$$F_\nu \stackrel{(8a)}{=}10^{\frac{2p^27.7p+0.8}{4p}}D_{28}^2\left[E_{53}^{\frac{1}{4}(12p^2)}n_{,0}^{\frac{1}{2}(p2)}\epsilon _{e,1}^{(p1)(3p)}\epsilon _{B,4}^{\frac{1}{4}(p^2+2p+4)}\right]^{\frac{1}{4p}}\nu _{17.5}^{\frac{p}{2}}T_{d,1}^{\frac{3p}{4}+\frac{1}{4p}}\mathrm{mJy}(2<p<3).$$
(B12)
The case given in equation (B12) and labeled $`(8a)`$ corresponds to the same frequency ordering as for case $`(8)`$, but the cooling break $`\nu _c`$ evolution (eq. ) is determined by the IC losses, i.e. $`T_r<T<T_y`$ and $`Y_a>1`$.
## C Wind External Medium ($`s=2`$)
Following the same exercise as above and using the relevant equations, the following results can be obtained for the wind model:
$$F_\nu \stackrel{(1)}{=}0.03D_{28}^2(Y_r+1)^1E_{53}A_{}^2\epsilon _{e,1}^1\epsilon _{B,2}^1\nu _{9.7}^2T_{d,1}^2\mathrm{mJy},$$
(C13)
$$F_\nu \stackrel{(2)}{=}70D_{28}^2(Y_r+1)^{2/3}E_{53}^{1/3}A_{}^{5/3}\epsilon _{B,2}\nu _{12}^{1/3}T_{d,1}^{2/3}\mathrm{mJy},$$
(C14)
$$F_\nu \stackrel{(5)}{=}0.07D_{28}^2E_{53}A_{}^1\epsilon _{e,1}\nu _{9.7}^2T_{d,1}^1\mathrm{mJy},$$
(C15)
$$F_\nu \stackrel{(6)}{=}9D_{28}^2E_{53}^{1/3}A_{}\epsilon _{e,1}^{2/3}\epsilon _{B,3}^{1/3}\nu _{12}^{1/3}T_d^0\mathrm{mJy},$$
(C16)
$$F_\nu \stackrel{(7)}{=}10^{2.31.2p}D_{28}^2E_{53}^{\frac{p+1}{4}}A_{}\epsilon _{e,1}^{p1}\epsilon _{B,4}^{\frac{p+1}{4}}\nu _{14.6}^{\frac{p1}{2}}T_d^{\frac{3p1}{4}}\mathrm{mJy},$$
(C17)
$$F_\nu \stackrel{(8a)}{=}10^{\frac{1.9p^28.6p+5.4}{4p}}D_{28}^2E_{53}^{\frac{p+2}{4}}\left[A_{}^{(p2)}\epsilon _{e,1}^{(p1)(3p)}\epsilon _{B,4}^{\frac{1}{4}(p^2+2p+4)}\right]^{\frac{1}{4p}}\nu _{17.5}^{\frac{p}{2}}T_{d,1}^{\frac{3p}{4}+\frac{p}{2(4p)}}\mathrm{mJy}(2<p<3).$$
(C18)
For the remainder of the cases it can be shown that the fluxes for the wind model differ from those for obtained in the $`s=0`$ model only by a constant factor:
$$\frac{F_\nu (s=2)}{F_\nu (s=0)}\stackrel{(3)}{=}0.60\frac{Y_r(s=0)}{Y_r(s=2)},$$
(C19)
$$\frac{F_\nu (s=2)}{F_\nu (s=0)}\stackrel{(4)}{=}1.24\left(\frac{17}{72}\right)^{p/4}\frac{Y_r(s=0)}{Y_r(s=2)},$$
(C20)
$$\frac{F_\nu (s=2)}{F_\nu (s=0)}\stackrel{(8)}{=}1.35\left(\frac{17}{72}\right)^{p/4}.$$
(C21)
This means that observations made at frequencies above the cooling break are unable to distinguish between the two models for the external medium, with the exception of the case where the electron cooling is dominated by IC losses (case $`8a`$).
The reference frequency in the equations above was chosen $`\nu =4\times 10^{14}`$ Hz whenever this optical frequency falls in one of the cases labeled in equations (65)–(66), for the model parameters given in Figures 1 and 2. The observer time was scaled to a value at which each case is more likely to occur.
The above light-curves must be corrected to account for cosmological effects. This is achieved by replacing $`T`$ with $`T/(1+z)`$ and $`\nu `$ with $`(1+z)\nu `$, where $`z`$ is the burst redshift.
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# 1 Introduction
## 1 Introduction
In the last years we have seen increasing evidences that string/M theory on AdS spaces are dual to large $`N`$ strongly coupled gauge theories . The most extensively studied cases are the dualities between Type IIB string theories on $`\mathrm{𝐀𝐝𝐒}_\mathrm{𝟓}\times M_5`$ for positively curved five dimensional Einstein manifolds and large $`N`$ strongly coupled four dimensional conformal gauge theories. The simplest example is $`\mathrm{𝐀𝐝𝐒}_\mathrm{𝟓}\times 𝐒^\mathrm{𝟓}`$. In this case, the dual field theory is $`𝒩=4`$ supersymmetric gauge theory. In , field theories with less supersymmetry have been studied as duals to string theory on orbifolds of $`𝐒^\mathrm{𝟓}`$. In , Klebanov and Witten studied the Einstein manifold $`𝐓^{1,1}=(SU(2)\times SU(2))/U(1)`$. This was the first example of five dimensional spaces which are not orbifolds of $`𝐒^\mathrm{𝟓}`$. Type IIB string theory compactified on this manifold is dual to an $`𝒩=1`$ superconformal $`SU(N)\times SU(N)`$ gauge theory with a quartic superpotential for bifundamental fields. In the T-dual picture, D3 branes probing a metric cone over $`𝐓^{1,1}`$ (which is the conifold) is either a brane configuration with D4 branes together with orthogonal NS branes or a brane box with D5 branes together with orthogonal NS branes . Other results on the conifold or the quotients of the conifold and their field theory duals were obtained in .
Another exciting development was study of the gravity dual of the field theory Renormalization group flow. The main point is that the radial coordinate of $`\mathrm{𝐀𝐝𝐒}_5`$ has the natural interpretation as an energy scale in field theory. Thus it becomes natural to consider type IIB supergravity interpolating solutions where the metric and the fields depend on the radial coordinate, and to interpret these solutions as RG flows in the dual field theory. Many ideas have emerged concerning different aspects of supersymmetric and non-supersymmetric RG flows for four dimensional theories .
Because the AdS/CFT involves the full string theory we should be able to go beyond the Supergravity approximation. In Klebanov and Nekrasov have studied the gravity duals of fractional branes in supersymmetric conifold and orbifold theories where the presence of the fractional branes breaks the conformal invariance and introduces an RG flow in field theories. One of the models considered in involves a large number of D3-branes at a conifold singularity whose near-horizon is a $`\mathrm{𝐀𝐝𝐒}_5\times 𝐓^{1,1}`$ background and the fluxes of $`B^{NS}`$ and $`B^{RR}`$ forms through the blow-up 2-cycle determine a difference in the coupling constants of the two group factors appearing on the world-volume of the D3 branes at singularities. The supergravity equations of motion together with the specific formulas for the 2-forms and 3-forms on $`𝐓^{1,1}`$ give a solution which reproduces the logarithmic flow of field theory beta function. Previous results were obtained in type 0B string where effective action uncertainties occur . In all the studies of RG flow from AdS/CFT, the RG flow was determined by turning on different operators in the field theory which break the conformal invariance. In supergravity this means turning on some of the supergravity scalar fields. Another way to break the conformal invariance is to introduce the twisted sectors of string theories.
In this paper we go one step further and study D3 branes on an orbifolded conifold which is the quotient of the conifold by $`𝐙_k\times 𝐙_l`$. Now the horizon $`X_5`$ will be singular along two disjoint, but linked circles and we need to resolve the singularities in order to obtain a smooth Einstein manifold $`\stackrel{~}{X_5}`$. We completely describe the resolution of the orbifolded conifold itself in two steps and discover that there are $`kl`$ isolated conifold singularities after the first step of the resolution. After the resolution we find a smooth Einstein manifold $`\stackrel{~}{X_5}`$ containing $`kl+k+k2`$ two-cycles. As $`\stackrel{~}{X_5}`$ approaches the exceptional fiber of the first resolution, $`kl`$ cycles are vanishing into the singular points, and $`k+l2`$ cycles deform to cycles in the fiber which separate the two circles of singularities. Near each singular point, $`\stackrel{~}{X_5}`$ can be approximated by $`𝐓^{1,1}`$ and $`kl`$ two-cycles of $`\stackrel{~}{X_5}`$ come from these $`𝐓^{1,1}`$’s. We then consider a large number of D3 branes probing this singularity, which corresponds to a brane box with D5 branes and orthogonal NS branes via T-duality . The D5 branes wrapped on 2-cycles of $`\stackrel{~}{X_5}`$ vanishing into the singular points of the partially resolved orbifolded conifold are the fractional D3 branes. We study the $`B^{NS}`$ and $`B^{RR}`$ fluxes through different 2-cycles of $`\stackrel{~}{X_5}`$ which give rise to a logarithmic flow for the field theory coupling constants. This agrees with the field theory expectations for the RG flow.
In section 2 we study the geometry of the orbifolded conifold and describe how to obtain a smooth horizon from the singular horizon $`𝐓^{1,1}/𝐙_k\times 𝐙_l`$. We also identify the different fractional D3 branes in the singularity picture with different components of brane interval or brane box configurations obtained by T-dualities. In section 3 we describe the supergravity dual to the field theory Renormalization Group flows.
## 2 Geometry and Brane Configurations of Orbifolded Conifolds
In this section we study the geometry of the orbifolded conifolds $`𝒞_{kl}`$ in detail and we make connections with brane configurations obtained by T-dualities. In particular, we study the resolutions of the orbifolded conifolds $`𝒞_{kl}`$ and the associated fractional branes in the T-dual picture. At the end of the section, we describe the homological cycles of the resolved horizon of $`𝒞_{kl}`$, and thus extends the results of . This play an important role in the study of fluxes in the next section.
Consider a singular Calabi-Yau threefold $`Y_6`$ which is a metric cone over a five dimensional Einstein manifold $`M_5`$. Then the metric near the apex of the cone will be
$$ds_{Y_6}^2=dr^2+r^2ds_{M_5}^2.$$
(2.1)
Here the apex is located at $`r=0`$ and $`M_5`$ is called the horizon of the cone $`Y_6`$. If $`N`$ parallel D3 branes are placed at the apex of the cone $`Y_6`$, the resulting ten dimensional spacetime has the metric
$$ds^2=R^2\left[\frac{r^2}{R^4}(dt^2+dx_1^2+dx_2^2+dx_3^2)+\frac{dr^2}{r^2}+ds_5^2\right],R^4g_sN(\alpha ^{})^2.$$
(2.2)
The near-horizon ($`r0`$) limit of the geometry is $`\mathrm{𝐀𝐝𝐒}_\mathrm{𝟓}\times M_5`$. Type IIB theory on this background is conjectured to be dual to the conformal limit of the field theory on the D3 branes.
In , an example of such duality has been discussed in the case of a conifold whose horizon is $`𝐓^{1,1}=(SU(2)\times SU(2))/U(1)`$. The conifold is a three dimensional hypersurface singularity in $`𝐂^4`$ defined by:
$`𝒞:z_1z_2z_3z_4=0.`$ (2.3)
The conifold can be realized as a holomorphic quotient of $`𝐂^4`$ by the $`𝐂^{}`$ action given by
$`(A_1,A_2,B_1,B_2)(\lambda A_1,\lambda A_2,\lambda ^1B_1,\lambda ^1B_2)\text{ for }\lambda 𝐂^{}.`$ (2.4)
Thus the charge matrix is the transpose of $`Q^{^{}}=(1,1,1,1)`$ and $`\mathrm{\Delta }=\sigma `$ will be a convex polyhedral cone in $`𝐍_𝐑^{^{}}=𝐑^3`$ generated by $`v_1,v_2,v_3,v_4𝐍^{^{}}=𝐙^3`$ where
$`v_1=(1,0,0),v_2=(0,1,0),v_3=(0,0,1),v_4=(1,1,1).`$ (2.5)
The isomorphism between the conifold $`𝒞`$ and the holomorphic quotient is given by
$`z_1=A_1B_1,z_2=A_2B_2,z_3=A_1B_2,z_4=A_2B_1.`$ (2.6)
To identify the horizon from this point of view, note that we can divide by the scaling $`z_isz_i`$ (with real positive $`s`$) by setting $`|A_1|^2+|A_2|^2=|B_1|^2+|B_2|^2=1`$. This gives us $`𝐒^3\times 𝐒^3=SU(2)\times SU(2)`$. Then dividing by the $`U(1)`$ action
$`(A_1,A_2,B_1,B_2)(e^{i\alpha }A_1,e^{i\alpha }A_2,e^{i\alpha }B_1,e^{i\alpha }B_2),`$ (2.7)
we obtain $`𝐓^{1,1}=(SU(2)\times SU(2))/U(1)`$.
The five dimensional manifold $`𝐓^{1,1}`$ has 2-cycles and 3-cycles. Besides the D3 branes orthogonal to $`𝐓^{1,1}`$, there are wrapped D3 branes over the 3-cycles of $`𝐓^{1,1}`$ (which correspond to “dibaryon” operators ) and wrapped D5 branes over 2-cycles of $`𝐓^{1,1}`$ (which to correspond to domain walls in $`\mathrm{𝐀𝐝𝐒}_\mathrm{𝟓}`$ and to fractional D3 branes ). Because $`𝐓^{1,1}`$ is $`𝐒^\mathrm{𝟐}\times 𝐒^\mathrm{𝟑}`$, we can identify the 2-cycle with $`𝐒^\mathrm{𝟐}`$ and the 3-cycle with $`𝐒^\mathrm{𝟑}`$. These two cycles are orthogonal so the D3 brane wrapped on $`𝐒^\mathrm{𝟑}`$ is orthogonal to the D5 brane wrapped on $`𝐒^\mathrm{𝟐}`$, therefore when they cross each other a fundamental string is created as explained in and the gauge group becomes $`SU(N+1)\times SU(N)`$.
The geometrical picture is T-dual to different types of brane configuration. By one T-duality one can obtain the brane interval picture with D4 branes wrapped on a circle and by two T-dualities one obtains the brane box picture with D5 branes wrapped on a 2-torus. The fractional branes have also been identified in the brane interval picture in . The idea was to interpret the conifold (2.3) as a $`𝐂^{}`$ fibration over the $`𝐂^2`$ parameterized by $`z_3,z_4`$. By performing T-duality along the $`U(1)`$-orbit in the $`𝐂^{}`$-fiber, we obtain from the degenerate fibers $`z_1=0`$ and $`z_2=0`$, two NS fivebranes extended, say, $`x^0x^1x^2x^3x^4x^5`$ and $`x^0x^1x^2x^3x^8x^9`$ directions which we denote by NS and NS’ branes. The D3 branes located at the singular point transform into D4 branes wrapping a circle which is transverse to the NS fivebranes. $`𝐓^{1,1}`$ has a $`U(1)`$-fibration over $`𝐏^1\times 𝐏^1`$ and a two cycle $`𝐒^\mathrm{𝟐}`$ of $`𝐓^{1,1}`$ can be identified to the difference of two homologically distinct spheres coming from $`𝐏^1\times 𝐏^1`$. After identifying $`𝐏^1\times 𝐏^1`$ with the exceptional locus in the full resolution of the conifold, D5 brane wrapping the two cycle $`𝐒^\mathrm{𝟐}`$ will transform as a D4 brane wrapping on one interval between two NS-branes. This is a fractional brane in the interval model (Figure 1).
One of the goals of this paper is to study the fractional branes in the brane box model for a quotient of the conifold. To do this, we start by taking a further quotient of the conifold $`𝒞`$ by a discrete group $`𝐙_k\times 𝐙_l`$. Here $`𝐙_k`$ acts on $`A_i,B_j`$ by
$`(A_1,A_2,B_1,B_2)(e^{2\pi i/k}A_1,A_2,e^{2\pi i/k}B_1,B_2),`$ (2.8)
and $`𝐙_l`$ acts by
$`(A_1,A_2,B_1,B_2)(e^{2\pi i/l}A_1,A_2,B_1,e^{2\pi i/l}B_2).`$ (2.9)
Thus they will act on the conifold $`𝒞`$ by
$`(z_1,z_2,z_3,z_4)(z_1,z_2,e^{2\pi i/k}z_3,e^{2\pi i/k}z_4)`$ (2.10)
and
$`(z_1,z_2,z_3,z_4)(e^{2\pi i/l}z_1,e^{2\pi i/l}z_2,z_3,z_4).`$ (2.11)
Its quotient is called the orbifolded conifold (or the hyper-quotient of the conifold) and denoted by $`𝒞_{kl}`$.
Note that the action (2.8) leaves a complex two space $`A_1=B_1=0`$ fixed in $`𝐂^4`$ and this is isomorphic to $`𝐂^1`$ given by $`z_1=z_3=z_4=0`$ on the conifold $`𝒞`$ after dividing by the $`U(1)`$ action. Similarly, the action (2.9) leaves fixed a complex two space $`A_1=B_2=0`$ in $`𝐂^4`$ and it is isomorphic to $`𝐂^1`$ given by $`z_1=z_2=z_3=0`$ on the conifold $`𝒞`$ after dividing by the $`U(1)`$ action. Furthermore, the action (2.8) descends to the horizon $`𝐓^{1,1}`$ and leaves the following circle fixed:
$`|z_2|^2=1,z_1=z_3=z_4=0`$ (2.12)
or equivalently $`|A_2|^2=|B_2|^2=1,A_1=B_1=0(modU(1))`$. Similarly, the action (2.9) leaves the following circle fixed
$`|z_4|^2=1,z_1=z_2=z_3=0`$ (2.13)
or equivalently $`|A_2|^2=|B_1|^2=1,A_1=B_2=0(modU(1))`$. Hence the horizon $`X_5:=𝐓^{1,1}/𝐙_k\times 𝐙_l`$ of the orbifolded conifold is singular along these two circles. These two circles are separated but linked. The horizon $`X_5`$ has $`𝐀_{k1}`$ singularity along the circle (2.12) and $`𝐀_{l1}`$ singularity along the circle (2.13). String theory in the back ground $`\mathrm{𝐀𝐝𝐒}_5\times X_5`$ has massless fields which are localized along these two linked circles. As discussed in , these massless fields are the twisted modes and they propagate on $`\mathrm{𝐀𝐝𝐒}_5\times 𝐒^1\mathrm{𝐀𝐝𝐒}_5\times 𝐒^1`$ where $`𝐒^1𝐒^1`$ are the circles of singularities. As we shall see below, there are $`k+l2`$ 2-cycles which separate these two circles. The fluxes of the NSNS and RR two forms through these cycles give rise to scalars which live in the $`\mathrm{𝐀𝐝𝐒}_5\times 𝐒^1`$ space and are the same as the scalars introduced in section 3 of .
To put the actions (2.6), (2.8) and (2.9) on an equal footing, consider the over-lattice $`𝐍=𝐍^{^{}}+\frac{1}{k}(v_3v_1)+\frac{1}{l}(v_4v_1)`$. Now the lattice points $`\sigma 𝐍`$ of $`\sigma `$ in $`𝐍`$ are generated by $`(k+1)(l+1)`$ lattice points as a semigroup. The discrete group $`𝐙_k\times 𝐙_l𝐍/𝐍^{^{}}`$ will act on the conifold $`𝐂^4//U(1)`$ and its quotient will be the symplectic reduction $`𝐂^{(k+1)(l+1)}//U(1)^{(k+1)(l+1)3}`$. The new toric diagram for $`𝒞_{kl}`$ will also lie on a plane at a distance from the origin and the toric diagram on the plane for $`𝒞_{23}`$ is shown in Figure 1.
In suitable coordinates, the orbifolded conifold will be given by
$`𝒞_{kl}:xy=z^l,uv=z^k.`$ (2.14)
As we have seen above, the horizon $`X_5`$ is singular. To obtain a smooth Einstein manifold from $`X_5`$, we will resolve the singularities of $`𝒞_{kl}`$ itself. We resolve the singular threefold $`𝒞_{kl}`$ in two steps. In the first step, we choose a partial resolution, denoted by $`\stackrel{~}{𝒞_{kl}}`$, of the orbifolded conifold $`𝒞_{kl}`$ for which the horizon will be smooth, but the Calabi-Yau threefold $`\stackrel{~}{𝒞_{kl}}`$ will have $`kl`$ number of isolated singular points. Around each singular point, the Calabi-Yau space $`\stackrel{~}{𝒞_{kl}}`$ is locally a metric cone over an Einstein manifold $`𝐓^{1,1}`$. In terms of the toric diagram, the partial resolution we have chosen is obtained by adding all possible vertical and horizontal arrows to the toric diagram of $`𝒞_{kl}`$. For example, the toric diagram for $`\stackrel{~}{𝒞_{23}}`$ is given as in Figure 3.
We are going to describe in detail each step but let us discuss first some features and make the connection of the result with the T-dual brane configurations. The partially resolved space $`\stackrel{~}{𝒞_{kl}}`$ is covered by $`kl`$ squares and each square in the toric diagram represents an ordinary conifold. Thus the metric near each singular point can be written locally as follows
$`ds^2={\displaystyle \frac{1}{9}}(d\varphi +\mathrm{cos}\theta _1d\varphi _1+\mathrm{cos}\theta _2d\varphi _2)^2+{\displaystyle \frac{1}{6}}{\displaystyle \underset{a=1}{\overset{2}{}}}(d\theta _a^2+\mathrm{sin}^2\theta _ad\varphi _a^2).`$ (2.15)
Note that $`𝒞_{kl}`$ can be regarded as a $`𝐂^{}\times 𝐂^{}`$ fibration over the $`z`$-plane via (2.14). By taking $`T`$-duality along $`U(1)\times U(1)`$ orbit in $`𝐂^{}\times 𝐂^{}`$, we will have two types of NS branes extended in, say, $`x^0x^1x^2x^3x^4x^5`$ and $`x^0x^1x^2x^3x^8x^9`$ directions, where $`x^4,x^8`$ are compact directions coming from the degenerate $`U(1)\times U(1)`$-orbit. The separation of the NS branes along the $`x^8`$ direction and similarly that of the NS’ branes along the $`x^4`$ direction can be achieved by partially resolving $`𝒞_{kl}`$. But there will be $`kl`$ intersections of the NS and NS’ branes on the $`x^4x^8`$ torus which correspond to the singular points of the partially resolved orbifolded conifold $`\stackrel{~}{𝒞_{kl}}`$. By replacing these singular points by $`𝐏^1`$’s, we resolve the singularities of $`\stackrel{~}{𝒞_{kl}}`$. In the field theory, the process of blowing-up and turning the NS-NS B-fluxes through the $`kl`$ $`𝐏^1`$ cycles means turning on gauge couplings. In brane box configurations it means replacing the intersection of NS and NS’ branes with ‘diamonds’ , the size of the diamonds being given by the fluxes of $`NSNS`$ fields through the blow-up cycles. In this T-duality, D3 branes become D5 branes which fill the $`x^4x^8`$ directions and the resulting brane configuration is as a brane box model shown in Figure 4 . Its components are diamonds and boxes and a stripe is an horizontal or vertical line of boxes (in Figure 4 we have represented an horizontal stripe).
In the second step we completely resolve the singularities of $`\stackrel{~}{𝒞_{kl}}`$ by replacing each of the singular points by a copy of $`𝐏^1`$ as explained before, procedure called a small resolution. In terms of the toric diagram, this corresponds to joining a pair of the diagonal vertices (but not both pairs) by a line segment in each square. Let us denote this completely resolved threefold by $`\widehat{𝒞_{kl}}`$. For example, the toric diagram for $`\widehat{𝒞_{23}}`$ is given as in Figure 3.
Let
$`\stackrel{~}{\pi }:\stackrel{~}{𝒞_{kl}}𝒞_{kl},\widehat{\pi }:\widehat{𝒞_{kl}}\stackrel{~}{𝒞_{kl}}`$ (2.16)
be the first and the second resolution discussed above. Let $`𝐨`$ be the apex of the cone $`𝒞_{kl}`$ and
$`\stackrel{~}{X_5}=\stackrel{~}{\pi }^1(X_5).`$ (2.17)
Note that $`\stackrel{~}{𝒞_{kl}}`$ is covered by a $`kl`$ number of the ordinary conifolds corresponding to the squares in the toric diagram. Hence $`\stackrel{~}{𝒞_{kl}}`$ has $`kl`$ isolated singular points corresponding to the apexes of these ordinary conifolds. These singular points lie on the fiber $`\stackrel{~}{\pi }^1(𝐨)`$. Moreover, we will explicitly show that the exceptional fiber $`\stackrel{~}{\pi }^1(𝐨)`$ consists of $`(k1)(l1)`$ copies of $`𝐏^1\times 𝐏^1`$. The map $`\widehat{\pi }`$ modifies only the singular points of $`\stackrel{~}{𝒞_{kl}}`$ replacing each of them by a copy of the projective space $`𝐏^1`$. Thus $`\widehat{\pi }`$ is an isomorphism outside $`(\widehat{\pi }\stackrel{~}{\pi })^1(𝐨)`$. In particular, we have
$`\widehat{\pi }^1(\stackrel{~}{X_5})\stackrel{~}{X_5}.`$ (2.18)
and $`\stackrel{~}{X_5}`$ is smooth. As we mentioned above, $`\stackrel{~}{𝒞_{kl}}`$ is smooth outside the $`kl`$ ordinary conifold singular points. Thus if we put a large number of D3 branes at one of $`kl`$ isolated singular points, denoted by $`x_{ij},1ik,1jl`$, then the near-horizon limit of the geometry will be $`\mathrm{𝐀𝐝𝐒}_5\times 𝐓^{1,1}`$.
The 5 dimensional manifold $`\stackrel{~}{X_5}`$ can be regarded as a smoothing of the singular Einstein manifold $`X_5`$. As we will see later, there will be $`kl+k+l2`$ number of 2-cyles and 3-cycles in $`\stackrel{~}{X_5}`$ where $`kl`$ is the number of the cycle coming from the horizon of each singular point of $`\stackrel{~}{𝒞_{kl}}`$ and $`k+l2`$ number of them comes by separating the above discussed two fixed circles in $`X_5=𝐓^{1,1}/𝐙_k\times 𝐙_l`$. As mentioned above, the first kind of these cycles corresponds to the ‘diamonds’ and the second kind corresponds to the ‘stripes’ in the brane box model. They correspond to stripes of boxes instead of individual boxes because we need to consider curves of either $`𝐀_{k1}`$ or $`𝐀_{l1}`$ singularity.
Before starting the actual discussion concerning the identification of the different 2-cycles, we present another proof for the fact that the Einstein manifold $`𝐓^{1,1}`$ is homeomorphic to $`𝐒^3\times 𝐒^2`$ . By changing coordinates
$`z_1=w_1+iw_2,z_2=w_1iw_2,z_3=w_3+iw_4,z_4=w_3iw_4,`$ (2.19)
we can rewrite the conifold equation (2.3) as:
$`w_1^2+w_2^2+w_3^2+w_4^2=0.`$ (2.20)
Since the Einstein manifold $`𝐓^{1,1}`$ can be realized as a horizon (link) of the the conifold singularity, $`𝐓^{1,1}`$ is described by the intersection of (2.20) and the seven sphere in $`𝐂^4`$ given by
$`|w_1|^2+|w_2|^2+|w_3|^2+|w_4|^2=1.`$ (2.21)
¿From (2.20) and (2.21), we see that $`𝐓^{1,1}`$ is given by
$`x_1^2+x_2^2+x_3^2+x_4^2=y_1^2+y_2^2+y_3^2+y_4^2=1/2,`$
$`x_1y_1+x_2y_2+x_3y_3+x_4y_4=0,`$ (2.22)
where $`x_i`$ and $`y_i`$ are the real and imaginary parts of $`w_i`$. Thus the $`y_i`$’s describe a bundle of two spheres in the tangent bundle of $`𝐒^\mathrm{𝟑}`$ given by the coordinates $`x_i`$’s. Hence $`𝐓^{1,1}`$ is a sphere bundle $`𝐒^\mathrm{𝟐}`$ over $`𝐒^\mathrm{𝟑}`$. Since $`𝐒^\mathrm{𝟑}`$ is parallelizable , a sphere bundle $`𝐒^\mathrm{𝟐}`$ over $`𝐒^\mathrm{𝟑}`$ is trivial and $`𝐓^{1,1}`$ is diffeomorphic to $`𝐒^\mathrm{𝟐}\times 𝐒^\mathrm{𝟑}`$. In fact, the frame for the sphere bundle over $`𝐒^\mathrm{𝟑}`$ can be given by
$`\{(x_2,x_1,x_4,x_3),(x_3,x_4,x_1,x_2),(x_4,x_3,x_2,x_1)\}.`$ (2.23)
Next, we want to study the exceptional fiber $`\stackrel{~}{\pi }^1(𝐨)`$. We will illustrate the general situation with $`\stackrel{~}{𝒞_{23}}`$. To facilitate understanding, let us choose a basis for the lattice $`𝐍`$ so that the coordinates of the lattice points are as follows:
$`\begin{array}{cccc}v_1=(1,0,1),\hfill & v_2=(1,1,1),\hfill & v_3=(1,2,1),\hfill & v_4=(1,3,1),\hfill \\ v_5=(1,0,0),\hfill & v_6=(1,1,0),\hfill & v_7=(1,2,0),\hfill & v_8=(1,3,0),\hfill \\ v_9=(1,0,1),\hfill & v_{10}=(1,1,1),\hfill & v_{11}=(1,2,1).\hfill & v_{12}=(1,3,1).\hfill \end{array}`$ (2.25)
The Figure 6 shows the coordinate rings corresponding to the various squares of the toric diagram of $`\stackrel{~}{𝒞_{23}}`$.
By tedious but direct computations from the Figure 6, one can see that the fiber $`\stackrel{~}{\pi }^1(𝐨)`$ consists of a union of $`(k1)(l1)`$ numbers of $`𝐏^1\times 𝐏^1`$. Moreover the adjacent components of the fiber meet along $`𝐏^1`$. If we denote each $`𝐏^1\times 𝐏^1`$ by a vertex and we join two of vertices if they meet, then we get a lattice of size $`(k1)`$ by $`(l1)`$. The Figure 7 shows what they look like for $`\stackrel{~}{𝒞_{23}}`$.
In Figure 7, each square represents $`𝐏^1\times 𝐏^1`$ and the singular points of $`\stackrel{~}{𝒞_{23}}`$ are denoted by black dots. The fiber consists of two $`𝐏^1\times 𝐏^1`$ and they meet along $`𝐏^1`$. ¿From this picture, one can see that the second betti number $`h_2(\stackrel{~}{\pi }^1(𝐨))=k+l2`$. The D5 brane wrapping one of these $`k+l2`$ spheres corresponds the fractional D3 brane coming from the $`𝐙_k\times 𝐙_l`$ twisted sector of the type IIB supergravity on $`\mathrm{𝐀𝐝𝐒}_5\times 𝐓^{1,1}/𝐙_k\times 𝐙_l`$. In the brane box model, this type of a fractional D3 brane turns into a D5 brane living on a stripe (See Figure 4).
As we mentioned above, a full resolution $`\widehat{𝒞_{kl}}`$ of $`𝒞_{kl}`$ can be obtained by replacing the singular points of $`\stackrel{~}{𝒞_{kl}}`$ by copies of $`𝐏^1`$. Hence the exceptional fiber $`(\widehat{\pi }\stackrel{~}{\pi })^1(𝐨)`$ will acquire $`kl`$ copies of $`𝐏^1`$. This can be achieved by blowing up $`kl`$ points on $`\stackrel{~}{\pi ^1}(𝐨)`$. Thus the second betti number of the exceptional fiber $`(\widehat{\pi }\stackrel{~}{\pi })^1(𝐨)`$ will be $`k+l2+kl`$.
To study cycles on the Einstein manifold $`\stackrel{~}{X_5}`$, consider an inclusion map
$`i:\stackrel{~}{X_5}\widehat{𝒞_{kl}}.`$ (2.26)
First note that $`\widehat{𝒞_{kl}}`$ contacts to the exceptional fiber $`(\widehat{\pi }\stackrel{~}{\pi })^1(𝐨)`$, which we will denote by $`E`$. This contraction can be constructed by lifting a conical structure of $`𝒞_{kl}`$. We now regard $`\stackrel{~}{X_5}`$ as a boundary of a smooth 6 dimensional manifold $`\widehat{𝒞_{kl}}`$. Consider a long sequence of homology groups with $`𝐐`$ coefficients:
$`H_4(\widehat{𝒞_{kl}},\stackrel{~}{X_5})\stackrel{}{}H_3(\stackrel{~}{X_5})\stackrel{i_{}}{}H_3(\widehat{𝒞_{kl}})\stackrel{j_{}}{}H_3(\widehat{𝒞_{kl}},\stackrel{~}{X_5})\stackrel{}{}H_2(\stackrel{~}{X_5})\stackrel{i_{}}{}H_2(\widehat{𝒞_{kl}})`$ (2.27)
where $`j_{}`$ is induced by the inclusion $`j:\widehat{𝒞_{kl}}(\widehat{𝒞_{kl}},\stackrel{~}{X_5})`$. Via the universal coefficient theorem and the Poincaré duality, we have
$`H_3(\widehat{𝒞_{kl}},\stackrel{~}{X_5})H^3(\widehat{𝒞_{kl}})H_3(\widehat{𝒞_{kl}})=0`$ (2.28)
since $`\widehat{𝒞_{kl}}`$ can be deformed to $`E`$. Hence in the following commutative diagram, the top horizontal arrow will be injective and the bottom horizontal arrow will be surjective.
Here the vertical arrows are isomorphisms because of the universal coefficient theorem and the Poincaré dualities. Therefore, we conclude that
$`H_2(\stackrel{~}{X_5},𝐐)H_2(\widehat{𝒞_{kl}},𝐐).`$ (2.29)
Moreover, we want to see the origin of these two cycles. Let $`L_{ij}`$ be the horizon of $`\stackrel{~}{𝒞_{kl}}`$ at $`x_{ij}`$. Then $`L_{ij}`$ is isomorphic to $`𝐓^{1,1}`$ and $`L_{ij}`$ does not change under the resolution $`\widehat{\pi }:\widehat{𝒞_{kl}}\stackrel{~}{𝒞_{kl}}`$. From (2.29), we see that there is a natural inclusion
$`H_2(L_{ij})H_2(\widehat{𝒞_{kl}})H_2(\stackrel{~}{X_5}).`$ (2.30)
Thus we may regard the 2-cycle of $`H_2(L_{ij})`$ as a 2-cycle of $`H_2(\stackrel{~}{X_5})`$. We denote this 2-cycle by $`C_{ij}^2`$. By Poincaré duality, we may also regard the 3-cycle of $`H_3(L_{ij})`$ as a 3-cycle of $`H_3(\stackrel{~}{X_5})`$, which will be denoted by $`C_{ij}^3`$. Note that there are $`kl`$ contributions of 2-cycles from $`H_2(L_{ij})`$.
Moreover on each open neighborhood of $`x_{ij}`$ represented by a square, we can choose a basis for one-forms
$`e_{ij}^\psi =\frac{1}{3}\left(d\psi +\mathrm{cos}\theta _1\varphi _1+\mathrm{cos}\theta _2\varphi _2\right)`$ (2.31)
$`e_{ij}^{\theta _1}=\frac{1}{\sqrt{6}}d\theta _1e_{ij}^{\varphi _1}=\frac{1}{\sqrt{6}}\mathrm{sin}\theta _1d\varphi _1`$ (2.32)
$`e_{ij}^{\theta _2}=\frac{1}{\sqrt{6}}d\theta _2e_{ij}^{\varphi _2}=\frac{1}{\sqrt{6}}\mathrm{sin}\theta _2d\varphi _2,`$ (2.33)
so that the harmonic representatives of the second and third cohomology groups can be written as
$`e_{ij}^{\theta _1}e_{ij}^{\varphi _1}e_{ij}^{\theta _2}e_{ij}^{\varphi _2}`$ $`H^2(L_{ij})`$ (2.34)
$`e_{ij}^\psi e_{ij}^{\theta _1}e_{ij}^{\varphi _1}e_{ij}^\psi e_{ij}^{\theta _2}e_{ij}^{\varphi _2}`$ $`H^3(L_{ij}).`$ (2.35)
On the other hand, from the inclusion of $`\stackrel{~}{X_5}`$ into the partially resolved conifold $`\stackrel{~}{𝒞_{kl}}`$, we obtain a map
$`H_2(\stackrel{~}{X_5})H_2(\stackrel{~}{𝒞_{kl}}).`$ (2.36)
Since we obtain $`\stackrel{~}{𝒞_{kl}}`$ from $`\widehat{𝒞_{kl}}`$ by collapsing the blown-up two spheres which is a smooth deformation of the spheres in $`L_{ij}`$, we have the following exact sequence:
$`0{\displaystyle \underset{i=1,\mathrm{},k,j=1,\mathrm{},l}{}}H_2(L_{ij})H_2(\stackrel{~}{X_5})H_2(\stackrel{~}{𝒞_{kl}})0.`$ (2.37)
Therefore we have obtained a concrete description of the cycles of $`\stackrel{~}{X_5}`$ in terms of $`kl`$ cycles from $`H_2(L_{ij})`$ and the $`k+l2`$ cycles $`H_2(\stackrel{~}{𝒞_{kl}})`$. The cycles in $`H_2(\stackrel{~}{𝒞_{kl}})`$ are separating the singular points of $`\stackrel{~}{𝒞_{kl}}`$. hence the corresponding cycles in $`H_2(\stackrel{~}{X_5})`$ will separate the two fixed circles of $`X_5`$. This generalizes the results from the conifold and allows us to study the fluxes of NS-NS and R-R two forms in order to obtain logarithmic RG flow in the next section.
## 3 Fractional Branes and RG Flows
In this section we are going to extensively use the mathematical results of the previous section and identify the fractional branes as small perturbations of the string background. This will allow us to study the interpolation between the background with or without fractional D3 branes. This description will be shown to reproduce the logarithmic flow of gauge couplings,being in complete agreement with results of field theory.
We begin with a brief review of where the RG flow determined by the fractional D3 branes was considered. Their result is a particular example of our case for $`k=l=1`$. The coupling constant of field theory are written in terms of the two-form charges on the vanishing sphere of the singularity:
$$\tau =C_0+i\frac{1}{g^2}=_{C^2}B^{RR}+i_{C^2}B^{NS}$$
(3.1)
where $`B^{RR},B^{NS}`$ are the R-R and NS-NS 2-form potentials. At the conifold point the values of the B-fields are fixed and the coupling constant is $`g^2e^\varphi /2`$. By wrapping $`M`$ D5 branes over the 2-cycle of $`𝐓^{1,1}`$ in addition to $`N`$ regular D3 branes orthogonal to the conifold, the string background will contain $`M`$ units of RR 3-form flux through the 3-cycle of $`𝐓^{1,1}`$:
$$_{C^3}H^{RR}=M.$$
(3.2)
The equations of motion imply that $`H^{NS}`$ should be proportional to $`M`$ and to a product of the closed 2-form on $`𝐓^{1,1}`$ and a one form which involves $`dr`$ and taken to be $`\frac{df}{dr}dr=df(r)`$. Hence the two form potential $`B^{NS}`$ will be:
$$B^{NS}=e^\varphi f(r)\omega _2,\text{where}f(r)M\mathrm{log}r$$
(3.3)
where $`\omega _2`$ is the closed form on $`𝐓^{1,1}`$. The R-R scalar $`C_0=0`$ and the dilaton are set to have constant values.
The fractional D3 branes, obtained by wrapping D5 branes on the 2-cycle of $`𝐓^{1,1}`$ represent domain walls in $`\mathrm{𝐀𝐝𝐒}_5`$ and are obtained by wrapping D5 branes on the 2-cycle of $`𝐓^{1,1}`$. The relation between the two coupling constants of the field theory on the D3 branes writes in the presence of M D5 branes wrapped on the 2-cycle as
$$\frac{1}{g_1^2}\frac{1}{g_2^2}e^\varphi (M\mathrm{log}r\frac{1}{2})e^\varphi M\mathrm{log}r$$
(3.4)
where the last relation is true in the large M approximation. This agrees with the field theory logarithmic RG flow equation in a non-conformal theory, the conformality being broken by the presence of the fractional D3 branes.
We now proceed to the case of the orbifolded conifold. The field theory on the world-volume of the $`N`$ coincident D3 branes probing the singularity $`𝒞_{kl}`$ has been obtained in . It is an $`𝒩=1`$ chiral supersymmetric gauge theory with the gauge group
$`{\displaystyle \underset{i=1}{\overset{k}{}}}{\displaystyle \underset{j=1}{\overset{l}{}}}SU(N)_{i,j}\times {\displaystyle \underset{i=1}{\overset{k}{}}}{\displaystyle \underset{j=1}{\overset{l}{}}}SU(N)_{I,j}^{}`$ (3.5)
and with matter fields
| Field | Representation |
| --- | --- |
| $`(A_1)_{i+1,j+1;i,j}`$ | $`(\text{ }\text{ }\text{ }\text{ }\text{ }_{i+1,j+1},\overline{\text{ }\text{ }\text{ }\text{ }\text{ }}_{i,j}^{})`$ |
| $`(A_2)_{i,j;i,j}`$ | $`(\text{ }\text{ }\text{ }\text{ }\text{ }_{i,j},\overline{\text{ }\text{ }\text{ }\text{ }\text{ }}_{i,j}^{})`$ |
| $`(B_1)_{i,j;i,j+1}`$ | $`(\text{ }\text{ }\text{ }\text{ }\text{ }_{i,j}^{},\overline{\text{ }\text{ }\text{ }\text{ }\text{ }}_{i,j+1})`$ |
| $`(B_2)_{i,j;i+1,j}`$ | $`(\text{ }\text{ }\text{ }\text{ }\text{ }_{i,j}^{},\overline{\text{ }\text{ }\text{ }\text{ }\text{ }}_{i+1,j})`$ |
Moreover there is a quartic superpotential
$`W{\displaystyle }((A_1)_{i+1,j+1;i,j}(B_1)_{i,j;i,j+1}(A_2)_{i,j+1;i,j+1}(B_2)_{i,j+1;i+1,j+1}`$ (3.6)
$`(A_1)_{i+1,j+1;i,j}(B_1)_{i,j;i+1,j}(A_2)_{i+1,j;i+1,j}(B_2)_{i+1,j;i+1,j+1})`$
As explained in the section 2, by taking T-duality, we obtain the brane box configuration consisting of $`k`$ NS branes and $`l`$ NS’ branes whose intersections are smoothen out by diamonds. The singular point of $`𝒞_{kl}`$ splits into $`kl`$ ordinary conifold singularities $`x_{ij}`$ on $`\stackrel{~}{𝒞_{kl}}`$ under the resolution $`\stackrel{~}{\pi }:\stackrel{~}{𝒞_{kl}}𝒞_{kl}`$ as in equation (2.16). Hence if we put a large $`N`$ number of D3 branes at each singular point $`x_{ij}`$, the the near-horizon limit of the geometry of $`\stackrel{~}{𝒞_{kl}}`$ at $`x_{ij}`$ will be $`\mathrm{𝐀𝐝𝐒}_5\times 𝐓^{1,1}`$. If we wrap a D5 brane over the 2-cycle of $`\stackrel{~}{X_5}`$ corresponding to the 2-cycle of $`L_{ij}`$, we obtain a fractional D3 brane which is a domain wall in $`\mathrm{𝐀𝐝𝐒}_5`$ since it lies in the orthogonal direction to the D3 branes placed on the singular point $`x_{ij}`$. When we put one fractional brane together with $`N`$ regular D3 branes, we will change the $`(i,j)`$-th copy of the $`SU(N)`$ gauge group and the gauge group will change to $`SU(N+1)\times SU(N)^{kl1}\times SU(N)^{{}_{}{}^{}kl}`$ on the other side of the domain wall. The evidence for this claim is similar to the one of i.e. by studying the behavior of wrapped D3-branes on 3-cycles of $`L_{ij}`$ when they cross domain walls.
Before going further, we need to make a crucial observation concerning the constraint imposed by the consistency of the field theory on the worldvolume of the D3 branes. Since our field theory is chiral and can have anomalies, it is important to be careful with the way we introduce the fractional D3 branes. ¿From the geometrical discussion, it appears that there is no restriction on introducing fractional D3 branes i.e. on wrapping D5 branes on different 2-cycles of the horizon. In the brane box picture this would mean that there is no restriction on the number of D5 branes on different diamonds. If there is no integer D3 brane in the theory, we can introduce any number of fractional D3 branes which correspond to D5 branes in a specific diamond. In the presence of D3 branes orthogonal to the conifold (integer D3 branes), we cannot put fractional D3 branes in only one diamond. This is because if we put one D5 brane in the $`(i,j)`$-th diamond, the gauge groups in the $`(i,j)`$-th and $`(i+1,j+1)`$-th boxes have one supplementary anti-fundamental field and the gauge groups in the $`(i,j1)`$-th and $`(i1,j)`$-th boxes have one supplementary fundamental field, these four gauge theories becoming anomalous. See Figure 8 where we represent the $`ij`$ diamond and fields which are in the fundamental or anti-fundamental representations. This determines a specific way to introduce fractional branes. We need to have either $`k`$ fractional D3 branes corresponding to D5 branes on a row of diamonds or $`l`$ fractional D3 branes corresponding to D5 branes on a column of diamonds. Then all the gauge groups in different boxes have the same number of fundamental and anti-fundamental fields and are anomaly free. When we discussed about the twisted sector in section 2, we saw that they also correspond to D5 branes on rows or columns of boxes so the filling of boxes and diamonds is similar.
We can now proceed to obtain the main goal of this section i.e. to compare the $`\beta `$-function calculation in field theory living on the world-volume of the integer D3 branes with the solution of supergravity equations of motion in the presence of fractional D3 branes. To start with, we need to discuss more about the “dibaryon” operators and the domain walls in $`AdS_5`$. Because $`\stackrel{~}{X_5}`$ contains $`kl`$ copies of Einstein manifolds $`𝐓^{1,1}`$, we can wrap D3 branes over the 3-cycle of each $`𝐓^{1,1}`$ to obtain $`kl`$ types of “dibaryons”. Besides we have integer D3 branes and D5 branes wrapped on each of the $`kl`$ blow-up 2-cycles at each of the $`kl`$ singular points, the latter ones being the domain walls. For $`M_{i_0j_0}`$ D5 branes wrapped around the $`(i_0,j_0)`$-th 2-cycle, the gauge group changes from $`SU(N)^{kl}\times SU(N)^{{}_{}{}^{}kl}`$ to
$$SU(N+M_{i_0j_0})\underset{i=1}{\overset{k}{}}\underset{j=1}{\overset{l}{}}SU(N)_{ij}\times SU(N)^{{}_{}{}^{}kl}$$
(3.7)
where the pair $`(i,j)`$ does not take the value $`(i_0,j_0)`$. As discussed before, field theory results require D5 branes on either rows or columns of diamonds. and this means that we need to have $`M_{i_0j},j=1,\mathrm{},l`$ D5 branes wrapped around the $`i_0,j=1,\mathrm{},l`$ cycles, the gauge group changing to
$$\underset{j=1}{\overset{l}{}}SU(N+M_{i_0j})\underset{i=1,ii_0}{\overset{k}{}}\underset{j=1}{\overset{l}{}}SU(N)_{ij}\times SU(N)^{{}_{}{}^{}kl}$$
(3.8)
For $`M_{ij_0},i=1,\mathrm{},k`$ D5 branes wrapped around the $`i=1,\mathrm{},l;j_0`$ cycles the gauge group changes to
$$\underset{i=1}{\overset{k}{}}SU(N+M_{ij_0})\underset{i}{\overset{k}{}}\underset{j=1,jj_0}{\overset{l}{}}SU(N)_{ij}\times SU(N)^{{}_{}{}^{}kl}$$
(3.9)
We can now proceed to construct the Type IIB dual to the $`𝒩=1`$ supersymmetric field theory with the gauge group (3.8) or (3.9). We have $`kl`$ fluxes of RR 3-form through the 3-cycles $`C_{ij}`$ which are Hodge duals to the 2-cycles surrounding the singular points $`x_{ij}`$ for $`i=1,\mathrm{},k,j=1,\mathrm{},l`$:
$$_{C_{ij}^3}H^{RR}=M_{ij},i=1,\mathrm{},k,j=1,\mathrm{},l$$
(3.10)
Here we are identifying a 2-cycle of $`H_2(L_{ij})`$ with a 2-cycle of $`H_2(\stackrel{~}{X_5})`$. To obey the above observed rule, one needs to turn fluxes through all $`C_{i_0j}^3,j=1,\mathrm{},l`$ or all $`C_{ij_0}^3,i=1,\mathrm{},k`$ cycles with fluxes equal to $`M_{i_0j}`$ or $`M_{ij_0}`$ respectively, $`i_0,j_0`$ being some fixed indices.
We can now use the results of our previous section where we have completely identified the 2-cycles and the 3-cycles so the result is that the $`H^{RR}`$ which we need to turn on are:
$`H^{RR}{\displaystyle \underset{j}{}}M_{i_0j}e_{i_0j}^\psi (e_{i_0j}^{\theta _1}e_{i_0j}^{\varphi _1}e_{i_0j}^{\theta _2}e_{i_0j}^{\varphi _2}),\text{for fixed}i_0\text{and}j=1,\mathrm{},l`$ (3.11)
or
$`H^{RR}{\displaystyle \underset{i}{}}M_{ij_0}e_{ij_0}^\psi (e_{ij_0}^{\theta _1}e_{ij_0}^{\varphi _1}e_{ij_0}^{\theta _2}e_{ij_0}^{\varphi _2}),\text{for fixed}j_0\text{and}i=1,\mathrm{},k.`$ (3.12)
We now consider the Type IIB SUGRA equations of motion with the 2-form gauge potentials in the $`\mathrm{𝐀𝐝𝐒}_5\times X_5`$ background with constant $`\tau =C_0+ie^\varphi `$:
$`dG=iF_5G.`$ (3.13)
Here $`G`$ is the complex 3-form field strength,
$`G=H^{RR}+\tau H^{NSNS},`$ (3.14)
which satisfies the Bianchi identity $`dG=0`$.
If we choose $`C_0=0`$ and a constant dilaton, it follows from (3.11), (3.12) and (3.13) that
$`e^\varphi H^{NSNS}{\displaystyle \underset{j}{}}df_{i_0j}(r)(e_{i_0j}^{\theta _1}e_{i_0j}^{\varphi _1}e_{i_0j}^{\theta _2}e_{i_0j}^{\varphi _2}),\text{for fixed}i_0`$ (3.15)
or
$`e^\varphi H^{NSNS}{\displaystyle \underset{i}{}}df_{ij_0}(r)(e_{ij_0}^{\theta _1}e_{ij_0}^{\varphi _1}e_{ij_0}^{\theta _2}e_{ij_0}^{\varphi _2}),\text{for fixed}j_0`$ (3.16)
are two solutions for the NS 3-form corresponding to the specific choice for $`H^{RR}`$. Since $`F_5=\text{vol}(\mathrm{𝐀𝐝𝐒}_5)+\text{vol}(X_5)`$ and $`dH^{RR}=e^\varphi F_5H^{NSNS}`$ , we conclude
$`F_5H^{NSNS}=0`$ (3.17)
and the real part of (3.13) is satisfied for all $`f_{ij}`$. ¿From the imaginary part we have either
$`{\displaystyle \frac{1}{r^3}}{\displaystyle \frac{d}{dr}}\left(r^5{\displaystyle \frac{d}{dr}}f_{i_0j}(r)\right)M_{i_0j},\text{for fixed}i_0`$ (3.18)
or
$`{\displaystyle \frac{1}{r^3}}{\displaystyle \frac{d}{dr}}\left(r^5{\displaystyle \frac{d}{dr}}f_{ij_0}(r)\right)M_{ij_0},\text{for fixed}j_0`$ (3.19)
Thus we need to turn on an NS form as
$`B^{NSNS}e^\varphi {\displaystyle \underset{j}{}}M_{i_0j}\omega _{i_0j}\mathrm{log}r,\text{for fixed}i_0`$ (3.20)
or
$`B^{NSNS}e^\varphi {\displaystyle \underset{i}{}}M_{ij_0}\omega _{ij_0}\mathrm{log}r,\text{for fixed}j_0`$ (3.21)
The gauge couplings of the gauge theories are modified in the presence of the fluxes of the $`B^{NS}`$ through the various 2-cycles. The gauge coupling without B-flux is related to the string coupling constant as $`g^2=\frac{1}{2g_s}`$ where $`g_s`$ is the string coupling constant. If all the diamonds and boxes have the same area, then field theories corresponding to D5 branes on boxes and diamonds have the same coupling constant and this the meaning of $`g`$ in the previous formula. Since the B-fields (inverse of the gauge couplings) are areas on the torus, by changing the B-fluxes through the $`(i_0,j),j=1,\mathrm{},l`$ or $`(i,j_0),i=1,\mathrm{},k`$ cycles we change the areas of the diamonds. If we wrap D5 branes on the $`(i_0,j),j=1,\mathrm{},l`$ or $`(i,j_0),i=1,\mathrm{},k`$ cycles, the fluxes of B-field through the corresponding cycles modify and the gauge couplings change acoording to $`\frac{1}{g_s}_{C_{i_0j}^2}B^{NSNS},j=1,\mathrm{},l`$ or $`\frac{1}{g_s}_{C_{ij_0}^2}B^{NSNS},i=1,\mathrm{},k`$. The connection to the RG flow in field theory uses the relation:
$$\frac{1}{g_{i_0^2j}}\frac{1}{g^2}\frac{1}{g_s}(_{C_{i_0j}^2}B^{NSNS}1/2),\text{for fixed}i_0\text{and}j=1,\mathrm{},l$$
(3.22)
or
$$\frac{1}{g_{ij_0}^2}\frac{1}{g^2}\frac{1}{g_s}(_{C_{ij_0}^2}B^{NSNS}1/2),\text{for fixed}j_0\text{and}i=1,\mathrm{},k$$
(3.23)
The previous results for $`B^{NSNS}`$ allow us to rewrite (3.22) as
$$\frac{1}{g_{i_0j}^2}\frac{1}{g^2}M_{i_0j}\mathrm{log}r,\text{for fixed}i_0\text{and}j=1,\mathrm{},l$$
(3.24)
and (3.23) can be written as
$$\frac{1}{g_{ij_0}^2}\frac{1}{g^2}M_{ij_0}\mathrm{log}r,\text{for fixed}j_0\text{and}i=1,\mathrm{},k$$
(3.25)
Because the coordinate $`r`$ is seen as a field theory scale in the AdS/CFT conjecture, relations (3.24) and (3.25) give the supergravity dual of the scale dependence of the difference between the gauge couplings.
The above results are obtained by solving the supergravity equations of motion and we are now going to compare them with $`\beta `$-function calculations in $`𝒩=1`$ supersymmetric field theory which give:
$`{\displaystyle \frac{d}{d\mathrm{log}(\mathrm{\Lambda }/\mu )}}{\displaystyle \frac{1}{g_{i_0j}^2}}`$ $``$ $`3(N+M_{i_0j})2N(1\gamma _A\gamma _B)`$ (3.26)
$`{\displaystyle \frac{d}{d\mathrm{log}(\mathrm{\Lambda }/\mu )}}{\displaystyle \frac{1}{g^2}}`$ $``$ $`3N2(N+M_{i_0j})(1\gamma _A\gamma _B),\text{for fixed}i_0\text{and}j=1,\mathrm{},l`$
For each diamond $`(i_0j)`$ which belongs to the $`i_0`$-th row, the fields $`A`$ and $`B`$ which enter in the previous equations correspond to the bi-fundamental representations $`(A_1)`$ in $`(\text{ }\text{ }\text{ }\text{ }\text{ }_{i_0,j},\overline{\text{ }\text{ }\text{ }\text{ }\text{ }}_{i_01,j1}^{})`$ of $`SU(N+M_{i_0j})_{i_0,j}\times SU(N)_{i_01,j1}^{}`$, $`(A_2)`$ in $`(\text{ }\text{ }\text{ }\text{ }\text{ }_{i_0,j},\overline{\text{ }\text{ }\text{ }\text{ }\text{ }}_{i_0,j}^{})`$ of $`SU(N+M_{i_0j})_{i_0,j}\times SU(N)_{i_0,j}^{}`$, $`(B_1)`$ in $`(\text{ }\text{ }\text{ }\text{ }\text{ }_{i_01,j}^{},\overline{\text{ }\text{ }\text{ }\text{ }\text{ }}_{i_0,j})`$ of $`SU(N+M_{i_0j})_{i_0,j}\times SU(N)_{i_01,j}^{}`$ and $`(B_2)`$ in $`(\text{ }\text{ }\text{ }\text{ }\text{ }_{i_0,j1}^{},\overline{\text{ }\text{ }\text{ }\text{ }\text{ }}_{i_0,j})`$ of $`SU(N+M_{i_0j})_{i_0,j}\times SU(N)_{i_0,j1}^{}`$ where $`SU(N)_{i,j}^{}`$ represents the gauge group on the $`(i,j)`$ box and we use the fact that the $`(i_0,j),(i_0,j1),(i_01,j)`$ and $`(i_01,j1)`$ boxes are adjacent to the $`(i_0,j)`$ diamond. The same for the $`(i,j_0)`$ diamond, we obtain the formulas:
$`{\displaystyle \frac{d}{d\mathrm{log}(\mathrm{\Lambda }/\mu )}}{\displaystyle \frac{1}{g_{ij_0}^2}}3(N+M_{ij_0})2N(1\gamma _A\gamma _B)`$ (3.27)
$`{\displaystyle \frac{d}{d\mathrm{log}(\mathrm{\Lambda }/\mu )}}{\displaystyle \frac{1}{g^2}}3N2(N+M_{ij_0})(1\gamma _A\gamma _B)`$
where $`\gamma `$ are the anomalous dimensions of the fields $`A_1,A_2,B_1,B_2`$ and near the fixed point $`\gamma `$ close to $`1/4`$. By subtracting the second equation from the first in both (3.26) and (3.27), we obtain either
$$\frac{1}{g_{i_0j}^2}\frac{1}{g^2}M_{i_0j}[3+2(1\gamma _A\gamma _B)]\mathrm{log}(\mathrm{\Lambda }/\mu )$$
(3.28)
or
$$\frac{1}{g_{ij_0}^2}\frac{1}{g^2}M_{ij_0}[3+2(1\gamma _A\gamma _B)]\mathrm{log}(\mathrm{\Lambda }/\mu )$$
(3.29)
We use the identification of the spacetime radial coordinate $`r`$ with the field theory scale and we see that the Type IIB supergravity solution has reproduced the field theoretic beta function, this establishing the gravity dual of the logarithmic RG flow in the $`𝒩=1`$ supersymmetric $`_{i=1}^k_{j=1}^lSU(N+M_{ij})\times SU(N)^{{}_{}{}^{}kl}`$ gauge theory on $`N`$ regular and $`M_{ij},i=1\mathrm{},k,j=1,\mathrm{},l`$ fractional D3 branes. The agreement is between $`\frac{1}{g_{ij}^2N}\frac{1}{g^2N}`$ at order $`M/N`$ in the large $`N`$ limit.
In the authors have obtained an analytic form for the gravitational RG flow in the gauged 5-d Supergravity in the case of $`\mathrm{𝐀𝐝𝐒}_5\times 𝐓^{1,1}`$. Their study concerned the back-reaction of the metric and 5-form fields. Their results could be generalized to our case with the difference that the local geometry around each of the $`kl`$ singular points should be used instead of the global geometry.
In the case $`l=1`$, the orbifolded $`𝐙_k\times 𝐙_l`$ conifold becomes a generalized $`𝐙_k`$ conifold. The horizon of the generalized conifold is singular and we need to partially resolve in order to obtain a smooth Einstein manifold horizon with $`k`$ singular points. The procedure is just a particular case of our general recipe but the field theory on the worldvolume of D3 branes is not chiral in this case.
## 4 Acknowledgments
We would like to thank Keshav Dasgupta and Prabhakar Rao for stimulating discussions. The work of K. Oh is supported in part by NSF grant PHY-9970664.
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# Yaroslavl State UniversityPreprint YARU-HE-00/03hep-ph/0003216 Lepton pair production by high-energy neutrino in an external electromagnetic field
## Abstract
The process of the lepton pair production by a neutrino propagating in an external electromagnetic field is investigated in the framework of the Standard Model. Relatively simple exact expression for the probability as the single integral is obtained, which is suitable for a quantitative analysis.
Nowadays, it is well established that a medium makes an active influence on the quantum processes. It stimulates a constantly growing interest in the particle physics in medium, especially in view of possible astrophysical manifestations. It should be noted that a consideration of an intense electromagnetic field as the medium, along with a dense matter, is physically justified indeed. Really, the field strengths inside the astrophysical objects can reach the critical Schwinger value $`B_e=m_e^2/e4.4110^{13}`$ G <sup>1</sup><sup>1</sup>1We use natural units in which $`c=\mathrm{}=1.`$, and even could exceed it essentially. On the other hand, the situation is possible when the so-called field dynamical parameter $`\chi `$ of the relativistic particle <sup>2</sup><sup>2</sup>2Its definition is: $`\chi =e(p_\alpha F_{\alpha \beta }F_{\beta \sigma }p_\sigma )^{1/2}/m^3`$, where $`p_\alpha `$ is the particle four-momentum, $`F_{\alpha \beta }`$ is the electromagnetic field tensor. propagating in a relatively weak electromagnetic field, $`F<B_e`$ ($`F=`$ and/or $`B`$), could appear rather high. In this case the field in the particle rest frame can exceed essentially the critical value and is very close to the crossed field ($`\stackrel{}{}\stackrel{}{B}`$, $`=B`$). Thus, the calculation in a constant crossed field is the relativistic limit of the calculation in an arbitrary relatively weak smooth field. Consequently, the results obtained in a crossed field possesses a great extent of generality, and acquires interest by itself.
It is known that such intense electromagnetic fields allow the processes which are kinematically forbidden in vacuum, such as the creation of the lepton pair by a neutrino, $`\nu \nu \mathrm{}^{}\mathrm{}^+(\mathrm{}=e,\mu ,\tau )`$. It should be noted, that the $`\mathrm{}^{}\mathrm{}^+`$ pair can have a sufficiently large space-like total momentum in an electromagnetic field, due to the specific kinematics of a charged particle in this field. Therefore, the process with a relativistic neutrino becomes purely diagonal with respect to the neutrino flavor and is insensitive to its mass and to the mixing in the lepton sector.
The theoretical study of the process of the electron - positron pair production by a neutrino in the crossed field limit has a rather long history . A correct type of the dependence of the probability on the dynamical parameter $`\chi `$ in the leading log approximation, namely, $`\chi ^2\mathrm{ln}\chi `$, was found in the paper , where the numerical coefficient was wrong, however. In the succeeding papers attempts were made to adjust this coefficient and to find the next post-logarithm terms which could appear quite essential when $`\mathrm{ln}\chi `$ was not very large.
According to the definition of the problem in the crossed field approximation, one should speak about the ultrarelativistic neutrino only, which exists as the left-handed one due to the chiral type of its interaction in the frame of the Standard Model, even if the neutrino mass is non-zero. A lack of understanding of the fact that nonpolarized ultrarelativistic neutrinos did not exist in Nature, often led to the appearance of an erroneous extra factor 1/2 in expressions for the probabilities of the processes with initial neutrinos, due to an unphysical averaging over the neutrino polarizations, see e.g. the papers .
The results for the probability of the process $`\nu \nu e^{}e^+`$ in the crossed field, which were obtained in the listed papers, had essential distinctions. This probability for the case of a high-energy neutrino ($`\chi 1`$) can be presented in the following form
$$w(\nu \nu e^{}e^+)=Kw_0\chi ^2\left(\mathrm{ln}\chi \frac{1}{2}ln3\gamma _E+\mathrm{\Delta }\right),$$
(1)
where
$$w_0=\frac{G_F^2(g_V^2+g_A^2)m_e^6}{27\pi ^3E_\nu },$$
(2)
$`\gamma _E=0.577\mathrm{}`$ is the Euler constant, $`g_V,g_A`$ are the constants in the effective local Lagrangian of the $`\nu \nu ee`$ interaction, see Eq. (3) below. In the recent paper devoted to the study of the massive neutrino decay $`\nu _i\nu _je^{}e^+(m_i>m_j+2m_e)`$ in an external field, a comparance was also made of various results for the process probability. However, the statement which was made in about a mutual agreement of the results was incorrect, in our opinion. Really, the constants $`K`$ and $`\mathrm{\Delta }`$ introduced in Eq. (1), were obtained by the authors as follows, see Table 1.
Note that in the papers calculations were performed for the case of the electron – neutrino interaction via the $`W`$ \- boson only. To compare our Eq. (1) with the results of these papers, one should suppose in this formula $`g_V=g_A=1`$ , and $`g_V=g_A=|U_{ei}U_{e3}|`$ , correspondingly. Loss of the factor $`m_e/E_\nu `$ in the resulting expressions for the probability in Ref. was not the numerical mistake but rather the physical one, because it led to the loss of Lorentz invariance of the product of the probability and the neutrino energy.
As it was mentioned above, the formula (1) for the probability described a rather particular case of $`\mathrm{ln}\chi 1`$. On the other hand, the situation is realized under some physical conditions where the dynamical parameter takes the values which are moderately large, namely, $`\chi 1`$, but $`\mathrm{ln}\chi 1`$. The crossed-field approximation is valid in this case, but the condition $`\mathrm{ln}\chi 1`$ fails, and a consideration is necesary in Eq. (1) of the next terms of expansion over the inverse powers of the large dynamical parameter $`\chi `$. The expressions for the probability at an arbitrary value of the $`\chi `$ parameter, presented in some of the listed papers, have a cumbersome form of multiple integrals, which are not suitable for an analysis.
Hence the problem is urgent of obtaining the probability of the lepton pair ($`e^{}e^+`$ or $`\mu ^{}\mu ^+`$) production by a neutrino propagating in the crossed field for an arbitrary value of the $`\chi `$ parameter. In this paper we present our result for the probability of the process $`\nu \nu \mathrm{}^{}\mathrm{}^+`$ which is rather simple and suitable for a quantitative analysis.
We will consider the case of relatively low momentum transfers ($`|q^2|m_W^2`$). Under this condition, the weak interaction of neutrinos with charged leptons can be considered in the local limit by using the effective Lagrangian
$$=\frac{G_F}{\sqrt{2}}\left[\overline{\mathrm{}}\gamma _\alpha (g_Vg_A\gamma _5)\mathrm{}\right]\left[\overline{\nu }\gamma ^\alpha (1\gamma _5)\nu \right],$$
(3)
where $`g_V=\pm 1/2+2sin^2\theta _W,g_A=\pm 1/2`$. Here the upper signs correspond to the situation when the neutrino flavor coinsides with the lepton $`\mathrm{}`$ flavor ($`\nu =\nu _{\mathrm{}}`$), in this case both $`Z`$ and $`W`$ boson exchange takes part in a process. The lower signs correspond to the case $`\nu =\nu _{\mathrm{}^{}},\mathrm{}^{}\mathrm{}`$, when the $`Z`$ boson exchange is only presented in the Lagrangian (3).
We omit the details of calculations which can be found e.g. in our paper , and present here the result for the probability in a form of the single integral containing the Airy function:
$`w(\nu \nu \mathrm{}^{}\mathrm{}^+)={\displaystyle \frac{G_F^2(g_V^2+g_A^2)m_{\mathrm{}}^6\chi _{\mathrm{}}^2}{27\pi ^4E_\nu }}{\displaystyle \underset{0}{\overset{1}{}}}u^2duz\mathrm{\Phi }(z)\{{\displaystyle \frac{4}{1u^2}}(2L(u){\displaystyle \frac{29}{24}})`$
$`{\displaystyle \frac{15}{2}}L(u){\displaystyle \frac{47}{48}}+{\displaystyle \frac{1}{8}}\left(1+(1u^2)L(u)\right)\left(33{\displaystyle \frac{47}{4}}(1u^2)\right)`$
$`+{\displaystyle \frac{9}{16}}{\displaystyle \frac{g_A^2}{g_V^2+g_A^2}}[48L(u)+2(1+(1u^2)L(u))(283(1u^2))]\}.`$ (4)
Here $`\chi _{\mathrm{}}`$ is the dynamical parameter of the lepton with the mass $`m_{\mathrm{}}`$,
$`\chi _{\mathrm{}}=e(p_\alpha F_{\alpha \beta }F_{\beta \sigma }p_\sigma )^{1/2}/m_{\mathrm{}}^3`$, $`\mathrm{\Phi }(z)`$ is the Airy function
$$\mathrm{\Phi }(z)=\underset{0}{\overset{\mathrm{}}{}}𝑑t\mathrm{cos}\left(zt+\frac{t^3}{3}\right),z=\left(\frac{4}{\chi _{\mathrm{}}(1u^2)}\right)^{2/3},L(u)=\frac{1}{2u}\mathrm{ln}\frac{1+u}{1u}.$$
(5)
In the case when $`\chi _{\mathrm{}}1`$, one immediately obtains from Eq. (4) the formula for the probability containing the well-known exponential suppression:
$$w(\chi _{\mathrm{}}1)=\frac{3\sqrt{6}G_F^2m_{\mathrm{}}^6}{(16\pi )^3E_\nu }(3g_V^2+13g_A^2)\chi _{\mathrm{}}^4\mathrm{exp}\left(\frac{8}{3\chi _{\mathrm{}}}\right),$$
(6)
which agrees with corresponding formula of Ref. .
In the case when $`\chi _{\mathrm{}}1`$, it is easy to obtain from Eq. (4) the formula (1), where $`K=1,\mathrm{\Delta }=29/24`$, in agreement with the result of Ref. . It is not difficult also to find from our Eq. (4) the next term of expansion over the inverse powers of the dynamical parameter $`\chi _{\mathrm{}}`$. One obtains:
$`w(\chi _{\mathrm{}}1)`$ $`=`$ $`{\displaystyle \frac{G_F^2(g_V^2+g_A^2)m_{\mathrm{}}^6\chi _{\mathrm{}}^2}{27\pi ^3E_\nu }}\{\mathrm{ln}\chi _{\mathrm{}}{\displaystyle \frac{1}{2}}ln3\gamma _E{\displaystyle \frac{29}{24}}`$ (7)
$``$ $`{\displaystyle \frac{1}{\chi _{\mathrm{}}^{2/3}}}{\displaystyle \frac{9}{56}}{\displaystyle \frac{3^{1/3}\pi ^2}{\left[\mathrm{\Gamma }\left(\frac{2}{3}\right)\right]^4}}{\displaystyle \frac{19g_V^263g_A^2}{g_V^2+g_A^2}}\},`$
where $`\mathrm{\Gamma }(x)`$ is the gamma function.
It is seen from Eq. (7) that the correction term $`\chi _{\mathrm{}}^{2/3}`$ is not universal. It is relatively small and negative in the case when the neutrino flavor coinsides with the charged lepton flavor $`(\nu _e\nu _ee^{}e^+,\nu _\mu \nu _\mu \mu ^{}\mu ^+)`$. When the flavors of the neutrino and of the charged lepton are different $`(\nu _\mu \nu _\mu e^{}e^+,\nu _e\nu _e\mu ^{}\mu ^+)`$, the correction term is positive and relatively large.
The dependence of the probability of the process $`\nu \nu \mathrm{}^{}\mathrm{}^+`$ on the dynamical parameter $`\chi _{\mathrm{}}`$ in the region where its value is moderately large, is represented in Figs. 1 and 2. One can see that the correction term $`\chi _{\mathrm{}}^{2/3}`$ is more likely to worsen than to improve the presentation of the probability in this region. As the analysis shows, this term has a sense just of the correction for the values $`\chi _{\mathrm{}}10^5`$ only.
Therefore, our exact formula (4) should be used in a detailed analysis of the probability of the lepton pair production by a neutrino propagating in an external electromagnetic field, when the value of the dynamical parameter $`\chi _{\mathrm{}}`$ is moderately large.
Acknowledgements
We thank A.V. Borisov for useful discussion which stimulated this research.
The work is supported in part by the Russian Foundation for Basic Research under the Grant N 98-02-16694.
Figure captions
Fig. 1
The dependence of the probability of the process $`\nu \nu \mathrm{}^{}\mathrm{}^+`$ on the moderately large dynamical parameter $`\chi _{\mathrm{}}`$ in the case when the neutrino flavor coinsides with the charged lepton flavor $`(\nu _e\nu _ee^{}e^+,\mathrm{})`$: a) from the exact formula (4); b) from the approximate formula (1); c) from the formula (7) with the ’correction’ $`\chi _{\mathrm{}}^{2/3}`$.
Fig. 2
The same as in Fig. 1, in the case when the flavors of the neutrino and of the charged lepton are different $`(\nu _\mu \nu _\mu e^{}e^+,\mathrm{})`$.
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# Abstract
## Abstract
We observed comets 122P/1995 S1 (deVico) and C/1995 O1 (Hale-Bopp) with high spectral resolving power in order to determine the ratio of N$`{}_{}{}^{+}{}_{2}{}^{}`$/CO<sup>+</sup> in their comae. While we clearly detected the CO<sup>+</sup> in both of these comets, no N$`{}_{}{}^{+}{}_{2}{}^{}`$ was detected in either comet. From these spectra, we derive sensitive upper limits for N$`{}_{}{}^{+}{}_{2}{}^{}`$/CO<sup>+</sup>. These upper limits are substantially below other reported detections of N$`{}_{}{}^{+}{}_{2}{}^{}`$/CO<sup>+</sup> in other comets. We discuss the prior N$`{}_{}{}^{+}{}_{2}{}^{}`$ detections and compare them with our observations. The abundance of N<sub>2</sub> in comets is important to our understanding of the condensation of ices in the solar nebula. In addition, N<sub>2</sub> is a tracer of Ar so study of N<sub>2</sub> allows an understanding of the role of comets for delivering volatiles to the terrestrial planets. It appears that many, if not most, comets are depleted in N<sub>2</sub> and it will be necessary to search for a mechanism for depleting this molecule in order to be consistent with current models of the solar nebula.
### Introduction
Nitrogen is one of the more abundant elements in the universe and is therefore assumed to be an important constituent of the solar nebula and of the comets. Nitrogen probably exists in comets in the form of N<sub>2</sub> and other nitrogen-bearing molecules, including NH<sub>3</sub>. Indeed, the ratio N<sub>2</sub>/NH<sub>3</sub> is a sensitive indicator of conditions in the solar nebula. Lewis and Prinn (1980) point out that at high temperatures and low pressures the dominant equilibrium species of carbon, oxygen and nitrogen would be N<sub>2</sub>, CO, and H<sub>2</sub>O. Only as temperature and pressure regimes change will CH<sub>4</sub> and NH<sub>3</sub> be produced and N<sub>2</sub> and CO be depleted. They conclude that the conversion of N<sub>2</sub> to NH<sub>3</sub> and CO to CH<sub>4</sub> would be sufficiently slow relative to radial mixing in the primitive solar nebula so that only small amounts of NH<sub>3</sub> and CH<sub>4</sub> should be present. Conditions in the circumplanetary nebulae would be sufficiently different so that jovian planets might have increased NH<sub>3</sub> and CH<sub>4</sub> abundances (Prinn and Fegley 1981).
Comets delivered some of the volatiles that we see today in the atmospheres of the terrestrial planets, but it is not certain how important a source of volatiles the comets represent. Owen and Bar-Nun (1995a,b) have pointed out that N<sub>2</sub> is an important guide to the volatile abundances of comets because it is trapped and released by amorphous ice in a manner which is similar to argon (Bar-Nun et al. 1988). Using N<sub>2</sub> as a guide to the argon, one can determine the extent to which comets enriched the volatile and noble gas components of the terrestrial planets. Ices formed at low temperatures will trap gas from the surrounding nebula, fractionating the original mixture as a function of the local temperature. Thus, they suggest that comets which formed near Uranus and Neptune, at temperatures around 50 K, would be the source of noble gases for Earth and Mars, while the higher quantities of neon and argon in the atmosphere of Venus, compared with Earth, would require comets formed at colder temperatures, such as in the Kuiper belt, to be the deliverers of some of the volatiles.
We detect such species as NH, NH<sub>2</sub> and CN in every comet, so evidence of nitrogen carriers is easily available. Most of these species and their parents are chemically reactive in the comae of comets. Molecular nitrogen should be less reactive than species such as NH<sub>3</sub> or HCN. While spacecraft have flown past comet Halley with mass spectrometers onboard, measurement of N<sub>2</sub> is difficult with mass spectrometry since both N<sub>2</sub> and CO occupy the mass 28 bin of these instruments (cf. Everhardt et al. 1987 for a discussion of CO and N<sub>2</sub> from Giotto observations of Halley). Thus, disentangling the quantity of N<sub>2</sub> from the CO is very model dependent.
This leaves the field of ground-based spectroscopy for determining the quantity of molecular nitrogen. Ground-based studies of molecular nitrogen are very difficult, however, because of the N<sub>2</sub> abundance of the Earth’s atmosphere. To circumvent the difficulty in observing N<sub>2</sub>, past observations have concentrated on the N<sub>2</sub> ion, primarily through observations of the N$`{}_{}{}^{+}{}_{2}{}^{}`$ (0,0) band at 3914Å. This band is extremely weak and is expected to be seen only in the tails of comets. Care must be taken when observing this band since N$`{}_{}{}^{+}{}_{2}{}^{}`$ emission is also excited in the atmosphere of the Earth, especially near dusk and dawn, when comets are often observed. Auroral activity will also excite this band in the terrestrial atmosphere. Additionally, this weak feature can easily be confused with other, nearby, cometary emissions. Thus, accurate measurement of N$`{}_{}{}^{+}{}_{2}{}^{}`$ in cometary spectra requires both good spatial and spectral resolution to separate the features from that of the Earth and other cometary features.
### Observations
We observed comets deVico and Hale-Bopp with the 2DCoude spectrograph (Tull et al. 1995) on the 2.7-m Harlan J. Smith telescope of McDonald Observatory. The 2DCoude has two operating modes. The “lower” resolution mode has a resolving power, R=60,000. In this mode, spectral coverage is complete from around 3800-5800Å and coverage continues to 1 $`\mu `$m with increasing interorder gaps. Typically, 60–65 spectral orders are observed. Therefore, in the blue, many molecular bands can be observed simultaneously, regardless of the exact grating setting. In the “high” resolution mode, R=200,000, but the coverage is much less complete than the lower resolution mode. Typically, high resolution covers 10–15 orders of approximately 15Å each. Thus, care must be taken to center key features on a spectral order and many features remain unobserved.
For this project, we observed comet Hale-Bopp in the high resolution mode, carefully centering the portion of the order containing the N$`{}_{}{}^{+}{}_{2}{}^{}`$ band on the CCD, while also observing the CO<sup>+</sup> (2,0) and (3,0) bands on two other orders. Since the CH<sup>+</sup> (0,0) band occurs at a wavelength coincident with the CO<sup>+</sup> (2,0) band, this ion was also observed. Comet deVico was observed in the lower resolution mode and the same three ions were observed. Table I gives the circumstances of the observations. For all observations, the slit was 8.2 arcsec long. For the Hale-Bopp observations, the slit was 0.34 arcsec wide, while it was 1.2 arcsec wide for deVico. Each slit width projects to two pixels on the CCD for the resolving power of the observations. Different positions within the coma were observed by moving the telescope around the sky under accurate computer control.
The data were reduced using the echelle package of iraf. Incandescent lamp observations were used to determine the flat field; ThAr lamp observations were used for calculating the dispersion curve. The rms errors of our fits for the dispersion curve are 0.24 mÅ for the Hale-Bopp spectra and 2.5 mÅ for the lower resolution deVico spectra. The solar spectrum was observed with an identical instrumental setup to that used for the comets by imaging the Sun through a diffuser on the roof of the spectrograph slit room and projecting this image through the slit in the same manner as objects observed through the telescope. Thus, we used an observed solar spectrum in our reductions.
Care was taken to preserve the relative flux levels of the spectra. The spectral orders were extracted by first tracing the order along the chip and carefully setting the edges of the apertures. Since the continua of the cometary spectra were rarely of high enough signal/noise to define well the aperture boundaries, it was assumed that the flat lamp boundaries were appropriate for the comet and only the position on the chip of the center of the order was computed for the cometary spectra. Extraction was done using variance weighting. We used a 1D fit for stars and solar spectra, while a 2D fit was used for cometary spectra (because of the emission line nature of the cometary spectra). At the end of the routine reduction, we had files containing $`n`$ spectra for each initial spectral image, where $`n`$ is the number of extracted orders in the image (60 for deVico and 13 for Hale-Bopp).
The solar spectrum observations were used to remove the underlying continuum from the cometary spectra. Comet deVico has very little solar continuum, but the continuum of Hale-Bopp was quite strong. We corrected the comet and the solar spectrum for the geocentric and heliocentric Doppler shifts so that both were on a common rest frame. Then, the solar spectrum was carefully weighted to match the continuum level of the comet in regions away from cometary emissions. Some scattered light might still remain, but the amount is minimal and was removed when the line intensities were calculated.
Figure 1 shows the spectral order of the CO<sup>+</sup> (2,0) and the CH<sup>+</sup> (0,0) bands in the spectrum obtained 100 arcsec tailward of the optocenter of Hale-Bopp. For both these ions, the predicted molecular transitions in the spectrum are marked. Inspection of this figure shows that only very low $`J`$levels are observed for CH<sup>+</sup>, while slightly higher $`J`$levels are observed for CO<sup>+</sup>. However, even for CO<sup>+</sup>, $`J`$levels above 10 or 11 are not seen.
Figure 2 shows the same spectral region for observations 100 arcsec tailward of the optocenter of comet deVico. The spectral coverage of an order is longer at the lower resolving power of the deVico observations. Even with this larger coverage, only one of the CO<sup>+</sup> (2,0) ladders is seen in these observations. Both CH<sup>+</sup> and CO<sup>+</sup> are again present, though the ratio of CH<sup>+</sup>/CO<sup>+</sup> may be slightly different in these two comets.
Figure 3 shows the Hale-Bopp spectrum obtained 10 arcsec tailward of the optocenter in the spectral order which should contain the N$`{}_{}{}^{+}{}_{2}{}^{}`$ (0,0) band. The N$`{}_{}{}^{+}{}_{2}{}^{}`$ transition is a B $`{}_{}{}^{2}\mathrm{\Sigma }`$–X $`{}_{}{}^{2}\mathrm{\Sigma }`$ transition and therefore does not have a Q-branch. The P- and R-branch line positions marked are from Dick et al. (1978) and have an accuracy of 0.01 cm$`{}_{}{}^{1}=2`$ mÅ. Although the solar-subtracted spectrum is somewhat noisy, there are no believable features. There might be a spike at 3909Å and a broader spike at 3913.5Å. Neither of these is coincident with any of the N$`{}_{}{}^{+}{}_{2}{}^{}`$ line positions. The errors in the wavelengths of our spectra, coupled with the N$`{}_{}{}^{+}{}_{2}{}^{}`$ laboratory errors would lead us to expect coincidence to 2 mÅ. The most believable feature is the broad feature starting at 3914.5Å and degrading redward. The positions of the C<sub>3</sub> (0,2,0)-(0,0,0) band transitions are marked underneath the spectrum. The feature seems to match well the R-branch bandhead of this C<sub>3</sub> band. Thus, we conclude that if the feature is real, it is some residual C<sub>3</sub> emission. Since this spectrum was obtained only 10,000 km from the optocenter, this does not seem an unlikely species to observe. Inspection of the spectrum obtained 100 arcsec tailward of the optocenter shows even less evidence for features. We therefore conclude that we did not detect any N$`{}_{}{}^{+}{}_{2}{}^{}`$ in the spectrum of Hale-Bopp.
Figure 4 shows the comparable spectral region for comet deVico, 100 arcsec from the optocenter. There appears to an upward fluctuation between 3913.5 and 3915.1Å, with spikes at 3913.9 and 3914.9Å. However, at R=60,000, it is impossible to tell if this feature is a molecular band and, if so, which way it degrades, or to differentiate whether it is C<sub>3</sub> or N$`{}_{}{}^{+}{}_{2}{}^{}`$. The N$`{}_{}{}^{+}{}_{2}{}^{}`$ P-branch bandhead occurs at 3914.3Å and distinct lines of the P-branch should be visible. We do not detect lines at the expected wavelengths, within the 3 mÅ wavelength uncertainties. On the strength of the Hale-Bopp observation, this feature could be C<sub>3</sub>. With the larger spectral coverage of the R=60,000 orders, the blue end of this order contains the CN (0,0) bandhead (not shown in Fig. 4). Thus, we were able to use the high signal/noise CN emission lines to confirm that there were no errors in our wavelength solution or in our Doppler shift corrections. The centers of the CN lines fell at the correct wavelengths, verifying that any spikes in the N$`{}_{}{}^{+}{}_{2}{}^{}`$ region would also have to be at predicted wavelengths.
For the deVico observations, our only spectrum off the optocenter was obtained more than 70,000 km into the tail. Thus, it is reasonable to ask about the likelihood that we observed C<sub>3</sub> this far from the optocenter. Figure 5 shows the optocenter spectrum from the same night. We show more of the order to illustrate the abundance of molecular emissions observed. The CH B–X (0,0) band is clearly detected along with several C<sub>3</sub> bands. Indeed, inspection of this plot shows there are several broad, unidentified features whose structure seems similar to the identified C<sub>3</sub> bands. Thus, it is likely that this order is riddled with C<sub>3</sub>. However, comparison of the strength of the strongest C<sub>3</sub> band, the (0,0,0) – (0,0,0) band, in the optocenter and tail spectra makes it unlikely that we detected the much weaker (0,2,0) – (0,0,0) R-branch bandhead in the tail spectrum. Thus, it is most likely we detected only noise in the deVico spectrum.
In summary, we do not believe that any N$`{}_{}{}^{+}{}_{2}{}^{}`$ was detected in the spectra of comets Hale-Bopp or deVico. The CO<sup>+</sup> and CH<sup>+</sup> were clearly detected in both comet’s spectra. We are therefore able to place limits on the important ratio of N$`{}_{}{}^{+}{}_{2}{}^{}`$/CO<sup>+</sup> in these two comets.
### Limits on N$`{}_{}{}^{+}{}_{2}{}^{}`$/CO<sup>+</sup>
For both comets Hale-Bopp and deVico, the CO<sup>+</sup> (2,0) band was clearly detected. Thus, we can derive an abundance of CO<sup>+</sup> from these data for comparison with other comets. However, coudé data can not be easily calibrated into absolute fluxes, so we must work with band intensities in detector counts. In addition, typical lower resolving power observations observe all branches of the CO<sup>+</sup> band, while we only observe the $`{}_{}{}^{2}\mathrm{\Pi }_{1/2}^{}`$(F<sub>2</sub>) branches. We took the simple approach of “integrating” the band by fitting a continuum and summing the counts in the band above the continuum. We limited our bandpass to just the region of the detected lines. These values are given in Table II. While a larger bandpass would be more comparable to prior low-resolution observations of comets, it is inappropriate for high resolution observations since larger bandpasses would increase the noise with no increase in signal. Typical low-resolution observations do not return to continuum in between the lines of different bands.
We do not believe that we have detected the N$`{}_{}{}^{+}{}_{2}{}^{}`$ in either comet. We computed upper limits by computing how much of a band could be hidden within the noise. We did this by computing the rms in a bandpass. Then, the upper limit is just $`1/2\times \mathrm{rms}\times \mathrm{bandpass}`$, in appropriate units. These are $`2\sigma `$ upper limits. For the same rms, more signal can be hidden in a large bandpass than a small bandpass. Since we do not know exactly how many lines would be likely to be detected, we cannot easily define the bandpass. We assumed a bandpass which would include the complete P-branch of N$`{}_{}{}^{+}{}_{2}{}^{}`$. The derived counts, which are $`2\sigma `$ upper limits for what could be hidden in the noise, are listed in Table II, column 3.
With the use of a few assumptions and simplifications, we can use our values to derive N$`{}_{}{}^{+}{}_{2}{}^{}`$/CO<sup>+</sup> for these two comets. Examination of the solar spectra in these two regions, in comparison with the published atlas of Kurucz et al. (1984), shows that the sensitivity of the N$`{}_{}{}^{+}{}_{2}{}^{}`$ order is lower than that of the CO<sup>+</sup> order. To match their sensitivity, we would need to multiply the N$`{}_{}{}^{+}{}_{2}{}^{}`$ upper limits by a factor of 1.7. However, the calibration of the solar spectrum depends on details such as activity, so this factor is uncertain. We do have observations of $`\alpha `$ Lyr with the same instrumental setup as for deVico, but since the N$`{}_{}{}^{+}{}_{2}{}^{}`$ band occurs in the Balmer decrement, where the $`\alpha `$ Lyr flux changes rapidly with wavelength, the $`\alpha `$ Lyr flux is not calibrated in this region. Assuming a smooth decrease through the Balmer decrement, we confirm that a correction factor of 1.5–1.7 for the N$`{}_{}{}^{+}{}_{2}{}^{}`$ counts would be appropriate. We therefore adopt a factor of 1.6.
Once the band intensity is known, the column density can be computed using
$$N=L/g_{\nu ^{}\nu ^{\prime \prime }}$$
where N is the column density, L is the integrated band intensity and $`g_{\nu ^{}\nu ^{\prime \prime }}`$ is the excitation factor. We used excitation factors of $`7.0\times 10^2`$ photons sec<sup>-1</sup> mol<sup>-1</sup> for the N$`{}_{}{}^{+}{}_{2}{}^{}`$ (0,0) band (Lutz et al. 1993) and $`3.55\times 10^3`$ photons sec<sup>-1</sup> mol<sup>-1</sup> for the CO<sup>+</sup> (2,0) band (the average value from Figure 2 of Magnani and A’Hearn 1986). Then,
$$\frac{\mathrm{N}_2^+}{\mathrm{CO}^+}=\frac{g_{\mathrm{CO}^+}}{g_{\mathrm{N}_2^+}}\frac{\mathrm{L}_{\mathrm{N}_2^+}}{\mathrm{L}_{\mathrm{CO}^+}}$$
For CO<sup>+</sup>, we observed only one of the two ladders. If we assume the two ladders are equal strength, we should multiply our CO<sup>+</sup> intensity by two for the calculation. We likewise need to multiply the N$`{}_{}{}^{+}{}_{2}{}^{}`$ upper limits by a factor of two since we have only measured the P-branch and the R-branch should have a similar intensity. In Table II, column 4, we list our upper limits for N$`{}_{}{}^{+}{}_{2}{}^{}`$/CO<sup>+</sup>, including using a factor of 1.6 to correct for the sensitivity difference of the two orders.
It would be impossible to hide much N$`{}_{}{}^{+}{}_{2}{}^{}`$ in our spectra. Figure 6 (upper panel) shows one of the Hale-Bopp spectra, as observed, and, in the lower panel, the same spectrum with a feature added which has enough integrated counts to yield the Halley N$`{}_{}{}^{+}{}_{2}{}^{}`$/CO<sup>+</sup> ratio (Wyckoff and Theobald 1989 – discussed below). We do not claim that this synthetic band is the exact shape that would be present, nor are the “lines” at exactly the N$`{}_{}{}^{+}{}_{2}{}^{}`$ wavelengths, but it gives an idea of the ease with which we would detect such a feature. Clearly, no feature this distinctive could be missed in our observations.
### Previous Observations of N$`{}_{}{}^{+}{}_{2}{}^{}`$
Most comets are not bright enough to be observed with the high spectral resolving powers that we used for deVico and Hale-Bopp. This was especially true in the past, when detectors, such as photographic plates, had much lower quantum efficiency than our current CCD detectors. Therefore, prior observations which have detected N$`{}_{}{}^{+}{}_{2}{}^{}`$ in cometary spectra have been obtained with lower resolution, often on photographic plates.
The N$`{}_{}{}^{+}{}_{2}{}^{}`$ feature is generally weak and is overwhelmed by other molecular emissions near the optocenter. In addition, since N$`{}_{}{}^{+}{}_{2}{}^{}`$ is an ion, it is entrained in the solar wind magnetic field and rapidly accelerated into the tail. Thus, spectra of the tail region are necessary for its definitive detection, yet tail spectra are generally of lower signal/noise than near-optocenter spectra since the cometary brightness falls with increasing cometocentric distance. Despite these difficulties, observations of comets exist which show the detection of N$`{}_{}{}^{+}{}_{2}{}^{}`$ in the tails of comets.
Only two of the prior reported observations are digitally measured spectra; the rest are estimates from digital spectra or are photographic spectra. Wyckoff and Theobald (1989) report a detection of N$`{}_{}{}^{+}{}_{2}{}^{}`$ in the tail of comet Halley at a cometocentric distance of $`3\times 10^5`$ km tailward. These observations were at much lower resolution than our observations. They detected a weak emission in the region from 3885–3950Å which they concluded was composed of contributions from the CO<sup>+</sup> (5,1), CO$`{}_{}{}^{+}{}_{2}{}^{}`$ (unassigned), N$`{}_{}{}^{+}{}_{2}{}^{}`$ (0,0) bands and an unidentified band. By modeling the combined feature, they were able to estimate the contribution of N$`{}_{}{}^{+}{}_{2}{}^{}`$ to the mixture. Using this estimate and the average for the CO<sup>+</sup> (2,0), (3,0) and (4,0) column densities, they derived a value of N$`{}_{}{}^{+}{}_{2}{}^{}`$/CO$`{}_{}{}^{+}=0.004`$. However, the excitation factor which Wyckoff and Theobald used for N$`{}_{}{}^{+}{}_{2}{}^{}`$ was not accurate and Wyckoff et al. (1991b) revised the value of the column density of N$`{}_{}{}^{+}{}_{2}{}^{}`$ using the excitation factors of Lutz (1989–a personal communication). This excitation factor is the same as that given in Lutz et al. (1993). If we apply the value from Lutz et al., then N$`{}_{}{}^{+}{}_{2}{}^{}`$/CO$`{}_{}{}^{+}=0.002`$. If only the (2,0) band column density of CO<sup>+</sup> is used, then N$`{}_{}{}^{+}{}_{2}{}^{}`$/CO$`{}_{}{}^{+}=0.003`$.
Lutz et al. (1993) reported observations of the tails of two comets obtained at low resolution ($`10`$Å). For comet Halley, they obtained spectra at $`2\times 10^4`$ and $`2\times 10^5`$ km from the optocenter in the tailward direction. They claim to have detected no N$`{}_{}{}^{+}{}_{2}{}^{}`$ emissions in the Halley tail spectra. However, they also did not detect the CO<sup>+</sup> emissions in several Halley spectra. Their derived upper limits for N$`{}_{}{}^{+}{}_{2}{}^{}`$/CO<sup>+</sup> when CO<sup>+</sup> was detected were higher than the Wyckoff and Theobald detection.
In addition to observations of Halley, Lutz et al. also observed comet C/1987 P1 (Bradfield=1987 XXIX). For this comet, spectra were obtained at $`2\times 10^4`$ and $`6\times 10^4`$ km from the optocenter. CO<sup>+</sup> was detected in both spectra, but N$`{}_{}{}^{+}{}_{2}{}^{}`$ was only detected at the larger cometocentric distance. At their resolution, the N$`{}_{}{}^{+}{}_{2}{}^{}`$ feature is on the wing of the CN (0,0) band. No mention is made of the possible contamination of this feature by the CO<sup>+</sup> (5,1) band. Assuming all of their measured band was N$`{}_{}{}^{+}{}_{2}{}^{}`$, Lutz et al. derive a value of N$`{}_{}{}^{+}{}_{2}{}^{}`$/CO$`{}_{}{}^{+}=0.02`$.
The vast majority of observations of cometary tails were photographic. Not only were they at lower resolving powers than our observations of deVico and Hale-Bopp, but photographic plates are even more difficult to calibrate! Non-uniformity in response and vignetting of the spectrograph slit cause difficulty interpreting these spectra. Still, there are many fine examples of photographic spectra and these can be used to determine the N$`{}_{}{}^{+}{}_{2}{}^{}`$/CO<sup>+</sup> ratio for some comets. The largest published collection of photographic spectra is that of Swings and Haser (1956). Examination of the plates in this atlas shows some comets for which CO<sup>+</sup> and N$`{}_{}{}^{+}{}_{2}{}^{}`$ are both apparent, while other comets show evidence of tails (i.e. CO<sup>+</sup>) but no N$`{}_{}{}^{+}{}_{2}{}^{}`$. Arpigny examined these and other photographic and digital spectra at his disposal and estimated the intensity ratio for the N$`{}_{}{}^{+}{}_{2}{}^{}`$/CO<sup>+</sup>, where the CO<sup>+</sup> band used was the (4,0) band because of its proximity to the N$`{}_{}{}^{+}{}_{2}{}^{}`$ emission (1999, personal communication). Table III lists the 12 comets which he determined had both N$`{}_{}{}^{+}{}_{2}{}^{}`$ and CO<sup>+</sup> in these spectra, along with his estimate of the ratio of the intensity of the two bands (column 2). In order to compare his intensity ratios with other observer’s column density ratios, it is necessary to multiply by the ratio of the excitation factors, as before. Since the (2,0) CO<sup>+</sup> was used in our work and in other published ratios, we converted the intensity ratios in Table III by using the relationship $`I(4,0)=0.6\times I(2,0)`$, where $`I(4,0)`$ is the intensity of the (4,0) band, $`I(2,0)`$ is the intensity of the (2,0) band, and the factor is taken from Table 4 of Magnani and A’Hearn (1986). The resultant column density ratios are given in column 3 of Table III. Arpigny’s estimates of the intensity ratios are consistent with the published numbers for comet Bester (Swings and Page 1950) and comet Humason (Greenstein 1962). It should be noted that Warner and Harding (1963) also observed comet Humason at a comparable heliocentric distance (however they only discuss CO<sup>+</sup>, not N$`{}_{}{}^{+}{}_{2}{}^{}`$). However, it is clear from Arpigny’s compilation that N$`{}_{}{}^{+}{}_{2}{}^{}`$ has been observed in previous spectra of some comets.
In addition, Arpigny reported four comets which had good spectra but for which no, or only very faint, evidence of a plasma tail existed. These comets are C/1948 V1 (Eclipse), C/1963 A1 (Ikeya), C/1968 N1 (Honda), and C/1975 N1 (Kobayashi-Berger-Milon). Arpigny points out that N$`{}_{}{}^{+}{}_{2}{}^{}`$ emission is always very weak, so we should not expect to see it when the CO<sup>+</sup> is weak or non-existent.
Our own examination of the atlas of Swings and Haser (1956) found five comets for which there was evidence of a tail but no evidence for N$`{}_{}{}^{+}{}_{2}{}^{}`$. The Big Comet of 1910 (1910 I) showed only continuum in the tail, so this was presumably a dust tail. Comets Halley (1910 II), Brooks (1911 V), Gale (1912 II), and Jurlof-Achmarof-Hassel (1939 III) showed evidence of weak CO<sup>+</sup> emission but no N$`{}_{}{}^{+}{}_{2}{}^{}`$ emission (Arpigny notes that N$`{}_{}{}^{+}{}_{2}{}^{}`$ was observed by Bobrovnikoff in spectra of Halley obtained in 1910, but these spectra are not included in the Swings and Haser Atlas). These would be similar to Hale-Bopp and deVico in the absence of N$`{}_{}{}^{+}{}_{2}{}^{}`$ while other ions are present. However, with the weakness of the CO<sup>+</sup> emissions in these four photographically observed comets, the N$`{}_{}{}^{+}{}_{2}{}^{}`$ emission is most probably below the plate sensitivity.
### Implications
In this paper, we have presented high resolution observations of two comets with which we were able to study the relative abundances of N$`{}_{}{}^{+}{}_{2}{}^{}`$ and CO<sup>+</sup>. These two ions are proxies for understanding the quantity of N<sub>2</sub> and CO, two of the least chemically reactive cometary coma species. Conversion from the quantity of the ions to the quantity of the neutrals is dependent on an understanding of the photodestruction branching ratios which are not well understood (Wyckoff and Theobald argue you must multiply the ion ratio by a factor of 2, while Lutz et al. find no factor necessary), so we will continue to discuss these species in terms of their ions. For both Hale-Bopp and deVico, CO<sup>+</sup> was easily detected but N$`{}_{}{}^{+}{}_{2}{}^{}`$ appears to be missing from the spectra. Thus, we have put very low upper limits on the ratio of N$`{}_{}{}^{+}{}_{2}{}^{}`$/CO<sup>+</sup>. We note, however, that there are previous observations of comets which show CO<sup>+</sup> but not N$`{}_{}{}^{+}{}_{2}{}^{}`$ for which sensitive upper limits cannot be derived.
The quantity of N<sub>2</sub> and CO expected in a comet depends on several factors including the temperature at which the ice was deposited, when in the history of the formation of the solar system the gases were trapped in the ice and the orbital history of the comet itself. Current models of the solar nebula have comets which now reside in the Oort cloud forming in the Uranus-Neptune region (cf. Weissman 1991; Duncan, Quinn and Tremaine 1987). The temperature in this region was probably about $`50\pm 20`$ K (Boss et al. 1989). Thus, a first guess to the deposition temperature of cometary ices is 50K. This is consistent with laboratory experiments described by Owen and Bar-Nun (1995b).
The first direct measure of a deposition temperature for ice came with observations of deuterium in comet Hale-Bopp. Meier et al. (1998b) reported the detection of HDO in Hale-Bopp and determined a ratio of D/H=$`(3.3\pm 0.8)\times 10^4`$ in H<sub>2</sub>O. In addition, Meier et al. (1998a) detected DCN for the first time and derived a ratio of D/H=$`(2.3\pm 0.4)\times 10^3`$ in HCN. Note that the D/H ratio is different for these two species, with D/H measured from HCN 7 times higher than from H<sub>2</sub>O. Since the D/H enrichment for different molecules is a strong function of temperature, Meier et al. (1998a) were able to derive a temperature for the cloud fragment in which this comet formed of no colder than $`30\pm 10`$ K.
Bar-Nun et al. (1988) performed laboratory experiments on deposition of various gases along with H<sub>2</sub>O ice and showed that CO is trapped 20 times more efficiently than N<sub>2</sub> in amorphous ice which formed at 50 K, when these two gases are present in equal abundances with CH<sub>4</sub> and Ar. This ratio changes slightly when only CO and N<sub>2</sub> are present in the gas (Notesco and Bar-Nun 1996, Table I) but generally shows enrichment factors of 15–30. From these laboratory experiments, Owen and Bar-Nun (1995a) concluded that icy planetesimals formed in the solar nebula at around 50 K, the temperature at which the studied comets should have formed, would have N<sub>2</sub>/CO$`0.06`$ in the gases trapped in the ice if N<sub>2</sub>/CO$`1`$ in the nebula. The predicted cometary ratio of N<sub>2</sub>/CO is much higher than our upper limits for deVico and Hale-Bopp and is higher even than the detections of Wyckoff and Theobald (1989) and Lutz et al. (1993), though some of the estimates might show ratios this high.
Several factors might mitigate this discrepancy. Prialnik and Bar-Nun (1990) point out that the gas/water vapor ratio is not necessarily representative of the ratio of ices in the nucleus. However, the laboratory experiments of Bar-Nun et al. (1988) have demonstrated that CO and N<sub>2</sub> should be released simultaneously in the same proportion as they exist in the ices. There is evidence for a source of CO at around 10,000 km from the nucleus (Eberhardt et al. 1987) which may be attributable to grains. In addition, Krankowsky (1991) points out that H<sub>2</sub>CO is probably an additional parent for CO. Thus, there may be additional mechanisms for the production of CO which do not exist for N<sub>2</sub>. While these factors might change the predicted ratio for N<sub>2</sub>/CO, Owen and Bar-Nun (1995a) go on to make the specific prediction that “future observations of dynamically new comets will show values of N<sub>2</sub>/CO systematically higher than those in well-established short-period comets”. They point out that as comets are continuously exposed to solar radiation in the inner solar system, they would be expected to lose any N<sub>2</sub> in their outer layers in a brief period of time.
Figure 7 shows all of the values and limits discussed in this paper, plotted as a function of 1/a<sub>o</sub>, the original semimajor axis \[four of the values for the estimates are the oscullating 1/a, as noted in Table III; for C/1987 P1 (Bradfield) 1/a<sub>o</sub>= 0.006380 (Marsden and Williams 1995); for Hale-Bopp 1/a<sub>o</sub>=0.00535 and for deVico 1/a<sub>osc</sub>=0.057 (Marsden personal communications 1999)\]. For the Wyckoff and Theobald Halley observation, we include the range of derived values.
If one ignores the Hale-Bopp upper limit, there would seem to be an increase of N$`{}_{}{}^{+}{}_{2}{}^{}`$/CO<sup>+</sup> with decreasing 1/a<sub>o</sub>, as was predicted by Owen and Bar-Nun. However, the trend is based mostly on Arpigny’s estimates, which are approximate. In addition, the Hale-Bopp upper limit can not be discarded since it is clear from inspection of the spectrum in Figure 6 that even as much N$`{}_{}{}^{+}{}_{2}{}^{}`$ as would be needed to equal the Halley N$`{}_{}{}^{+}{}_{2}{}^{}`$/CO<sup>+</sup> cannot be hidden in this spectrum. Conversely, while the values for the ratio can be questioned for the estimates, inspection of the atlas of Swings and Haser (1956) and spectra such as the Humason spectrum of Greenstein (1962) clearly show an emission which is coincident with the location of the N$`{}_{}{}^{+}{}_{2}{}^{}`$ feature. Thus, at least some comets have some N$`{}_{}{}^{+}{}_{2}{}^{}`$ in their spectra. However, at lower spectral resolution, blending of features may lead to spurious detections of N$`{}_{}{}^{+}{}_{2}{}^{}`$ and wrong estimates of the strength of this band. The potential for blending, coupled with the weakness of the N$`{}_{}{}^{+}{}_{2}{}^{}`$ feature even when it exists, point to a need for caution in interpretation.
Thus, at this point, we have contradictory evidence for the ratio of N<sub>2</sub>/CO. At least in the active outer regions of Hale-Bopp and deVico, these comets appear to be very depleted in N<sub>2</sub> relative to CO. The observations of Halley by Wyckoff and Theobald (1989) also pointed to a depletion of N<sub>2</sub> for Halley. Indeed, Wyckoff et al. (1991a) derived a nitrogen depletion for comet Halley of a factor of $`6`$ relative to the Sun. Owen and Bar-Nun (1995) have pointed out the strong temperature dependence of trapping of N<sub>2</sub> and CO. Comets for which the deposition temperature was greater than 50K could trap progressively less CO and N<sub>2</sub>. However, H<sub>2</sub>CO will continue to be trapped in comparable quantities to ices deposited at 50K. Thus, there would continue to be a source of CO, but the N<sub>2</sub> will be depleted relative to the CO.
Perhaps the solution to the quandary of depleted nitrogen is that our assumption that the solar nebula preferentially condenses nitrogen into N<sub>2</sub> instead of NH<sub>3</sub> is incorrect. However, observations of dense molecular clouds (Womack et al. 1992) have shown that N$`{}_{2}{}^{}>>`$ NH<sub>3</sub> for these potential star-forming sites. Another possibility is that the comets formed with much more molecular nitrogen but that it was depleted post-formation. It is certainly true that comets will deplete volatile gases in their outer layers as they pass close to the Sun, but this cannot be used as an explanation when comparing comets with similar orbital histories, such as Halley and deVico or Hale-Bopp and Bennett and C/1987P1 (Bradfield), which show discrepant ratios of N$`{}_{}{}^{+}{}_{2}{}^{}`$/CO<sup>+</sup>. Indeed, Engel et al. (1990) conclude that there might be some post-formational processing of Halley, but not to any large extent because of the low internal temperatures which are derived from the spin temperature of H<sub>2</sub>O (Mumma et al. 1993). They point out, however, that gas can be trapped in water ice efficiently, but only if the ice is amorphous, such as it is in the various laboratory experiments. They conclude that codeposition into amorphous ice is unlikely to be a favorable mechanism in forming cometary ices since the water ice condenses, for most solar nebula models, at 140–160 K, at temperatures where the ice is likely to be crystalline and to not adsorp volatiles readily. However, the D/H ratios of H<sub>2</sub>O and HCN suggest that cometary deposition temperatures were not as warm as Engel et al. posit.
In summary, in this paper we presented evidence of two comets for which no N$`{}_{}{}^{+}{}_{2}{}^{}`$ was detected, along with stringent upper limits, which would indicate that these comets are depleted in N<sub>2</sub> relative to CO. These observations are at odds with our understanding of the formation processes of ices in the solar nebula. Either a mechanism must be found to deplete the N<sub>2</sub> ice once formed or we must understand how a gaseous cloud with N$`{}_{2}{}^{}>>`$ NH<sub>3</sub> formed ices which do not contain much molecular nitrogen. Since we believe that N<sub>2</sub> will be deposited into H<sub>2</sub>O ice in a manner which is similar to Ar, understanding this process has important implications for understanding the role of comets for delivery of noble gases to the terrestrial planets. It is therefore important that more unambiguous observations of the ion tails of comets be obtained, when possible, to determine the intrinsic values of N<sub>2</sub>/CO in comets.
Acknowledgements
This work was funded by NASA Grant NAG5 4208. We thank Walter Huebner for encouraging the Hale-Bopp observations and Toby Owen for stimulating our examination of the data. We especially thank Claude Arpigny for graciously allowing us to use his estimates of earlier observations and for many helpful discussions.
### Figure Captions
Figure 1: The spectral region of the CO<sup>+</sup> (2,0) and CH<sup>+</sup> (0,0) bands for Hale-Bopp. The positions of lines within this bandpass are marked, though not all marked lines are present. The pattern of detected vs. non-detected lines can be explained by the excitation levels of these molecules.
Figure 2: The spectral region of the CO<sup>+</sup> (2,0) and CH<sup>+</sup> (0,0) bands for deVico. The spectral orders are longer for the R=60,000 mode so more CH<sup>+</sup> lines were detected. Note that CH<sup>+</sup> is stronger relative to CO<sup>+</sup> in deVico than in Hale-Bopp.
Figure 3: The spectral region of the N$`{}_{}{}^{+}{}_{2}{}^{}`$ (0,0) band for Hale-Bopp. In addition, the C<sub>3</sub> (0,2,0)–(0,0,0) band is in this spectral region. The expected positions of lines for both of these bands are marked. There appears to be a feature at 3914.5Å, which we tentatively attribute to C<sub>3</sub>, not N$`{}_{}{}^{+}{}_{2}{}^{}`$.
Figure 4: The spectral region of the N$`{}_{}{}^{+}{}_{2}{}^{}`$ (0,0) band for deVico. As with Hale-Bopp, there appears to be a feature at 3914.5Å. However, we believe this feature is just noise.
Figure 5: An optocenter spectrum of deVico, including the region of the (0,0) band of N$`{}_{}{}^{+}{}_{2}{}^{}`$. Several C<sub>3</sub> bands and the CH B–X (0,0) band are identified. Additional, unidentified, C<sub>3</sub> bands are probably present. We show this spectrum to show the abundance of C<sub>3</sub> in this comet.
Figure 6: Simulated N$`{}_{}{}^{+}{}_{2}{}^{}`$ data. The upper panel shows an actual spectrum of Hale-Bopp. The positions of the potential N$`{}_{}{}^{+}{}_{2}{}^{}`$ lines are marked beneath it. The lower panel shows this same spectrum but with a “fake” band replacing the data between 3913 and 3914Å. The fake band would have enough counts so that the ratio of N$`{}_{}{}^{+}{}_{2}{}^{}`$/CO<sup>+</sup> would equal the value detected for Halley. In the real data, there are some upward excursions which do not correspond to any N$`{}_{}{}^{+}{}_{2}{}^{}`$ lines. Even integrating just these upward excursions yields only 1/4 the necessary counts.
Figure 7: Values for N$`{}_{}{}^{+}{}_{2}{}^{}`$/CO<sup>+</sup> as a function of 1/a<sub>o</sub>. The various ratios and limits are plotted in order to examine whether a trend exists with dynamical age of the comets. The prediction for this ratio of Owen and Bar-Nun (1985a) is shown as a dotted line. The open symbols are values derived from estimates of features in photographic and digital spectra, while the closed symbols represent measured digital spectra. See text for a discussion.
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# The European Large Area ISO Survey I: Goals, Definition and Observations
## 1 Introduction
The Infrared Space Observatory was the natural successor to the Infrared Astronomical Satellite (IRAS), and has primarily been used to undertake detailed studies of individual objects and regions. However, ISO also provided an opportunity to perform survey work at sensitivities beyond the reach of IRAS. The IRAS survey was of profound significance for cosmology, extragalactic astrophysics and for the study of stars, star-forming regions and the interstellar medium in the Galaxy. The mapping of large-scale structure in the galaxy distribution, the discovery of ultra-luminous infrared galaxies (see the review by Sanders & Mirabel ) and of hyper-luminous infrared galaxies like IRAS F10214+4724 \[1991a\], and the detection of proto-planetary discs around fairly evolved stars, were all unexpected discoveries of the IRAS survey. The $`z=2.3`$ galaxy F10214+4724, was at the limit of detectability by IRAS ($`S_{60}\genfrac{}{}{0pt}{}{_>}{^{}}0.2`$Jy). Several other $`z>1`$ galaxies and quasars have now been found from follow-up of faint IRAS samples. Recent sub-mm surveys, in particular with SCUBA on the JCMT, (e.g. Smail et al. Hughes et al. Barger et al. Eales et al. Blain et al. ) are detecting sources which are probably very high redshift counterparts to these IRAS sources. Pointed observations of high redshift quasars and radio galaxies produce detections at sub-millimetre wavelengths in continuum and line emission, but mostly lie below the limit of the IRAS survey at far infrared wavelengths.
While designed as an observatory instrument, the huge improvement in sensitivity provided by ISO offered the opportunity to probe the galaxy population to higher redshift than IRAS and to make progress in understanding the obscured star formation history of the Universe. A significant fraction of the mission time was thus spent on field surveys. In this paper we describe the “European Large Area ISO Survey” (ELAIS) which represents the largest non-serendipitous survey conducted with ISO. This survey provides a link between the IRAS survey, the deeper ISO surveys and the sub-mm surveys.
ELAIS is a collaboration involving 25 European institutes, led from Imperial College. This project surveyed around 12 square degrees of the sky at 15$`\mu `$m and 90$`\mu `$m nearly 6 square degrees at 6.7$`\mu `$m together with a further one square degree at 175$`\mu `$m. The survey used the ISO Camera at the two shorter wavelengths and the ISO Photometer at the longer wavelengths. ELAIS was the largest open time project undertaken by ISO: a total of 375 hours of scientifically validated data have been produced. We have detected over 1000 extra-galactic objects and a similar number of Galactic sources. Around 200 of these objects have been re-observed with ISO to provide detailed mid/far infrared photometry.
This paper outlines the broad scientific objectives of this project and describes the selection of the observing modes and survey fields. It also details the execution of the ISO observations and briefly outlines the data reduction and data products. Finally we show how this survey complements other ISO surveys and summarise the extensive multi-wavelength programmes taking place in the ELAIS fields.
## 2 Key Scientific Goals
### 2.1 The Star Formation History of the Universe
The main extra-galactic population detected by IRAS was galaxies with high rates of star formation. These objects are now known to evolve with a strength comparable to Active Galactic Nuclei (AGN) (e.g. Oliver et al. ). The distance to which these objects were visible to IRAS was, however, insufficient to determine the nature of their evolution. The sensitivity of ISO allows us to detect these objects at much higher redshifts and thus obtain greater understanding of the cosmological history of star formation. The infrared luminosity provides a better estimate of the total star formation rate than optical and UV estimators (e.g. Madau et al.) as these monitor star formation only from regions with low obscuration and require large corrections for extinction . Another important star formation indicator for galaxies is the radio luminosity (e.g. Condon ). For galaxies obeying the well known far infrared radio correlation , the depth of the survey described here is well matched to that of sub-mJy radio surveys (e.g. Condon & Mitchell , Windhorst , Windhorst et al. , Hopkins et al. , Gruppioni, Mignoli & Zamorani ). Comparison of the global star formation rate determined in the infrared with other determinations from the optical and UV luminosity densities, $`H\alpha `$ luminosity density, radio luminosity density, etc. will give a direct estimate of the importance of dust obscuration, vitally important for models of cosmic evolution, as well as providing us with a reliable estimate for the total star formation rate. The ELAIS follow-up surveys (see section 7) will allow us to go a stage further and apply a number of these complementary star formation tracers to the same volume and in many cases on the same objects, thereby addressing the impact of dust extinction independently of any peculiarities to any particular survey volume.
Figures 1-3, show the predicted redshift distribution of star-forming galaxies in the ELAIS survey selected at 15$`\mu `$m, 90$`\mu `$m and 175$`\mu `$m. The predictions come from three different evolutionary models; the first model is that of Pearson & Rowan-Robinson , the second and third are models ‘A’ and ‘E’ from Guiderdoni et al. . All three models are extrapolations from IRAS data. The total number of objects of various different types predicted by two of these models and a third from Franceschini et al. , are also tabulated in Table 1. While the source counts from ELAIS alone may not be able to distinguish between such models, spectroscopic identifications, source classifications and the redshift distributions will.
### 2.2 Ultra-luminous Infrared Galaxies at High $`z`$
IRAS uncovered a population with enormous far infrared luminosities, $`L_{\mathrm{FIR}}`$ $`>10^{12}L_{}`$ (see the review by Sanders & Mirabel 1996). While somewhere between 20 and 50 per cent of these objects appear to have an AGN (Veilleux et al. , Sanders et al. , Veilleux et al. , Lawrence et al. ) it is still a source of controversy as to whether the illumination of the dust arises principally from an AGN or a star-burst. ISO spectra of samples of ultra-luminous infrared galaxies (Genzel et al. , Lutz et al. , Lutz, Veilleux & Genzel, Rigopoulou et al. ) appear to demonstrate that while some do require the photoionization energies typical for AGN to explain the obscured lines, most are consistent with star-burst models. Interestingly, most of these objects appear to be in interacting systems, suggesting a mechanism that could trigger either an AGN, a star-burst, or indeed both (e.g., Sanders et al. , Lawrence et al. , Leech et al. , Clements et al. ).
The area of this survey is small compared to that of IRAS so we would not expect to detect large numbers of these objects. The Pearson & Rowan-Robinson (1996) model would predict that we would detect between 40 and 80 of these objects, though models such as that of Guiderdoni et al. (1998), which takes into account the increase in temperature of the dust with increasing luminosity would predict more. Nevertheless such objects will be visible at greater distances than they were in IRAS and even a few examples at higher redshift would be interesting. Assuming a star-burst SED an object of $`L_{60}=10^{12}L_{}`$ ($`H_0=50,q_0=1/2`$) would be visible ($`S_{15}>3\mathrm{m}\mathrm{J}\mathrm{y}`$) in the ELAIS survey to $`z=0.5`$ where it is only visible to $`z=0.26`$ in the IRAS Faint Source Catalog ($`S_{60}>0.2\mathrm{Jy}`$) and to $`z=0.15`$ in the IRAS Point Source Catalog ($`S_{60}>0.6\mathrm{Jy}`$). ELAIS thus allows us to study samples of these controversial objects at higher redshift where both AGN and star formation are known to be enhanced. Figure 4 shows the minimum 60$`\mu `$m lluminosity of a source which could be detected in both the ELAIS survey and the IRAS survey as a function of redshift.
### 2.3 Emission from Dusty Tori around AGN
The orientation-based unified models of AGN involve a central engine surrounded by an optically and geometrically thick torus (Antonucci & Miller , Scheuer , Barthel , Antonucci ). In this model the optical properties of the central regions are dependent on the inclination angle of the torus, with type 2 objects defined as those with the central nucleus obscured by the torus, and type 1 objects (such as quasars) as those with an unobscured view of the nucleus. Objects with radio jets have the jets aligned approximately with the torus symmetry axis. The scheme is very attractive in providing a single conceptual framework for what would otherwise appear to be extremely diverse populations, and the models have survived many observational tests and predictions. It is now widely accepted that the unified models are broadly correct at least to “first order” (e.g. Antonucci 1993) and that many if not most type 2 AGN contain obscured type 1 nuclei.
An important corollary of the unified models is the expectation that populations of obscured (i.e. type 2) AGN will be present all redshifts. These predicted populations are in general extremely difficult to identify observationally (e.g. Halpern & Moran ) except locally in low-luminosity AGN, and at high redshift ($`0<z\stackrel{<}{_{}}5`$) in the radio-loud AGN minority. Nevertheless, the strength and shape of the X-ray background has been taken as evidence of the existence of a large population of obscured quasars, outnumbering normal quasars by a factor of several (e.g. Comastri et al. ). Such a large population of obscured quasars may also explain the unexpectedly large population of local remnant black holes (Fabian & Iwasawa , Lawrence ). Even hard X-ray samples may miss the very heavily obscured objects, so an infrared-selected sample is the only reliable way to obtain a complete census of AGN. For example, it will be possible with ELAIS to make quantitative constraints on the dust distribution and torus column densities, as well as on the evolution of obscured quasar activity.
### 2.4 Dust in Normal Galaxies to Cosmological Distances
At the longer ISO wavelengths (90 and 175 $`\mu `$m) emission from the cool interstellar ‘cirrus’ dust in normal galaxies will be detectable in our survey in fainter and cooler objects than were accessible to IRAS. This will allow us to examine the temperature distribution functions and in particular look for unusually cool galaxies. Quantifying the distributions of such cool sources will be important for deep sub-mm surveys as there is considerable degeneracy between cool, low redshift and warm high redshift objects in this wavelength regime.
### 2.5 Circumstellar Dust Emission from Galactic Halo Stars
We expect to detect hundreds of stars at 6.7 and 15 microns and it will be of interest to check whether any show evidence of an infrared excess due to the presence of a circumstellar dust shell. Such shells are expected from late type stars due to mass-loss while on the red giant branch, from cometary clouds or from proto-planetary discs. At the high galactic latitudes of our survey, late type stars with circumstellar dust shells should be rare (e.g. Rowan-Robinson & Harris, ), so any detections of such shells could be especially interesting.
### 2.6 New classes of Galactic and Extra-Galactic Objects
F10214+4724 \[1991a\] was at the limit of IRAS sensitivity and new classes of objects may well be discovered at the limit of the ELAIS sensitivity. The lensing phenomenon which made F10214+4724 detectable by IRAS may become more prevalent at fainter fluxes, increasing the proportion of interesting objects.
### 2.7 The Extra-Galactic Background
The discovery of the $`140850`$$`\mu `$m far infrared background (Puget et al. , Fixsen et al. , Hauser et al. , Lagache et al ) from COBE data has shown that most of the light produced by extra-galactic objects has been reprocessed by dust and re-emitted in the far infrared and sub-mm. This discovery provides further strong motivation for studying the dust emission from objects at all redshifts and all far infrared wavelengths. It is possible to explain this far infrared background radiation with a number of evolution models that are consistent with the IRAS data. The constraints provided by ISO surveys such as ELAIS are expected to be able to rule out some of these a priori models. The motivation behind our 175$`\mu `$m survey was specifically to start to resolve this far infrared background into its constituent galaxies.
## 3 Survey Definition
### 3.1 Selection of survey wavelengths and area
In order to detect as many sources as efficiently as possible we restricted ourselves to two primary ISO broad band filters and aimed to cover as large an area as possible. We selected filters with central wavelengths at: 15$`\mu `$m which is particularly sensitive to AGN emission and 90$`\mu `$m which is sensitive to emission from star formation regions. At 90$`\mu `$m we aimed to reach the confusion limit and pre-flight sensitivity estimates led us to conclude that this could be achieved with an on-sky integration time of 20s. We decided to map the same area of sky at 15$`\mu `$m using a similar total observation time and this required on-sky integration times of 40s. In both cases these integration times were close to the minimum practical. A survey area of order 10 square degrees was chosen to produce a statistically meaningful sample of galaxies. This area and depth was ideal to complement the deep ISO-CAM surveys (Cesarsky et al. (1996), Elbaz et al. , Taniguchi et al. \[1997a\]) as discussed in Section 6.
A further justification for a large area survey is that many of the sources will be at relatively low redshift (e.g. an ultra-luminous star-burst would be detectable at $`z=0.5`$ as discussed in Section 2.2). Thus, unless our survey is of a sufficient area, the volume will be such that cosmic variance can be a significant problem, i.e. large-scale clustering means that the mean density within a survey volume may not be representative of the universal mean. To estimate this effect we use the galaxy power spectrum as compiled by Peacock & Dodds . From this we can estimate the variance in a survey of any given volume (we assume a cubical geometry, which means we will underestimate the variance). Figure 5 illustrates the area required to study populations out to a given redshift allowing for different amounts of cosmic variance. From this we can see that a survey of around 10 square degrees is required to measure the mean density of populations visible to $`z=0.5`$ with negligible errors ($`<10`$ per cent) due to large-scale structure. A survey with the area of ELAIS can also measure the mean density of populations $`z=0.25`$ with 20 per cent accuracy. Populations below $`z=0.15`$ would only have mean densities known to around 50 per cent. Figure 6 shows what fractional errors we would expect in mean quantities derived from ELAIS for populations that are visible to different depths.
During the mission we introduced two additional filters. The first of these was designed to provide constraints on the infrared spectral energy distribution of ELAIS sources from fields (around six square degrees) that would not have been observed in time for pointed ISO follow-up. For this aspect of the survey we selected the 6.7$`\mu `$m filter, which was the most sensitive for sources detected at 15$`\mu `$m. Naturally as well as providing improved spectral coverage of other ELAIS sources this also produced an independent source list which was sensitive to emission from normal galaxies. The second filter, centred at 175$`\mu `$m, was introduced specifically to explore the populations making up the far infrared background as discussed in Section 2.7.
A more detailed description of the survey parameters is given Section 3.4.
### 3.2 Time Awarded
Over the course of the ISO mission the ELAIS programme was awarded a total of 377 hours. This allocation was used not only to perform the basic blank field survey observations, discussed in Section 3.1, but also a number of other related programmes.
Principal among these was an ISO photometry programme to investigate around 200 sources that had been detected by ELAIS in the early parts of the mission. These observations were designed to provide constraints on the spectral energy distributions of the ELAIS sources but would also provide a serendipitous, though biased, survey in their own right. In addition we were awarded time to observe a number of sub-fields repeatedly to help quantify our reliability and completeness. We also performed eight ISO-PHOT calibration measurements on three known stars and three ELAIS sources, independently of the instrument team.
The amount of time actually spent and AOTs used on both the survey proper and the photometry programmes are summarised in Table 2.
### 3.3 Field Selection
The allocated observing time was sufficient to observed around 12 square degrees. The choice of where to distribute the ELAIS rasters on the sky was governed by a number of factors. Firstly, we decided not to group these all in a single contiguous region of the sky; this further reduces the impact of cosmic variance on the survey (see Section 3.1). Distributing the survey areas across the sky also has advantages for scheduling follow-up work. Cirrus confusion is a particular problem, so we selected regions with low IRAS 100$`\mu `$m intensities ($`I_{100}<1.5`$MJy/sr), using the maps of Rowan-Robinson et al. \[1991b\]. In recognition of the large amount of time required we decided to minimise scheduling conflict with other ISO observations by further restricting ourselves to regions of high visibility ($`>25`$ per cent) over the mission lifetime, while to reduce the impact of the Zodiacal background we only selected regions with high Ecliptic latitudes ($`|\beta |>40^{}`$). Finally, it was essential to avoid saturation of the ISO-CAM detectors, so we had to avoid any bright IRAS 12$`\mu `$m sources ($`S_{12}>0.6`$Jy). These requirements led us to selecting the four main fields detailed in the upper portion of Table 3. The location of all ELAIS fields are indicated in Figure 7 showing the Galactic Cirrus distribution, while in Figures 8-11 we show the nominal boundaries of each of the main survey fields overlayed on a Cirrus map.
Towards the end of the mission an additional field (S2) was selected with similar criteria, this field was multiply observed to provide reliability and completeness estimates.
A further 6 fields were selected as being of particular interest to warrant a single small ($`24{}_{}{}^{}\times 24^{}`$) raster. These were chosen either because of existing survey data or because the field contained a high redshift object and were thus more likely to contain high redshift ISO sources.
1. Phoenix: This field was the target of a deep radio survey and has been extensively followed up from the ground with imaging and spectroscopy.
2. Lockman 3: This was one of the deep ROSAT survey fields .
3. Sculptor: This field has been the subject of an extensive ground based optical survey programme (e.g. Galaz & De Lapparent ).
4. VLA 8: This field is centred on a $`z=2.394`$ radio galaxy and has been the target a deep Hubble Space Telescope observations .
5. TX0211-122: This object is at $`z=2.34`$ and was discovered in the Texas radio survey .
6. TX1436+157: This object is at $`z=2.538`$ and was discovered in the Texas radio survey . The ISO field centre is offset from the radio object as the B1950 equinox coordinates were entered rather than the J2000 coordinates.
These 6 regions are also described in the lower portion of Table 3.
### 3.4 Observation Parameters
Table 4 summarises the instrument parameters specified in the majority of our survey AOTs. Most are self explanatory.
For ISO-CAM the $`Gain`$ was set to 2 which was the standard used for most ISO-CAM observations. $`TINT`$ was the integration time per readout and $`NEXP`$ was the number of readouts per pointings (i.e. the total integration time per pointing is $`NEXP\times TINT`$). $`NSTAB`$ was the additional number of readouts for the first pointing of a raster added to allow the detector to stabilise.
With ISO-PHOT, $`TINT`$ was the total integration time per pointing.
The parameters related to the raster geometry ($`PFOV`$, $`NPIX`$, $`M,N`$, $`dM,dN`$) have the same meaning for each instrument. $`PFOV`$ is the nominal pixel field of view on the sky. $`NPIX`$ is the number of pixels along each axis of the detector array. $`M,N`$ are the number of steps in a raster while $`dM,dN`$ are the step sizes.
The ISO-CAM rasters were designed such that each sky position was observed twice in successive pointings to improve reliability. To reduce overheads we selected a very large raster size, $`40{}_{}{}^{}\times 40^{}`$. With the exception of small rasters and one test raster, the ISO-CAM parameters remained unchanged throughout the survey.
Since the ISO-PHOT internal calibration measurements were only performed at the beginning and end of a raster we chose these to be half the size of the ISO-CAM rasters ($`20{}_{}{}^{}\times 40^{}`$). We originally used ISO-PHOT with a non-overlapping raster pattern and switched to an overlapping mode during the mission, with most N1 and S1 observations performed in the non-overlapping mode.
The observation parameters for all survey observations are tabulated in Appendix A.
## 4 ISO Observations
The ISO observations for the ELAIS programme were executed from 12th March 1996 (revolution 116), 37 days after the beginning of routine operations (4th February 1996, revolution 79) until 17:44 on 8th April 1998 (revolution 875), 10 hours 44 minutes after the first signs of boil off had been detected and 5 hours 23 minutes before the last observations were performed.
In general the execution of the planned observations was very successful. Only three observations were reported as “failed”. Three observations were flagged as “aborted”, all three of these had been concatenated to “failed” observations but appear to have been successfully executed despite this.
The only significant problem in the execution of the survey observations occurred in N3. It transpired that there was a paucity of guide stars in this region and the mission planning team were unable to schedule many of the observations near the original dates requested. To accommodate this problem the sizes of the rasters were reduced and restrictions on the possible observation dates relaxed. However, in the last available observing window for N3 other ISO mission priorities, together with remaining guide star acquisition problems, interfered with the scheduling. The net result is that the coverage of the N3 region is patchy.
It may be that this guide star problems noticed in N3 may be related to an apparent offset of around 6” between the reference frame of the DSS and e.g. the APM catalogue in this field. The APM catalogue agrees very well with the Guide Star Catalogue v1.2 (`http://www-gsss.stsci.edu/gsc/gsc12/gsc12_form.html`).
### 4.1 Main Survey Observations
Table 5 indicates the area that has been surveyed at least once in any band in all of our fields. For the four large fields the separation of the raster pointings (40) is used to compute the area, i.e. 0.44 square degrees per raster. For the small fields which are not mosaiced the actual size of the raster is used.
The coverage, in terms of integration time per sky pixel, of the four main survey fields (N1-3 and S1) in each of the bands are shown in Figures 12, 13,14,15.
### 4.2 Duplicate Survey Observations
A number of sub-fields have been repeated on one or more occasions. This repetition will considerably aid in assessing the reliability and completeness of the survey. In addition this data will provide deeper survey regions which are good targets for more focussed follow-up campaigns and other exploitation.
Table 6 lists all the fields that have repeated observations together with the level of redundancy.
### 4.3 Photometry Programme
In addition to the survey observations, we also undertook a photometry programme to observe objects detected early on in the survey programme at other ISO wavelengths.
These objects were selected from the S1 and N1 survey regions which had been observed at an early stage in the campaign. 180 objects which had been detected at 15 $`\mu `$m were selected to be observed with ISO-CAM at 4.5, 6.7, 9, and 11 $`\mu `$m using the filters LW-1,LW-2,LW-4 and LW-7. 80 Objects were selected to be observed by ISO-PHOT at 60 and 175 $`\mu `$m, using the C60, and C160 filters.
The ISO-CAM observations were performed in concatenated chains of 10 pointings. At each pointing a $`2\times 1`$ raster was performed to ensure accurate photometry and reliable detections. The chains were arranged such that each of the 10 sources was located in a different position on the array (separated by around 18<sup>′′</sup>), this was to allow accurate sky flat-fielding over the course of the concatenated chain. The 120 ISO-CAM pointings in S1, and the 80 pointings in N1 were ordered to minimize the total path length, ensuring that sequential observations were as close to each other as possible, both spatially and temporally, improving the flat-fielding.
The ISO-PHOT observations were performed in chains of 15 pointings. On average the 15 pointings contained 5 source positions and 10 background positions. Like the ISO-CAM photometry observations, the ISO-PHOT source positions were ordered to minimize the total path length. The background pointings were chosen to be spaced along this path at reasonably regular intervals, while ensuring that there was at least one background position between every source position.
Other parameters from the AOTs for the photometry programme are summarised in Table 7.
## 5 Data Processing and Products
In order to provide targets early on in the campaign to allow follow-up programmes, both from the ground and with ISO, it was decided to perform an initial “Preliminary Analysis”. This was started long before the end of the mission, while the understanding of the behaviour of the instruments was naturally less than it is currently and will be superseded with a “Final Analysis” incorporating the best available knowledge post mission.
The Preliminary Analysis was conducted with the intention of producing reliable source lists at $`6.7,15`$ and $`90\mu `$m. The processing of the ISO-CAM survey observations is described in detail by Serjeant et al. and the reduction of the ISO-PHOT 90$`\mu `$m survey data will be discussed by Efstathiou et al. (1999, in preparation).
The Final Analysis is currently being undertaken. This is expected to produce better calibrated and fainter source lists than the Preliminary Analysis. The Final Analysis will also produce maps which can be used to determine fluxes or upper-limits for known sources. This analysis will not, however, be completed until early 2000.
The ELAIS products will comprise source catalogues at all wavelengths, 4.5, 6.7, 9, 11, 15, 60, 90, 175, together with maps from all the survey observations. Highly reliable sub-sets of the “Preliminary Analysis” catalogues were released to the community, via our WWW site (`http://athena.ph.ic.ac.uk/`), concurrent with the expiration of the propriety period on 10th August 1999.
### 5.1 Data Quality
The quality of the 15 $`\mu `$m ISO-CAM data is moderately uniform. Some rasters are more affected by cosmic rays than others but the total amount of data seriously affected by cosmic rays is small. The noise levels are within a factor of a few of those expected; a typical noise level is 0.2 ADU/sec/pixel per pointing.
The ISO-PHOT data is seriously affected by cosmic rays and detector drifts. We have used the fluctuations in the time sequence of each pixel as an estimate of the average noise level. The fluctuations per pointing were typically 3 per cent of the background level, though 3 of the 9 pixels were noisier with fluctuations typically 4 per cent of the background. A few observations showed higher noise due to increased cosmic ray hits. Our original AOTs employed an integration time of 20s. We subsequently decreased this to 12s to allow for an overlapping raster giving a factor of two redundancy with similar observation time. Importantly there does not appear to be a significant difference in noise levels per pointing despite the factor of two reduction in integration time, indicating that non-white noise in the pixel histories is dominant. The redundancy introduced by this new strategy could improve the signal to noise ratio for sources by as much as $`\sqrt{2}`$.
### 5.2 Preliminary Data Analysis
The processing for the ISO-PHOT and ISO-CAM data proceeds in a similar fashion. All data reduction used a combination of standard routines from the PHOT Interactive Analysis <sup>1</sup><sup>1</sup>1PIA is a joint development by the ESA Astrophysics Division and the ISO-PHOT consortium) software and the CAM Interactive Analysis together with purpose-built IDL routines. The frequency of glitches and other transient phenomena led to non-Gaussian and non-white-noise behaviour.
A number of data reduction techniques were tested at ICSTM, CEA/SACLAY, IAS and MPIA. Parallel pipeline processes for reducing the ISO-PHOT data were run at both ICSTM and MPIA. Data reduction techniques suitable for ISO-CAM data with multiple redundancy, such as the observations of the Hubble Deep Field , e.g. the Pattern REcognition Technique for ISO-CAM data were unsuccessful in processing this data. The most reliable approach for source extraction was found to be looking for source profiles in the time histories of individual pixels rather than by constructing sky maps.
For both instruments the data stream from each detector pixel was treated as an independent scan of the sky. These data streams were filtered to remove glitches and transients and averaged to produce a single measurement at each pointing position. Significant outliers remaining in the data streams were flagged as potential sources. For the ISO-CAM observations the redundancy of the pointings was used to provide confirmation of candidate sources. The data stream surrounding all remaining candidates was then examined independently by at least two observers to remove spurious detections. Sources that were acceptable to two or more observers were classified as good ($`REL=2`$) and those acceptable to only one observer were classified as marginal $`(REL=3)`$.
The fraction of spurious detections was high due to the non-Gaussian nature of the noise and relatively low thresholds applied. More than 13 thousand ISO-PHOT source candidates were examined as were just over 15 thousand ISO-CAM 15 $`\mu `$m candidates. At 6.7 $`\mu `$m the rejected fraction was lower and the candidate list was only 3 thousand. The final numbers of objects in the Preliminary Catalogue Version 1.3 are tabulated in Table 8
The “eye-balling” technique while laborious ensured that the resulting catalogues are highly reliable, as discussed in greater detail in Serjeant et al. and Efstathiou et al. (1999., in preparation). The sub-sets of the Preliminary Catalogues that were released to the community were those ISO-PHOT sources that had been confirmed by four observers, and those ISO-CAM sources that had been confirmed by two observers with fluxes above 4mJy, these sub-sets are exceptionally reliable.
A “Final Analysis” process has been developed which uses the transient correction techniques of Lari (1999, in preparation). These techniques have been shown to be excellent for reducing ISO-CAM data. While this is almost certainly the best procedure for reducing the ELAIS data, it is labour intensive and time consuming and we do not expect the “Final Analysis” to be finished until early 2000, hence the release of our “Preliminary” products.
### 5.3 Source Calibration
For the ISO-CAM observations we have of order 10 stars per raster and these provide a very good calibration. A preliminary analysis of the star fluxes (Crockett et al. 1999 in preparation, see also Serjeant et al. 1999) suggests that our raw instrumental units (ADU/g/s) need to be multiplied by a factor of 1.75 to give fluxes in mJy. This implies a $`50`$ per cent completeness limit of approximately $`3`$ mJy at $`15\mu `$m. The flux calibration is still uncertain at $`6.7\mu `$m, due largely to the uncertain aperture corrections to the under-sampled observations and the single-pixel detection algorithm, though PSF models and pre-flight sensitivity estimates suggest a $`50`$ per cent completeness level at less than $`1`$mJy. (See Serjeant et al. for more details.)
For the 90 $`\mu `$m survey the calibration proceeded as follows. The expected background was estimated using COBE and IRAS data and Zodiacal light models. These predictions were compared to the measurements of the background calibrated using the internal calibration device (FCS) allowing the predicted backgrounds to be corrected from an extended source to a point source calibration. These predictions were then used to scale the measured fluctuations above the background. Single pixel detections (“Point sources”) were then calibrated using the expected fraction of flux falling on a single pixel for a source placed arbitrarily with respect to the pixel centre. The fluxes of “extended sources” were calculated in a more complicated fashion and have great associated uncertainties. The fluxes were found to be in good agreement with model stellar fluxes in our own dedicated calibration measurements and with the fluxes of IRAS sources in the fields. This suggests a 5$`\sigma `$ noise level of 100mJy. This ISO-PHOT calibration, completeness and reliability estimates is discussed in detail by Efstathiou et al. (1999, in preparation) and Surace et al. (1999, in preparation).
## 6 Comparison with Other ISO Surveys
ISO carried out a variety of complementary surveys exploring the available parameter space of depth and area. Table 9 summaries the main extra-galactic blank-field surveys. With the exception of the two main serendipity surveys ELAIS covers the largest area and has produced the largest number of ISO sources. Figure 16 illustrates how deeper smaller area surveys are complemented by shallower wider area surveys.
## 7 Follow-up
An extensive follow-up programme is being undertaken, including observations in many bands from X-ray to radio. This programme will provide essential information for identifying the types of objects detected in the infrared, their luminosities, energy budgets and other detailed properties. As well as studying the properties of the objects detected by ISO a number of the follow-up surveys will provide independent source lists which will be extremely valuable in their own right, e.g. to investigate the differences between infrared and non infrared emitting objects.
### 7.1 Surveys
A number of follow-up programmes are in fact independent surveys at other wavelengths, carried out within the ELAIS survey area. These include:
1. Optical: $`R`$-band CCD surveys are essential to provide optical identifications for spectroscopic and related follow-up with improved astrometry, photometric accuracy and to fainter levels than those provided by the Second Sky Survey. Our principal southern field (S1) has been completely covered with the ESO/Danish 1.5m telescope to a depth of $`R23.5`$ (La Franca et al. 1999, in preparation), while all our northern fields N1-3 have been completed to a similar depth using the INT Wide Field Camera (Verma et al. 1999, in preparation). Other optical bands allow object classification and other more detailed investigations. Four square degrees within our northern fields have been observed to a depth of $`U22`$ (Verma et al. 1999, in preparation). In June 1999, we observed the central 1.2 square degrees of S1 in $`U`$ and 3 square degrees in $`I`$ using the ESO Wide Field Imager (Héraudeau et al. 1999, in preparation) with these we expect to reach $`U23,I23`$. McMahon et al. (`http://www.ast.cam.ac.uk/~rgm/int_sur/`) have covered around 9 square degrees of N1 to $`u23.3,g24.2,r23.5,i22.7,z21.1`$ and 2 square degrees in N2 to similar depths in $`g,r,i,z`$ as part of the ING wide field survey. The $`U`$-band surveys will be especially interesting as they will allow us to compute the $`U`$-band luminosity density (and hence star formation rate) in the same volume as we calculate the infrared luminosity density, providing a direct comparison between obscured and unobscured star formation estimators.
2. Near Infrared: A substantial area has been surveyed in the near infrared. In the $`H`$-band around 0.85 square degrees in N1 and N2 was surveyed using CIRSI on the INT (Gonzalez-Solares et al. 1999, in preparation). Approximately 0.5 square degree has been surveyed in both N1 and N2 in $`K^{}`$ using Omega Prime on the Calar Alto 3.5m (Rigopoulou et al. 1999, in preparation). The smaller, multiply repeated southern field S2 has been covered in $`K`$ with SOFI on the NTT (Héraudeau et al., 1999, in preparation).
3. Radio: 21cm radio data at sub-mJy level will allow identification of some of the most interesting objects which are expected to be very faint in the optical but would have detectable radio fluxes if they obey the usual radio to far infrared correlation. These surveys will also allow an independent estimate of the star formation rate within the same volume. The southern field S1 is completely covered to a depth of 0.3 mJy , the 6 square degrees in the northern fields has been covered to a depth of 0.2 mJy . A deeper survey in the south has been conducted on the smaller, multiply repeated field S2 (Gruppioni et al. 1999, in preparation).
4. X-ray: Almaini et al. have been awarded 150ks to do two deep Chandra pointings one in N1 and one in N2. La Franca et al. have also been awarded 200 ks on BeppoSAX to make 5 pointings covering around 2 square degrees in S1.
5. Sub-mm: The UK SCUBA Survey consortium (Rowan-Robinson et al., independent of ELAIS) are performing part of their shallow (8mJy) 850$`\mu `$msurvey in N1 and N2 and are aiming to cover 200 square arc minutes in each.
Additional multi-wavelength surveys of these fields are expected in the near future.
### 7.2 Photometry & Spectroscopy
We intend to obtain spectroscopic identifications for all (or the vast majority) of optical candidates for all ELAIS sources. This involves a two-pronged attack using multi-object spectroscopy for the brightest objects and single object spectroscopy using 4m class telescopes on the fainter objects. This will be principally to obtain the redshifts and thus luminosity but also for classification and to assess star formation rates. Some preliminary multi-fibre spectroscopy has been carried out with FLAIR on the UK Schmidt Telescope. This been supplemented by single object spectroscopy from the ESO/Danish 1.5m telescope to provide spectroscopy on a complete sample of 90$`\mu `$m selected sources (Linden-Vørnle et al. 1999, in preparation). A further 100 sources have been identified spectroscopically in a largely weathered-out run on the 2dF in September 1998 (Gruppioni et al. 1999, in preparation) and an additional night on the 2dF in August 1999 was also seriously hampered by weather (Oliver et al. 1999, in preparation). 40 spectra for fainter sources have already been taken with EMMI on the NTT and EFOSC2 on the ESO 3.6m Telescope (La Franca et al. 1999, in preparation).
Until now the ELAIS northern fields have been only moderately surveyed spectroscopically. We have used the Calar Alto 2.2m telescope and the Calar Alto Faint Object Spectrograph (CAFOS) for the sources brighter than 17 (Gonzalez-Solares et al. 1999, in preparation) and the Calar Alto 3.5m telescope and the Multi Object Spectrograph (MOSCA) for the fainter sources (Surace et al., 1999, in preparation). 29 ELAIS objects have been observed during the period May-July 1998 (of which 14 had $`m>19`$). These observations have been completed with 29 field galaxies chosen in the same region for comparison purpose. From these northern samples most sources show strong star-burst signatures up to $`z=0.5`$ though two AGNs and one $`z=1.2`$ QSO have been detected, these samples will be discussed in a forthcoming paper.
A number of programmes have been instigated to obtain more specific photometric and spectroscopic data of the infrared selected sources over a wider wavelength range. Some examples are detailed below:
1. We have observed (Héraudeau, Kotilainen, Surace et al., in preparation) about 150 sources in pointing observations in the S1 field using IRAC2 on the ESO/MPG 2.2m telescope October 1997, June 1998 and SOFI on the NTT October 1998
2. Near-infrared H+K band spectroscopy of a small subset of sources with SOFI on the NTT (Alexander et al., in preparation).
## 8 Conclusions
In this paper we have described the motivation behind ELAIS, the largest non-serendipitous survey performed by ISO. Our primary goals in conducting the survey were to determine the relative importance and recent evolution of the dust–obscured mode of star formation in galaxies, and to constrain AGN unification models, and we detailed above how these influenced our selection of survey fields and observational parameters. The fields that have been covered by ISO are also being extensively mapped from radio to X–ray wavelengths as part of a concerted ground–based follow–up programme, whose multi–wavelength coverage will make the ELAIS regions fertile ground for undertaking future astrophysical investigations extending well beyond our initial survey aims.
Subsequent papers in this series will discuss in detail the scientific results from the ELAIS “Preliminary Analysis” and “Final Analysis”. The first of these papers will include: discussions of the extra-galactic counts from the “Preliminary Analysis” at 7 and 15 $`\mu `$m(Serjeant et al. 1999), and at 90$`\mu `$m (Efstathiou et al., 1999, in preparation); discussion of the stellar calibration and counts (Crockett et al., in preparation); and a discussion of sources detected in the multiply–repeated areas (Oliver et al., in preparation). Preliminary ELAIS data products were released through our WWW page (`http://athena.ph.ic.ac.uk/`), which also contains further details on the programme and the follow-up campaign.
## Acknowledgments
This paper is based on observations with ISO, an ESA project, with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) and with participation of ISAS and NASA.
The ISO-CAM data presented in this paper was analysed using “CIA”, a joint development by the ESA Astrophysics Division and the ISO-CAM Consortium. The ISO-CAM Consortium is led by the ISO-CAM PI, C. Cesarsky, Direction des Sciences de la Matiere, C.E.A., France.
PIA is a joint development by the ESA Astrophysics Division and the ISOPHOT Consortium. The ISO-PHOT Consortium is led by the Max-Planck-Institut fuër Astronomie (MPIA), Heidelberg, Germany. Contributing ISO-PHOT Consortium institutes to the PIA development are: DIAS (Dublin Institute for advanced studies, Ireland) MPIK (Max-Planck-Institut fuër Kernphysik, Heidelberg, Germany), RAL (Rutherford Appleton Laboratory, Chilton, UK), AIP (Astronomisches Institut Potsdam, Germany), and MPIA.
This work in part was supported by PPARC (grant number GR/K98728) and by the EC TMR Network programme (FMRX-CT96-0068).
We would like to thank all the ISO staff at Vilspa both on the science team and on the Instrument Development Teams for their eternal patience in dealing with the wide variety of problems that a large programme like this presented.
## Appendix A Log of the ISO Observations
In Table A1 we present a list of all the observations performed by ISO as part of the ELAIS raster observations. We do not include the observations performed as part of the ISO photometric follow-up of ELAIS sources which are available on our WWW pages `http://athena.ph.ic.ac.uk/`. Table 14 details those observations which have had some instrument or telemetry problems as flagged at Vilspa or for which we have noted peculiarities.
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# Definitive Computation of Bernstein-Sato Polynomials
## 1. Introduction
Throughout this paper $`k`$ is a field of characteristic $`0`$, $`R_n(k)=k[x_1,\mathrm{},x_n]`$ is the ring of polynomials in $`n`$ variables and $`A_n(k)=kx_1,\mathrm{},x_n,_1,\mathrm{},_n`$ is the corresponding Weyl algebra, i.e. an associative $`k`$-algebra generated by $`x`$’s and $``$’s with the relations $`_ix_i=x_i_i+1`$ for all $`i`$.
For every polynomial $`f`$$`R_n(k)`$ there are $`b(s)k[s]`$ and $`Q(x,,s)A_n(k)[s]`$ such that
(1.1)
$$b(s)f^s=Q(x,,s)f^{s+1}.$$
The polynomials $`b(s)`$ for which equation (1.1) exists form an ideal in $`k[s]`$. The monic generator of this ideal is denoted by $`b_f(s)`$ and called the Bernstein-Sato polynomial of $`f.`$ The roots of $`b_f(s)`$ are rational, in particular, $`b(s)[s]`$, where $``$ is the field of rational numbers. A good introduction to $`D`$-modules may be found in \[B\].
The simplest characteristics of a polynomial $`f`$ are its degree $`d`$ and its number of variables $`n`$. This paper is motivated by the following natural question: what can one say about $`b_f(s)`$ in terms of $`n`$ and $`d`$? We give what may be regarded as a complete answer to this question. Namely, we describe an algorithm that for fixed $`n`$ and $`d`$ gives a complete list of all possible Bernstein-Sato polynomials and, for each polynomial $`b(s)`$ in this list, a complete description of the polynomials $`f`$ such that $`b_f(s)=b(s)`$.
Let $`P(n,d;k)`$ be the set of all the equivalence classes of the non-zero polynomials of degree at most $`d`$ in $`n`$ variables with coefficients in $`k`$ modulo the equivalence relation $`fgf=cg\text{ }`$ for some $`0ck`$. Note that $`b_f(s)=b_g(s)`$ if $`fg`$. We view $`P(n,d;k)`$ as the set of the $`k`$-rational points of the projective space $`(n,d;k)_k^{N1}`$ where $`N`$ is the number of monomials in $`n`$ variables of degree at most $`d`$. G. Lyubeznik \[L\] defined $`B(n,d)`$ as the set of all the Bernstein-Sato polynomials of all the polynomials from $`P(n,d;k)`$ as $`k`$ varies over all fields of characteristic $`0`$ and he proved that $`B(n,d)`$ is a finite set. He also asked if the subset of $`(n,d;k)`$ corresponding to a given element of $`B(n,d)`$ is constructible. In this paper we give an affirmative answer to Lyubeznik’s question. The constructible sets in question turn out to be definable over $``$, i.e. their defining equations and inequalities are the same for all fields $`k`$.
A crucial ingredient in our proof is the fact, very recently discovered by T. Oaku \[O\] that there is an algorithm that, given a polynomial $`f`$, returns its Bernstein-Sato polynomial $`b_f(s)`$. Using Oaku’s algorithm and our proof of the constructibility of the set of polynomials $`f`$ having a fixed $`b_f(s)`$ we have developed an algorithm for the *definitive computation* of the Bernstein-Sato polynomials for each pair $`(n,d)`$, by which we mean that our algorithm, given $`n`$ and $`d`$, returns the list of all the elements of $`B(n,d)`$ and for each $`b(s)B(n,d)`$, a finite number of locally closed sets $`V_i=V_i^{}V_i^{\prime \prime }`$, where $`V_i^{}`$ and $`V_i^{\prime \prime }`$ are Zariski closed subsets of $`(n,d;)`$ defined by explicit polynomial equations with rational coefficients, such that for every field $`k`$ of characteristic $`0`$, the subset of $`P(n,d;k)`$ having $`b(s)`$ as the Bernstein-Sato polynomial is the set of $`k`$-rational points of
$$S(b(s),k)=(_iV_i)_{}k(n,d;)_{}k=(n,d;k).$$
The definitive computation for fixed $`n`$ and $`d`$ could be useful in a number of ways. For example, it would produce an algorithm for the computation of $`b_f(s)`$ for $`fP(n,d;k)`$ that is likely to be considerably more efficient than all other currently available algorithms. It would also produce the smallest integer $`t`$ such that $`R_n(k)_f`$ is generated by $`\frac{1}{f^t}`$ as an $`A_n(k)`$-module for all $`f`$ of degree at most $`d`$ (this integer is denoted $`t(n,d)`$ in \[L\]). Moreover, using a similar technique we develop an algorithm for a *quasi-definitive computation* of the annihilator of $`\frac{1}{f^t}`$ in $`A_n(k)`$, which, provided $`t`$ is known, gives a presentation of $`R_n(k)_f`$ as an $`A_n(k)`$-module (see Example 2.4). We call it quasi-definitive because its output is not uniquely determined (see Remark 6.7). The last algorithm is particularly important for U. Walther’s algorithmic computation of local cohomology modules \[Wa\]. These applications are discussed in the next section. More applications will undoubtedly arise in the future.
The results of this paper are a part of my thesis. I would like to thank my advisor Gennady Lyubeznik for suggesting this problem to me.
## 2. Examples and Discussion
Our algorithms have been implemented as scripts written in the Macaulay 2 programming language (see \[M2\]). In this section we give some examples of actual computations and discuss possible uses of the results of computation.
###### Example 2.1.
If $`n=2`$ and $`d=2`$ then
$$f=a_{20}x^2+a_{11}xy+a_{02}y^2+a_{10}x+a_{01}y+a_{00},$$
so $`P(2,2;k)`$ is the set of the k-rational points of the projective space $`(2,2;k)=_k^5`$ with the homogeneous coordinate ring $`k[a_{ij}]`$, $`i,j=0,1,2`$. It takes our program less than 20 minutes on 300MHz Pentium-II machine to produce
$$B(2,2)=\{1,s+1,(s+1)^2,(s+1)(s+\frac{1}{2})\}$$
and give a description of the corresponding constructible sets of polynomials from $`B(2,2)`$ which is essentially equivalent to the following:
$``$ $`b_f(s)=1`$ iff $`fV_1=V_1^{}V_1^{\prime \prime }`$, where $`V_1^{}=V(a_{1,1},a_{0,1},a_{0,2},a_{1,0},a_{2,0})`$, while $`V_1^{\prime \prime }=V(a_{0,0})`$,
$``$ $`b_f(s)=s+1`$ iff $`fV_2=(V_2^{}V_2^{\prime \prime })(V_3^{}V_3^{\prime \prime })`$, where $`V_2^{}=V(0)`$, $`V_2^{\prime \prime }=V(\gamma _1)`$, $`V_3^{}=V(\gamma _2,\gamma _3,\gamma _4)`$, $`V_3^{\prime \prime }=V(\gamma _3,\gamma _4,\gamma _5,\gamma _6,\gamma _7,\gamma _8)`$,
$``$ $`b_f(s)=(s+1)^2`$ iff $`fV_4^{}V_4^{\prime \prime }`$, where $`V_4^{}=V(\gamma _1)`$, $`V_4^{\prime \prime }=V(\gamma _2,\gamma _3,\gamma _4)`$,
$``$ $`b_f(s)=(s+1)(s+\frac{1}{2})`$ iff $`fV_5^{}V_5^{\prime \prime }`$, where $`V_5^{}=V(\gamma _3,\gamma _4,\gamma _5,\gamma _6,\gamma _7,\gamma _8)`$, while $`V_5^{\prime \prime }=V(a_{1,1},a_{0,1},a_{0,2},a_{1,0},a_{2,0})`$,
where $`\gamma _i`$ may be looked up in this list:
$`\gamma _1=a_{0,2}a_{1,0}^2a_{0,1}a_{1,0}a_{1,1}+a_{0,0}a_{1,1}^2+a_{0,1}^2a_{2,0}4a_{0,0}a_{0,2}a_{2,0}`$,
$`\gamma _2=2a_{0,2}a_{1,0}a_{0,1}a_{1,1}`$,
$`\gamma _3=a_{1,0}a_{1,1}2a_{0,1}a_{2,0}`$,
$`\gamma _4=a_{1,1}^24a_{0,2}a_{2,0}`$,
$`\gamma _5=2a_{0,2}a_{1,0}a_{0,1}a_{1,1}`$,
$`\gamma _6=a_{0,1}^24a_{0,0}a_{0,2}`$,
$`\gamma _7=a_{0,1}a_{1,0}2a_{0,0}a_{1,1}`$,
$`\gamma _8=a_{1,0}^24a_{0,0}a_{2,0}`$.
It is not hard to see that this definitive computation agrees with the well-known result that $`b_f(s)=1`$ iff $`f`$ is constant, $`b_f(s)=s+1`$ iff $`f`$ is non-constant and non-singular, and $`b_f(s)=(s+1)^2`$ ( resp. $`b_f(s)=(s+1)(s+\frac{1}{2})`$ ) iff f can be reduced to $`xy`$ (resp. $`x^2`$) by a linear change of variables.
The definitive computation for fixed $`n`$ and $`d`$ is likely to lead to a considerably more efficient way of computing $`b_f(s)`$ for $`fP(n,d;k)`$. Namely, to compute $`b_f(s)`$ for a concrete polynomial $`f`$ one just has to “search the database”, i.e. check which of the constructible sets this polynomial belongs to. Since there are finitely many of them and each one is described by explicit equations and inequalities in the coefficients of $`f`$ and each $`f`$ belongs to a unique one, we get a straightforward algorithm for computing $`b_f(s)`$ for all $`fP(n,d;k)`$.
All other known algorithms for computing $`b_f(s)`$ involve Gr bner bases computations. Often $`b_f(s)`$ is not very big but its computation is enormous because of the "intermediate explosion" caused by the fact that Gr bner bases computations are very time and memory consuming. But the algorithm of “searching the database” does not involve any Gr bner bases at all! For this reason it is likely to be considerably more efficient in computing $`f`$ for $`fP(n,d;k)`$, especially if the field $`k`$ is the fraction field of some finitely generated $``$-algebra, so that ordinary arithmetic operations in $`k`$ and hence Gr bner bases computations are especially expensive.
Certainly the algorithm just described requires “setting up the database”. A definitive computation for $`n`$ and $`d`$ must be performed just once. This part may be done on a “powerful computer” (we have in mind implementing some parallel processing techniques) and the results of this computation may then be stored in a file accessible for “not-so-powerful” machines, which are capable of performing the “search the database” part. However a definitive computation even for rather small values of $`n`$ and $`d`$ with the modest computer resources at our disposal and with the current level of efficiency of our program faces its own "intermediate explosion" problem.
###### Example 2.2.
If $`n=2`$ and $`d=3`$ then
$`f`$ $`=`$ $`a_{3,0}x^3+a_{2,1}x^2y+a_{1,2}xy^2+a_{0,3}y^3`$
$`+`$ $`a_{2,0}x^2+a_{1,1}xy+a_{0,2}y^2+a_{1,0}x+a_{0,1}y+a_{0,0},`$
so $`P(2,3;k)`$ is the set of the k-rational points of $`(2,3;k)=_k^9`$ with the homogeneous coordinate ring that involves $`10`$ variables. Our program exhausts all available memory, 128Mb, of the computer after about 3 hours and stops without producing an answer. However, a somewhat creative use of our program enables us to give a complete list of all the elements of $`B(2,3)`$ (but not the explicit descriptions of the constructible sets corresponding to each element of $`B(2,3)`$):
Since for any nonsingular polynomial its Bernstein-Sato polynomial is equal to $`s+1`$, it remains to consider the case where our $`f(2,3;k)`$ possesses a singularity at some point $`(x_0,y_0)`$. Keeping in mind that the Bernstein-Sato polynomial is stable under any linear substitution of variables, we may get rid of its linear part via the substitution $`xxx_0`$, $`yyy_0`$, i.e. $`f`$ takes the form
$$f=(ax^3+bx^2y+cxy^2+dy^3)+(a^{}x^2+b^{}xy+c^{}y^2).$$
Now it is easy to see that by homogeneous linear transformation the quadratic part may be shaped to one of the forms $`0`$, $`xy`$, $`x^2`$. Therefore it is enough to compute the Bernstein-Sato polynomial for the following polynomials:
$`f_1`$ $`=`$ $`ax^3+bx^2y+cxy^2+dy^3,`$
$`f_2`$ $`=`$ $`(ax^3+bx^2y+cxy^2+dy^3)+xy,`$
$`f_3`$ $`=`$ $`(ax^3+bx^2y+cxy^2+dy^3)+x^2.`$
Our program returns the complete sets of possible Bernstein-Sato polynomials for $`f_1`$ in 22 minutes, for $`f_2`$ in 16 minutes and for $`f_3`$ in 21 minutes. Of course, in each of the three cases our program produces an explicit description of the corresponding constructible set in $`_k^3`$ (each of $`f_i`$ contains 4 indeterminate coefficients) for each element $`b(s)B_{f_i}`$. We omit these and list only the Bernstein-Sato polynomials:
$`B_{f_1}`$ $`=`$ $`\{(s+1)^2(s+{\displaystyle \frac{2}{3}})(s+{\displaystyle \frac{4}{3}}),`$
$`(s+1)^2(s+{\displaystyle \frac{1}{2}}),`$
$`(s+1)(s+{\displaystyle \frac{2}{3}})(s+{\displaystyle \frac{1}{3}})\};`$
$`B_{f_2}`$ $`=`$ $`\{(s+1)^2\};`$
$`B_{f_3}`$ $`=`$ $`\{(s+1)(s+{\displaystyle \frac{7}{6}})(s+{\displaystyle \frac{5}{6}}),`$
$`(s+1)^2(s+{\displaystyle \frac{3}{4}})(s+{\displaystyle \frac{5}{4}}),`$
$`(s+1)^2(s+{\displaystyle \frac{1}{2}}),`$
$`(s+1)(s+{\displaystyle \frac{1}{2}})\}.`$
Thus
$`B(2,3)`$ $`=`$ $`\{(s+1)^2(s+{\displaystyle \frac{2}{3}})(s+{\displaystyle \frac{4}{3}}),`$
$`(s+1)^2(s+{\displaystyle \frac{1}{2}}),`$
$`(s+1)(s+{\displaystyle \frac{2}{3}})(s+{\displaystyle \frac{1}{3}}),`$
$`(s+1)^2,`$
$`(s+1)(s+{\displaystyle \frac{7}{6}})(s+{\displaystyle \frac{5}{6}}),`$
$`(s+1)^2(s+{\displaystyle \frac{3}{4}})(s+{\displaystyle \frac{5}{4}}),`$
$`(s+1)(s+{\displaystyle \frac{1}{2}}),`$
$`s+1,`$
$`1\}.`$
As was mentioned above, only the efficiency of the algorithm and the current efficiency of computer hardware and software obstruct us from getting a complete description of the constructible sets that correspond to the polynomials above.
As was pointed out in \[L\], $`t(n,d)`$ (which is defined in the last paragraph of the preceding section) is the largest absolute value of all the negative integer roots of all the polynomials in $`B(n,d)`$. Thus we get
###### Corollary 2.3.
$`t(2,3)=1`$, i.e. if $`fR_2(k)`$ is of degree at most $`3`$, then $`\frac{1}{f}`$ generates $`R_2(k)_f`$ as an $`A_2(k)`$-module.
To compute the localization of $`R_n(k)`$ at a polynomial $`f0`$ one needs to compute $`\text{Ann}f^sA_n(k)[s]`$ and take $`N=\text{Ann}f^s|_{s=a}`$, where $`a`$ is the minimal integer root of $`b_f(s)`$. Then $`R_n(k)_f=A_n(k)/N`$ as an $`A_n(k)`$-module (see Section 5 below).
Using a technique similar to that for computing Bernstein-Sato polynomials, we constructed an algorithm for a quasi-definitive computation of $`\text{Ann}f^s`$ for all $`fP(n,d;k)`$. By this we mean an explicit subdivision of $`(n,d;k)`$ into a finite union of constructible subsets and for each such subset $`V,`$ an explicit finite set of elements $`\beta _1,\beta _2,\mathrm{}A_n(k)[a_{i_1\mathrm{}i_n}][s]`$ with $`i_1+\mathrm{}+i_nd`$, such that $`\text{Ann}(f^s)=(\beta _1^{},\beta _2^{},\mathrm{})`$ for every $`fV`$, where $`\beta _i^{}`$ is the image of $`\beta _i`$ under the specialization of the $`a_{i_1\mathrm{}i_n}`$ to the corresponding coefficients of $`f`$.
###### Example 2.4.
Here is what we got for $`P(2,2;k)`$ (See Example 2.1 for notation):
$``$ $`\text{Ann}(f^s)=(\beta _1,\beta _2,\beta _3)`$ if $`f(V_1^{}V_1^{\prime \prime })(V_2^{}(V_{2,1}^{\prime \prime }V_{2,2}^{\prime \prime }))`$, where $`V_1^{}=V(0)`$, $`V_1^{\prime \prime }=V(\gamma _1)`$, $`V_2^{}=V(\gamma _2,\gamma _3,\gamma _4)`$, $`V_{2,1}^{\prime \prime }=V(a_{1,1},a_{0,2},a_{0,1})`$ and $`V_{2,2}^{\prime \prime }=V(\gamma _2,\gamma _3,\gamma _4,\gamma _5,\gamma _6,\gamma _7)`$;
$``$ $`\text{Ann}(f^s)=(\beta _1,\beta _4)`$ if $`fV_3^{}V_3^{\prime \prime }`$, where $`V_3^{}=V(\gamma _1)`$, while $`V_3^{\prime \prime }=V(\gamma _2,\gamma _3,\gamma _4)`$;
$``$ $`\text{Ann}(f^s)=(\beta _5,\beta _6)`$ if $`fV_4^{}(V_{4,1}^{\prime \prime }V_{4,2}^{\prime \prime })`$, where $`V_4^{}=V(\gamma _2,\gamma _3,\gamma _4,\gamma _5,\gamma _6,\gamma _7)`$, $`V_{4,1}^{\prime \prime }=V(a_{1,0},a_{2,0},a_{1,1},\gamma _5)`$ and $`V_{4,2}^{\prime \prime }=(a_{1,1},a_{0,2},a_{0,1},\gamma _7)`$;
$``$ $`\text{Ann}(f^s)=(\beta _7,\beta _8)`$ if $`fV_5^{}V_5^{\prime \prime }`$, where $`V_5^{}=V(a_{1,0},a_{2,0},a_{1,1},\gamma _5)`$, while $`V_5^{\prime \prime }=V(a_{1,1},a_{0,1},a_{0,2},a_{1,0},a_{2,0})`$;
$``$ $`\text{Ann}(f^s)=(\beta _9,\beta _{10})`$ if $`fV_6^{}V_6^{\prime \prime }`$, where $`V_6^{}=V(a_{1,1},a_{0,2},a_{0,1},\gamma _7)`$, while $`V_5^{\prime \prime }=V(a_{1,1},a_{0,1},a_{0,2},a_{1,0},a_{2,0})`$;
$``$ $`\text{Ann}(f^s)=(\beta _9,\beta _{11})`$ if $`f(V_7^{}V_7^{\prime \prime })V_8^{}`$, where $`V_7^{}=V(a_{1,1},a_{0,2},a_{0,1})`$, $`V_7^{\prime \prime }=V(a_{1,1},a_{0,2},a_{0,1},\gamma _7)`$ and $`V_8^{}=V(a_{1,1},a_{0,1},a_{0,2},a_{1,0},a_{2,0})`$;
where the polynomials $`\beta _i`$ are listed below:
$`\beta _1=a_{1,1}x_1dx_1+2a_{0,2}x_2dx_12a_{2,0}x_1dx_2a_{1,1}x_2dx_2+a_{0,1}dx_1a_{1,0}dx_2`$,
$`\beta _2=a_{1,1}a_{2,0}x_1^2dx_1+a_{1,1}^2x_1x_2dx_1+a_{0,2}a_{1,1}x_2^2dx_12a_{2,0}^2x_1^2dx_22a_{1,1}a_{2,0}x_1x_2dx_2`$
$`2a_{0,2}a_{2,0}x_2^2dx_2a_{1,1}^2sx_2+4a_{0,2}a_{2,0}sx_2+a_{1,0}a_{1,1}x_1dx_1+a_{0,1}a_{1,1}x_2dx_1`$
$`2a_{1,0}a_{2,0}x_1dx_22a_{0,1}a_{2,0}x_2dx_2a_{1,0}a_{1,1}s+2a_{0,1}a_{2,0}s+a_{0,0}a_{1,1}dx_12a_{0,0}a_{2,0}dx_2`$,
$`\beta _3=a_{2,0}x_1^2dx_2+a_{1,1}x_1x_2dx_2+a_{0,2}x_2^2dx_2a_{1,1}sx_12a_{0,2}sx_2+a_{1,0}x_1dx_2`$
$`+a_{0,1}x_2dx_2a_{0,1}s+a_{0,0}dx_2`$,
$`\beta _4=a_{1,1}^2x_1dx_14a_{0,2}a_{2,0}x_1dx_1+a_{1,1}^2x_2dx_24a_{0,2}a_{2,0}x_2dx_22a_{1,1}^2s`$
$`+8a_{0,2}a_{2,0}s2a_{0,2}a_{1,0}dx_1+a_{0,1}a_{1,1}dx_1+a_{1,0}a_{1,1}dx_22a_{0,1}a_{2,0}dx_2`$,
$`\beta _5=a_{1,1}dx_12a_{2,0}dx_2`$,
$`\beta _6=2a_{2,0}x_1dx_2+a_{1,1}x_2dx_22a_{1,1}s+a_{1,0}dx_2`$,
$`\beta _7=dx_1`$,
$`\beta _8=2a_{0,2}x_2dx_24a_{0,2}s+a_{0,1}dx_2`$,
$`\beta _9=dx_2`$,
$`\beta _{10}=2a_{2,0}x_1dx_14a_{2,0}s+a_{1,0}dx_1`$,
$`\beta _{11}=a_{2,0}x_1^2dx_12a_{2,0}sx_1+a_{1,0}x_1dx_1a_{1,0}s+a_{0,0}dx_1`$,
and the polynomials $`\gamma _i`$ are in this list:
$`\gamma _1=a_{0,2}a_{1,0}^2a_{0,1}a_{1,0}a_{1,1}+a_{0,0}a_{1,1}^2+a_{0,1}^2a_{2,0}4a_{0,0}a_{0,2}a_{2,0}`$,
$`\gamma _2=2a_{0,2}a_{1,0}a_{0,1}a_{1,1}`$,
$`\gamma _3=a_{1,0}a_{1,1}2a_{0,1}a_{2,0}`$,
$`\gamma _4=a_{1,1}^24a_{0,2}a_{2,0}`$,
$`\gamma _5=a_{0,1}^24a_{0,0}a_{0,2}`$,
$`\gamma _6=a_{0,1}a_{1,0}2a_{0,0}a_{1,1}`$,
$`\gamma _7=a_{1,0}^24a_{0,0}a_{2,0}`$.
## 3. Constructible Sets
In this section we describe some of the properties of constructible sets that are used in the next section. We recall that a set is constructible iff it is a finite union of locally closed sets and a set is locally closed iff it is the difference of two closed sets.
###### Theorem 3.1.
Let $`C`$ be a constructible subset of a variety $`X`$. Then $`C`$ may be presented uniquely as a disjoint union $`_{i=1}^m(V_i^{}V_i^{\prime \prime })`$, where for all $`i`$ the sets $`V_i^{}`$ and $`V_i^{\prime \prime }`$ are closed, $`V_1^{}V_1^{\prime \prime }V_2^{}V_2^{\prime \prime }\mathrm{}V_m^{}V_m^{\prime \prime }`$ and no two consequent sets in this chain have an irreducible component in common. We call it a canonical presentation of $`C`$ as a union of locally closed subsets.
###### Proof.
Let $`d(C)`$ be the maximal dimension of an irreducible component in $`\overline{C}`$. The only possible choices for $`V_1^{}`$ and $`V_2^{}`$. Now $`V_1^{}=\overline{C}`$ and $`V_1^{\prime \prime }=\overline{V_1^{}C}`$ and let $`C_1=CV_1^{\prime \prime }`$. Note that $`d(C_1)<d(C)`$ and we may assume by induction on $`d`$ that the chain $`V_2^{}V_2^{\prime \prime }\mathrm{}V_m^{}V_m^{\prime \prime }`$ such that $`C^{}=_{i=2}^m(V_i^{}V_i^{\prime \prime })`$ exists and is unique. Then $`V_1^{}V_1^{\prime \prime }V_2^{}V_2^{\prime \prime }\mathrm{}V_m^{}V_m^{\prime \prime }`$ is the unique chain for $`C`$, which satisfies the condition in the statement. ∎
###### Remark 3.2.
There is an algorithmic way for constructing such a presentation, starting with $`C`$ presented as a union of nonempty sets $`W_\alpha (W_\alpha ^{(1)}\mathrm{}W_\alpha ^{(h_\alpha )})`$, where $`W_\alpha `$ and $`W_\alpha ^{(i)}`$ are closed irreducible subsets and $`W_\alpha W_\alpha ^{(i)}`$ for all $`i`$. Let $`d(C)=\mathrm{max}_\alpha dimW_\alpha `$ (which agrees with the definition in the proof of the theorem).
Let $`V_1^{}`$ be the union of all maximal elements in the set $`\{W_\alpha \}`$ and $`V_1^{\prime \prime }`$ be the union of all $`W_\alpha ^{(i)}`$ that are minimal with the following property: there is a set of pairs $`\{(\alpha _j,i_j)\}_{j=1}^l`$ such that $`W_{\alpha _1}`$ is a component of $`V_1^{}`$, $`W_{\alpha _l}^{(i_l)}=W_\alpha ^{(i)}`$ and $`W_{\alpha _j}^{(i_j)}W_{\alpha _j1}`$ for all $`j=2,\mathrm{},l`$ . Now $`d(C(V_1^{}V_1^{\prime \prime }))`$ is less than $`d(C)`$, therefore, we may assume again by induction on $`d`$ that we are able to construct the rest of $`V_i^{}`$ and $`V_i^{\prime \prime }`$.
###### Lemma 3.3.
Let $`X`$ be a variety and $`f:XY`$ a map into any finite set $`Y`$. Then $`f^1(y)`$ is constructible for every $`yY`$ iff for every closed irreducible subvariety $`X^{}X`$ there is an open $`UX^{}`$ such that $`f|_U`$ is a constant function.
###### Proof.
Assume the second part holds. Take any $`yY`$ and let $`Z=f^1(y)`$. Let $`n=dimX`$ and assume the lemma is proved for dimensions less then $`n`$. First of all, since $`X`$ is a finite union of its irreducible components, we may proceed assuming that $`X`$ is irreducible. Let $`U`$ be an open subset of $`X`$ such that $`f(u)=y^{}`$ for all $`uU`$. If $`y^{}y`$ then $`ZXU`$, which has dimension less than $`n`$ and, therefore, $`Z`$ is constructible by the induction assumption. If $`y=y^{}`$ then $`(ZU)(XU)`$ is constructible, hence so is $`Z=U(ZU)`$.
It remains to check the case $`dimX=0`$, in which $`X`$ is a finite set of points and is certainly constructible.
Conversely, assume that $`f^1(y)`$ is constructible for every $`yY`$. Let $`X^{}X`$ be a closed irreducible subvariety. Then $`X^{}=_{yY}(f^1(y)X^{})`$ and, since $`Y`$ is a finite set and $`X^{}`$ is irreducible, the closure of $`X_y^{}=f^1(y)X^{}`$ for some $`yY`$ is equal to $`X^{}`$. But $`X_y^{}`$ is constructible, hence it is open in its closure $`\overline{X_y^{}}=X^{}`$. ∎
## 4. Parametric Gr bner Bases
This section describes an approach to computing parametric Gr bner bases in Weyl algebras. A good source on computing Gr bner bases in non-commutative algebras is \[KR,We\]. For a discussion of parametric Gr bner bases, which leads to the notion of comprehensive Gr bner bases, see \[We\] for the commutative case and \[K,We\] for the case of solvable algebras. However, everything that is needed for this paper is stated and proved in this section.
Let $`C=k[\overline{a}]`$ ($`\overline{a}=\{a_1,\mathrm{},a_m\}`$) be the ring of parameters and $`R=C\overline{y},\overline{x},\overline{}`$ be the ring of non-commutative polynomials in $`\overline{y}=\{y_1,\mathrm{},y_l\}`$, $`\overline{x}=\{x_1,\mathrm{},x_n\}`$ and $`\overline{}=\{_1,\mathrm{},_n\}`$ with coefficients in $`C`$, where $`\overline{x}`$ and $`\overline{}`$ satisfy the same relations as in a Weyl algebra and $`\overline{y}`$ is contained in the center of $`R`$.
###### Definition 4.1.
For a prime $`P`$ in $`C`$, we shall call the natural map $`Ck(P)`$ as well as the induced map $`R=C\overline{y},\overline{x},\overline{}k(P)\overline{y},\overline{x},\overline{}`$, where $`k(P)`$ is the residue field at $`P`$, the *specialization* at the point $`P`$ and denote both maps by $`\sigma _P`$.
The next result is similar to Oaku’s Proposition 7 in \[O\].
Let $`<`$ be an order on monomials in $`a`$, $`y`$, $`x`$ and $``$ such that every $`a_i`$ is $`<<`$ than any of $`x_j`$, $`y_j`$ or $`_j`$ (i.e. the order $`<`$ eliminates $`x_j`$, $`y_j`$ and $`_j`$). Assume $`G`$ is a finite Gr bner basis in $`R`$, then we claim that $`\sigma _P(G)=\{\sigma _P(g)|gG\}`$ is a Gr bner basis in $`\sigma _P(R)`$ for “almost” every $`P\text{Spec}C`$. Namely,
###### Lemma 4.2.
For any $`G`$ there exists a polynomial $`hC`$ such that $`\sigma _P(G)`$ is a Gr bner basis for every $`P`$ not containing $`h`$.
###### Proof.
We need to make some definitions. For a polynomial $`f`$ let $`inM(f)`$ be the initial monomial $`inC(f)`$ the initial coefficient such that $`in(f)=inC(f)inM(f)`$ the initial term of $`f`$. Also for $`fR`$ let $`inM_{}(f)\overline{y},\overline{x},\overline{}`$ and $`inC_{}(f)C`$ be the initial monomial and the initial coefficient of $`f`$ viewed as a polynomial in $`x,y,`$ with coefficients in $`C`$ with respect to $``$, the restriction of $`<`$ to $`\overline{y},\overline{x},\overline{}`$.
One obvious observation is that a specialization $`\sigma _P:(R,<)(\sigma _P(R),)`$ preserves the order.
Let $`h=_{gG}inC_{}(g)C`$. Consider any $`P\text{Spec}C`$ not containing $`h`$. Take a polynomial $`f^{}`$ in ideal of $`\sigma _P(R)`$ generated by $`\sigma _P(G)`$, then there is $`f`$ such that $`f^{}=\sigma _P(f)`$ and $`inM_{}(f)=inM(f^{})`$. Since $`G`$ is a Gr bner basis in $`R`$, we have $`inM(g)|inM(f)`$ for some $`gG`$, which means that $`inM_{}(g)|inM_{}(f)`$. Now, $`inM(\sigma _P(g))=inM_{}(g)`$, because $`inC_{}(g)P`$. Thus $`inM(\sigma _P(g))|inM(\sigma _P(f))`$, which proves that $`\sigma _P(G)`$ is a Gr bner basis. ∎
###### Remark 4.3.
The statement of the lemma is true for reduced Gr bner bases as well.
The lemma leads to the following
###### Algorithm 4.4.
| Input: | $`F^{}`$: | a finite set of generators for a prime ideal $`QC`$. |
| --- | --- | --- |
| | $`F`$: | a finite set of generators of a left ideal $`IR`$ containing $`QA_n`$, |
| Output: | $`G`$: | a (reduced) Gr bner basis in $`R`$ with respect to $`<`$, |
| | $`h`$: | a polynomial in $`C`$, which we shall call an *exceptional polynomial*, |
| | | such that for any $`P\text{Spec}(k[a_1,\mathrm{},a_m])`$, $`PQ`$ and $`hP`$ |
| | | the ideal $`\sigma _P(I)\sigma _P(R)`$ has a $`\sigma _P(G)`$ as a (reduced) Gr bner |
| | | basis with respect to $``$. |
1. Compute a Gr bner basis $`G`$ of $`I+QR`$ (which is generated by $`FF^{}`$) .
2. Return $`G`$ and $`h=_{gGQ}inC_{}(g)`$.
###### Remark 4.5.
If all polynomials in $`F^{}`$ and all $`C`$-coefficients of all elements of $`F^{}`$ are homogeneous, then so is the exceptional polynomial $`h`$.
## 5. Oaku’s Algorithm
The original algorithm of T.Oaku for computing the Bernstein-Sato polynomial appeared in \[O\]. However there exist several modifications of the algorithm (see \[S,S,T\] for example). For our needs a version of the algorithm described in \[Wa\] will be utilized.
Let $`fR_n(k)`$. Denote by $`\text{Ann}f^s`$ the ideal of all elements in $`A_n(k)[s]`$ annihilating $`f^s`$. The following algorithm is Algorithm 4.4. from \[Wa\] with $`L=(_1,\mathrm{},_n)`$.
###### Algorithm 5.1.
| Input: | $`f`$: | a polynomial in $`R_n(k)`$ , |
| --- | --- | --- |
| Output: | $`\{P_j^{}\}`$: | generators of $`\text{Ann}f^s`$ |
1. Set $`Q=\{_i+\frac{df}{dx_i}_t,t\}`$.
2. Homogenize all $`q_iQ`$ using the new variable $`y_1`$ with respect to the weight $`w`$, where $`w(t)=w(y_1)=1`$, $`w(_t)=w(y_2)=1`$, $`w(x_i)=w(_i)=0`$. Denote the homogenized elements $`q_i^h`$.
3. Compute a Gr bner basis for the ideal generated by $`q_1^h,\mathrm{},q_r^h,\mathrm{\hspace{0.17em}1}y_1y_2`$ in $`A_{n+1}[y_1,y_2]`$ with respect to an order eliminating $`y_1,y_2`$.
4. Select the operators $`\{p_j\}_1^b`$ in this basis which do not contain $`y_1,y_2`$.
5. For each $`p_j`$, if $`w(p_j)>0`$ then replace $`p_j`$ by $`p_j^{}=_t^{w(P_j)}p_j`$ else replace $`p_j`$ by $`p_j^{}=t^{w(P_j)}p_j`$.
6. Return the operators $`\{p_j^{}\}_1^b`$.
The following is Algorithm 4.6 in \[Wa\].
###### Algorithm 5.2.
| Input: | $`f:`$ | o polynomial in $`R_n(k)`$, |
| --- | --- | --- |
| Output: | $`b_f(s)`$ | the Bernstein-Sato polynomial of $`f`$. |
1. Determine $`\text{Ann}f^s`$ following Algorithm 5.1.
2. Find a reduced Gr bner basis for the ideal $`\text{Ann}f^s+A_n[s]f`$ using an order that eliminates $`x`$ and $``$.
3. Return the unique element in the basis contained in $`k[s]`$.
## 6. The Main Results
Consider $`(n,d;k)`$ with the coordinate ring $`C=k[\overline{a}]`$, where $`\overline{a}=\{a_\alpha :|\alpha |d\}`$. Let $`f=_{|\alpha |d}a_\alpha x^\alpha `$.
###### Definition 6.1.
Let $`b(s)B(n,d)`$. The corresponding set $`S(b(s),k)(n,d;k)`$ in $`𝔸_k^N`$ is defined as the set of all the points $`P(n,d;k)`$ such that $`b_{\sigma _P(f)}(s)=b(s)`$. (We view points in $`(n,d;k)`$ as homogeneous primes in $`C`$. See Definition 4.1 for $`\sigma _P(f)`$.)
Let $`Q`$ be a homogeneous prime in $`C`$. Then $`\sigma _Q(f)`$ is a polynomial with coefficients in a field, hence $`b_{f_Q}(s)`$ may be computed. What would happen if we run Algorithm 5.2 trying to compute $`b_{f_Q}(s)`$ “lifting from $`k(Q)`$, the fraction field of $`C/Q`$, to $`C`$” every single step of the algorithm? Notice that $`\sigma _Q:Ck(Q)`$ has $`C/Q`$ as its image. Since the steps of the algorithm that do not involve Gr bner bases computation do not involve division either, we have to worry only about the two steps that deal with Gr bner bases. Suppose for these two steps we used 4.4 with $`F^{}`$ is a set generating $`Q`$, in particular we obtained the exceptional polynomials $`h_1`$ and $`h_2`$. Set $`h=h_1h_2`$, then the output, which is going to be $`b_{\sigma _Q(f)}(s)`$, is also the Bernstein-Sato polynomial of $`\sigma _P(f)`$ for every $`PQ`$ such that $`hP`$. Thus we have
###### Algorithm 6.2.
| Input: | $`f`$: | a polynomial in $`R_n(C)`$, |
| --- | --- | --- |
| | $`F^{}`$: | generators of a homogeneous prime ideal, |
| Output: | $`b(s)`$: | a polynomial in $`[s]`$, |
| | $`H`$: | generators of a homogeneous ideal in $`C`$ such that |
| | | $`b(s)=b_{\sigma _P(f)}(s)`$ for every point $`PV^{}\backslash V^{\prime \prime }`$, |
| | | where $`V^{}=V(F^{})`$ and $`V^{\prime \prime }=V(H)`$ ($`V^{\prime \prime }V^{}(n,d;k)`$). |
1. Compute the polynomial $`b(s)`$ and the exceptional polynomial $`h`$ as described above.
2. Return $`b(s)`$ and $`\{h\}F^{}`$.
###### Remark 6.3.
If we consider $`C^{}=Ck^{}`$ and $`f1R_n(C^{})`$, where $`k^{}`$ is an extension of $`k`$, then $`b(s)`$ is the Bernstein-Sato polynomial for any point in $`(V^{}_kk^{})(V^{\prime \prime }_kk^{})`$.
The next theorem gives an affirmative answer to Lyubeznik’s question about the constructibility of the set $`S(b(s),k)`$ of Definition 6.1.
###### Theorem 6.4.
The set $`S(b(s),k)`$ is constructible for every $`b(s)`$.
###### Proof.
The proof follows from the above algorithm. For the function $`\varphi :(n,d;k)B(n,d)`$, $`\varphi (P)=b_{\sigma _P(f)}(s)`$ the following is true. For every projective $`V^{}(n,d;k)`$ there is an open set $`U=V^{}V^{\prime \prime }V^{}`$ such that $`f|_U`$ is a constant function. Therefore we may apply Lemma 3.3. ∎
Algorithm 6.2 leads to the main algorithm and theorem of the paper.
###### Algorithm 6.5.
Input: $`n,d`$.
Output: The set of pairs $`L=\{(b(s),S(b(s)))|b(s)B(n,d)\}`$, where $`S(b(s))=S(b(s),)(n,d;)`$.
* Set $`L:=\mathrm{}`$, $`f:=_{|\alpha |d}a_\alpha x^\alpha `$ .
* Define the recursive procedure BSP($`Q`$), where $`Q\text{Spec}([\overline{a}])`$.
| BSP($`Q`$) := { |
| --- |
| Apply Algorithm 6.2 to $`V(Q)`$ and $`f`$ |
| to get an ideal $`I`$ in $`C`$ and $`b(s)[s]`$; |
| IF there is a pair $`(b(s),S)L`$ |
| THEN replace it by $`(b(s),S(V(Q)V(I)))`$ |
| ELSE $`L:=L\{(b(s),V(Q)V(I))\}`$; |
| IF $`V(I)\mathrm{}`$ THEN { |
| Find the minimal primes $`\{Q_i\}`$ associated to $`I`$; |
| FOR each $`Q_i`$ DO BSP($`Q_i`$) ; |
| } |
| } |
* Run BSP($`0`$).
###### Remark.
This algorithm returns some presentations for constructible sets $`S(b(s),)`$, the canonical presentations for which may be obtained by using the algorithm discussed in Remark 3.2.
###### Corollary 6.6.
The set $`S(b(s),k)`$ is defined over $``$, i.e. there exist ideals $`I_i[\overline{a}]`$ and $`J_i[\overline{a}]`$ $`(i=1,\mathrm{},m)`$ such that for any field $`k`$
$$S(b(s),k)=\underset{i}{}(V_i^{}V_i^{\prime \prime }),$$
where $`V_i^{}=V(k[\overline{a}]I_i)`$ is the zero set of the extension of $`I_i`$ and $`V_i^{\prime \prime }=V(k[\overline{a}]J_i)`$ is the zero set of the extension of $`J_i`$.
###### Proof.
Follows from the algorithm and Remark 6.3.
The annihilators $`\text{Ann}(f^s)`$ are computed using Algorithm 5.1 and the same technique as in the algorithm above. The output is a set of pairs $`\{(I_i,V_i)\}`$, where $`I_i`$ are the ideals in $`A_n(k)[\overline{a}][s]`$ and $`V_i`$ are locally closed sets, such that for any polynomial $`f`$ with coefficients in $`k`$ that corresponds to a point $`PV_i`$ the ideal $`\text{Ann}(f^s)`$ equals $`\sigma _P(I_i)`$, the ideal $`I_i`$ specialized to $`P`$.
After doing the above steps, the real life algorithm that produces Example 2.4 compresses its output in the following way. If $`(I_i,V_i)`$ and $`(I_j,V_j)`$ are two different pairs such that $`\sigma _P(I_i)=\sigma _P(I_j)`$ for all $`PV_j`$ then these two are replaced by the pair $`(I_i,V_iV_j)`$.
###### Remark 6.7.
The stratification of the parameter space produced by such computation is not unique. So we have to use prefix “quasi“ in “quasi-definitive computation”, because the annihilators, as opposed to Bernstein-Sato polynomials, depend on the parameters making it possible to slice the space of parameters in many ways.
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# Eleven dimensional supergravity as a constrained topological field theory
## 1 Introduction
Eleven dimensional supergravity has been an object of fascination since it was first discovered more than twenty years ago. It is the highest dimension in which a supergravity theory exists, and it contains all of the lower dimensional supergravities as dimensional reductions. It has also been proposed as a description of the classical limit of a phase of $``$ theory, which is not itself described by any consistent string theory. Anyone wanting to understand the dynamical structure of the different supergravities, and their inter-relations with each other under reductions and various kinds of duality transformations are recommended to start with the $`11`$ dimensional theory.
It is thus very interesting to wonder whether $`11`$ dimensional supergravity might be made into a quantum theory directly, using some non-perturbative approach. It is true that it is non-renormalizable in perturbation theory, but so are general relativity and the $`N=1,2`$ supergravities in four dimensions, and this has not prevented a great deal of progress being made understanding the exact structure of their canonical and path integral quantum theories-. A beautiful and robust structure was found that could not have been seen in perturbation theory, which leads to definite physical predictions such as the discrete spectra of the area and volume operators. The eigenstates of these operators provide a basis of states for background independent quantum theories of gravity, which are the spin network basis. These are constructed in terms of simple combinatorial and representation theory data, are closely related to the fusion algebras of conformal field theory and exist also for supersymmetric theories.
It is also clear now that these theories have rigorously defined hamiltonian and path integral formulations and so do exist as examples of diffeomorphism invariant quantum field theories. This existence is independent of the question of whether they have classical limits which reproduce classical general relativity or supergravity. This is presently an open question for general relativity and supergravity in $`d=4`$. Even if the particular dynamics of quantum general relativity in $`d=4`$ leads to a theory which is sensible at the Planck scale but lacks a good classical limit, it is still reasonable to conjecture that the robust and theory independent features of the background independent formulation of quantum gravity uncovered in loop quantum gravity will play a role in a background independent formulation of $``$ theory.
There is a simple reason to expect that the classical limit will exist for a background independent of $`11`$ dimensional supergravity, but not for $`d=3+1`$ general relativity, which is that whenever such a limit exists one expects a sensible perturbation theory to exist around any classical background that arises in a classical limit. There is no such consistent perturbation theory for pure four dimensional general relativity but there is one for many, if not all, consistent compactifications of $`11`$ dimensional supergravity-these are of course exactly the perturbative string theories. In fact, there is a simple argument that where such a limit exists the weakly coupled excitations will be described by a string theory.
It is then sensible to suppose that one way to uncover the structure of the states and operators that will go into the background independent formulation of $``$ theory is to make an exact canonical quantization of $`11`$ dimensional supergravity and discover what structures play the role of the spin networks and super-spin networks in $`D=3+1`$. The first step in any such attempt must be to have a suitable form for the action. What is required to make progress in a non-perturbative approach to quantization is to have a first order polynomial form of the action, which has the property that it arises from a topological field theory by the imposition of a set of constraints. The reason is it is by now understood that all the beautiful results which have followed from the use of the Ashtekar-Sen and related variables are due ultimately to the fact that in $`4`$ dimensions general relativity and supergravity (at least up to $`N=2`$) can be understood as arising in this way from topological field theories<sup>1</sup><sup>1</sup>1Indeed, such a form of the theory was understood earlier, by Plebanski. The significance of this form was only realized after Sen discovered the equivalent Hamiltonian variables, in an attempt to understand the canonical structure of supergravity. This was then formalized by Ashtekar for the canonical theory and in for the lagrangian theory. Attempts to relate this form to topological field theory then led to a rediscovery of Plebanski’s action and its extension to supergravity.. A further reason for taking this route is that in the presence of appropriate boundary conditions it leads to holographic formulations of quantum general relativity and supergravity in which the Bekenstein bound is realized naturally because the boundary theories are built from the state spaces of Chern-Simons theory<sup>2</sup><sup>2</sup>2This method has also been used recently to construct an explicit description of the quantum geometry of the black hole horizon..
The main goal of this paper is then to present a formulation of the $`11`$ dimensional supergravity as a constrained topological field theory. In the next section we introduce an $`11`$ dimensional topological field theory which has the property that its algebra of first class constraints, which generate the gauge transformations of the theory, reproduces exactly the $`11`$ dimensional super-Poincare algebra, including its central charges. The theory is introduced in the next section and the canonical analysis is given in section 4. In section 3 we show how the eleven dimensional supergravity action arises by imposing a certain set of constraints on the topological field theory. We then arrive at the supergravity action in the form given by Fre and d’Auria and Fre. A feature of that form of the theory, which plays an important role as well in our formulation, is the presence of two abelian potentials, which are a six form and a three form. Although we are not completely certain of the correct way to express it, we believe it likely that the formulation of a gravitational theory as a constrained topological field theory is closely related to the idea of formulating it in terms of free differential algebras, which was pursued by Fre and collaborators<sup>3</sup><sup>3</sup>3Another approach to the relationship between topological field theory and supergravity is described in ..
This paper represents only the first step of a program of quantization of $`11`$ dimensional supergravity. Still to be investigated is the implications for the canonical structure of the full $`11`$ dimensional supergravity action and the possible existence of a quantization using the methods of loop quantum gravity.
In the next section we introduce the topological quantum field theory in $`11`$ dimensions and in section 4 we derive its canonical formalism and compute its constraint algebra. We find that it is first class and that the algebra of constraints does reproduce weakly<sup>4</sup><sup>4</sup>4That is up to the constraints that say that the curvatures vanish in a topological field theory the super-Poincare algebra in $`11`$ dimensions with central charges. Along the way, in section 3, we show how constraints and lagrange multiplier terms may be added to the topological field theory to arrive at the full $`11`$ dimensional supergravity action, in the form given by Fre.
## 2 TQFT for the $`11`$ dimensional super-Poincare algebra
The first step in our construction is to find the TQFT whose algebra of constraints reproduces the supersymmetry algebra of supergravity in $`11`$ dimensions, including the central extensions. This algebra has the form<sup>5</sup><sup>5</sup>5We list the convention used in this paper. $`(i)`$ Indices $`\mu ,\nu ,\mathrm{}=0,1,\mathrm{},10`$ are spacetime indices while $`i,j,\mathrm{}=1,2,\mathrm{},10`$ are used for spatial indices; $`(ii)`$ The capital $`A,B,\mathrm{}=1,2,\mathrm{},32`$ stand for the $`Sp(32)`$ spinor indices; $`(iii)`$ $`a,b,\mathrm{},=0,1,\mathrm{},10`$ represent $`SO(10,1)`$ indices. We sometimes use a condensed index notation in which $`a_p=a_1a_2\mathrm{}a_p`$.
$$\{Q^A,Q^B\}=\mathrm{\Gamma }_a^{AB}G^a+\mathrm{\Gamma }_{ab}^{AB}G^{ab}+\mathrm{\Gamma }_{a_5}^{AB}G^{a_5},$$
(1)
where
$$\mathrm{\Gamma }_{a_p}^{AB}:=\mathrm{\Gamma }_{[a_1}\mathrm{\Gamma }_{a_2}\mathrm{}\mathrm{\Gamma }_{a_p]}^{AB}.$$
(2)
Here $`G^a`$ must be the generator of spacetime diffeomorphisms, so that $`G^0`$ must be related to the Hamiltonian and $`G^i`$ to the diffeomorphism constraints of the theory, while $`G^{ab}`$ and $`G^{a_5}`$ are the central extensions. This can be understood to be a contraction of $`Osp(1|32)`$.
One place to begin is with the fields of a gauge theory for the superalgebra $`Osp(1|32)`$ in $`10+1`$ dimensions. The generators of the superalgebra consist of the translation generator $`G_a`$, Lorentz generator $`J_{ab}`$, supersymmetry generators $`Q^A`$ and five-index antisymmetric generator $`G^{[abcde]}G^{a_5}`$. One starting point would be to define a 1-form superconnection associated with those generators:
$$𝒜_\mu :=A_\mu ^{ab}J_{ab}+e_\mu ^aG_a+\Psi _\mu ^AQ_A+A_\mu ^{a_\mathit{5}}G_{a_\mathit{5}},$$
(3)
where $`\mathrm{\Psi }^A`$ are $`32`$ component Majorana spinors. We may then introduce a $`BF`$ action as
$$=_M𝑑x^{\mathit{11}},$$
(4)
where $``$ is a super nine form and $`F`$ is the curvature of $`Osp(1|32)`$ defined by
$$=d𝒜+𝒜𝒜.$$
(5)
Unfortunately this route seems not to lead to the standard $`11`$ dimensional supergravity<sup>6</sup><sup>6</sup>6But it is of interest and has been pursued in .. The difference seems to be that the relevant gauge group for 11 dimensional supergravity, at least at the classical level, is the Super-poincare group, which is a contraction of $`Osp(1|32)`$. As a result, the central generators are realized in supergravity by functionals of a three form abelian gauge field rather than by a standard component of a connection one-form. The basic mystery of the construction of $`11`$ dimensional supergravity (as well as many of its lower dimensional reductions) is how such a field may be seen to arise from a gauge theoretic structure as does the connection of spacetime. From the point of view in which gravitational theories are understood as constrained topological field theories, this mystery can be solved, for topological field theories can indeed be constructed based on gauge theories of Abelian $`p`$-forms and they can be quantized using the methods of loop quantum gravity. Indeed as we shall see here, one can construct a topological field theory whose constraint algebra realizes perfectly the full super-Poincare algebra, with the abelian $`p`$ forms realizing the central charges<sup>7</sup><sup>7</sup>7Thus, while we solve the problem of encoding the central charges in an algebra of canonical constraints, in a way that leads under suitable constraints to $`11`$ dimensional supergravity, we do not solve the problem of what all this may have to do with $`Osp(1|32)`$..
In order to achieve this we have found it necessary to introduce not only the three form gauge field $`a_{\mu \nu \rho }`$ but its dual, which is a six form field $`b_{\alpha \beta \gamma \delta ϵ\varphi }`$. This leads us to a formalism which is similar to that of D’Auria and Fre. Indeed, as our results show, there is likely a close relationship between their conception of supergravity based on Cartan integrable algebras and the more recent conception of a gravitational theory as a constrained topological field theory.
We now introduce our $`11`$ dimensional super-Poincare topological field theory.
To begin with we define the topological field theory associated to the $`11`$ dimensional super-Poincare algebra, unextended by central charges. The gauge field is then of the form,
$$𝒜_\mu :=A_\mu ^{ab}J_{ab}+e_\mu ^aG_a+\mathrm{\Psi }_\mu ^AQ_A.$$
(6)
The components of the curvature are given by
$$F^{ab}=dA^{ab}A^{ac}A_c^b,$$
(7)
$$F^a=de^aA^{ab}e_b\frac{i}{2}\mathrm{\Psi }_A\mathrm{\Gamma }_{}^{aA}{}_{B}{}^{}\mathrm{\Psi }^B,$$
(8)
$$F^A=D\mathrm{\Psi }^A=d\mathrm{\Psi }^A\frac{1}{4}A^{ab}\mathrm{\Gamma }_{abB}^A\mathrm{\Psi }^B.$$
(9)
To represent the central charges we introduce the three form $`a_{\mu _3}`$ and its dual which is a 6-form field $`b_{\mu _6}`$. We introduce their curvatures,
$$F^{}=db15daa\frac{ı}{2}\rho ^5,$$
(10)
$$F^{\mathrm{}}=da\frac{1}{2}\rho ^2,$$
(11)
where
$$\rho _{a_p}^p:=\mathrm{\Psi }_A\mathrm{\Gamma }_{}^{a_pA}{}_{B}{}^{}\mathrm{\Psi }^BE_{a_p}^p,$$
(12)
and
$$E_{a_p}^p=e_{a_1}e_{a_2}\mathrm{}e_{a_p}.$$
(13)
Using these curvatures we then write a TQFT,
$$I^{TQFT}=\frac{1}{g^2}B_{ab}F^{ab}+B_aF^a+B_AF^A+B^{\mathrm{}}F^{\mathrm{}}+B^{}F^{}.$$
(14)
The $`B`$’s are lagrange multipliers which have the form degree indicated. The field equations are
$$F^{ab}=F^a=F^A=F^{\mathrm{}}=F^{}=0,$$
(15)
and
$$𝒟B^{ab}e^{[a}B^{b]}\frac{1}{4}\mathrm{\Psi }_A\mathrm{\Gamma }_{}^{abA}{}_{B}{}^{}B^B=0,$$
(16)
$$𝒟B^a\mathrm{\Psi }_A\mathrm{\Gamma }_{}^{abA}{}_{B}{}^{}\mathrm{\Psi }^Be^bB^{\mathrm{}}\frac{5i}{2}\mathrm{\Psi }_A\mathrm{\Gamma }_{}^{ab_4A}{}_{B}{}^{}\mathrm{\Psi }^Be_{b_4}B^{}=0,$$
(17)
$$𝒟B^Ai\mathrm{\Gamma }_{aB}^A\mathrm{\Psi }^BB^a\mathrm{\Gamma }_{abB}^A\mathrm{\Psi }^BE^{ab}B^{\mathrm{}}i\mathrm{\Gamma }_{a_5B}^A\mathrm{\Psi }^BE^{a_5}B^{}=0,$$
(18)
$$dB^{}=0,$$
(19)
$$dB^{\mathrm{}}30daB^{}+15adB^{}=0.$$
(20)
In section (4) we will discuss the canonical formulation of the action (14) and show that its algebra of first class constraints replicates the superalgebra (1).
## 3 Constraining the $`TQFT`$ to get supergravity in $`11`$ dimensions
We obtain an action for $`11`$ dimensional supergravity by adding constraint and lagrange multiplier terms,
$$I_{11D}^{SUGRA}=I^{TQFT}+I^{CONST.}+I^{F_4},$$
(21)
where
$`I^{CONST.}`$ $`=`$ $`{\displaystyle \lambda _{ab}}(B^{ab}{\displaystyle \frac{1}{9l^9}}E^{ab})+\lambda _A(B^A{\displaystyle \frac{2}{l^8}}\mathrm{\Gamma }_{a_8B}^A\mathrm{\Psi }^BE^{a_8})`$ (22)
$`+\lambda _a(B^a{\displaystyle \frac{7ı}{30}}E_7^{ab_6}\mathrm{\Psi }_A\mathrm{\Gamma }_B^{c_5A}\mathrm{\Psi }^Bϵ_{b_6c_5})`$
$`+\lambda _{}(B^{}56da)+\lambda _{\mathrm{}}(B^{\mathrm{}}+56ı\rho ^5),`$
and
$$I^{F_4}=2F_{a_4}R^{\mathrm{}}E_{b_7}^7ϵ^{a_4b_7}+\frac{1}{330}F_{a_4}F^{a_4}E_{11}^{}.$$
(23)
It is not difficult to see that the variations of the lagrange multipliers $`\lambda `$ reproduce the $`D=11`$ supergravity action in the form given by Fre.
$`I^{SG}`$ $`=`$ $`{\displaystyle \frac{1}{g^2}}{\displaystyle \frac{1}{9l^9}E^{ab}}F^{ab}+{\displaystyle \frac{7ı}{30}}E_7^{ab_6}\mathrm{\Psi }_A\mathrm{\Gamma }_B^{c_5A}\mathrm{\Psi }^Bϵ_{b_6c_5}F^a+{\displaystyle \frac{2}{l^8}}\mathrm{\Gamma }_{a_8B}^A\mathrm{\Psi }^BE^{a_8}F^A`$ (24)
$`+56ı\rho ^5F^{\mathrm{}}+56daF^{}`$
$`2F_{a_4}R^{\mathrm{}}E_{b_7}^7ϵ^{a_4b_7}+{\displaystyle \frac{1}{330}}F_{a_4}F^{a_4}E_{11}^{}.`$
Elimination of the lagrange multipliers then leads to the theory in the original form.
We not that the lagrange multiplier terms $`I^{F_4}`$ are necessary to get the $`da^2`$ terms in the supergravity action. Were they absent the $`a`$ field would have dynamics only from the Chern-Simons like terms $`adada`$. One interesting question that this approach should be able to answer is why supersymmetry requires both the Maxwell and Chern-Simons like terms in the supergravity action.
## 4 Canonical formulation of the $`11D`$ TQFT
We now describe to the canonical decomposition of the TQFT given by (14). Our main goal here is to find the algebra of its constraints. We assume that the eleven dimensional spacetime $`^{11}`$ has the form $`^{11}=\mathrm{\Sigma }^{10}\times R`$ where $`\mathrm{\Sigma }^{10}`$ is a compact ten dimensional manifold. We then make a $`10+1`$ decomposition to find that<sup>8</sup><sup>8</sup>8Note that the decomposition is much simpler than in the standard case as there is no metric and hence no lapse and shift. The decomposition is purely a matter of pulling back forms.
$`I^{TQFT}`$ $`=`$ $`{\displaystyle }dt{\displaystyle _{\mathrm{\Sigma }^{10}}}d^{10}x\{{\displaystyle \frac{1}{2}}\pi ^{iab}\dot{A}_{iab}+\pi ^{ia}\dot{e}_{ia}+\pi ^{iA}\dot{\mathrm{\Psi }}_{iA}+{\displaystyle \frac{1}{6!}}p^{i_6}\dot{b}_{i_6}+{\displaystyle \frac{1}{3!}}R^{i_3}\dot{a}_{i_3}`$ (25)
$`+B_0^{ab}F_{ab}+B_0^AF_A+B_0^aF_a+B_0^{\mathrm{}}F^{\mathrm{}}+B_0^{}F^{}`$
$`+A_{0ab}J^{ab}+\mathrm{\Psi }_{0A}Q^A+e_{0a}G^a+a_{0ij}G_2^{ij}+b_{0i_5}G_5^{i_5}\},`$
where now all forms are in the $`10`$ dimensional space.
The expressions for the canonical momenta are,
$`\pi _{ab}^i`$ $`=`$ $`{\displaystyle \frac{2}{g^2}}(B_{ab}^{})^i,`$
$`\pi _a^i`$ $`=`$ $`{\displaystyle \frac{1}{g^2}}(B_a^{})^i,`$
$`\pi _A^i`$ $`=`$ $`{\displaystyle \frac{1}{g^2}}(B_A^{})^i,`$
$`p^{i_6}`$ $`=`$ $`{\displaystyle \frac{1}{g^2}}(B_{}^{})^{i_6},`$
$`R^{i_3}`$ $`=`$ $`{\displaystyle \frac{1}{g^2}}(B_{\mathrm{}}^{})^{i_3}{\displaystyle \frac{15}{3!}}p^{i_3j_3}a_{j_3}.`$ (26)
The constraints take the form
$`Q^A`$ $`=`$ $`D_i\pi ^{iA}ı\pi ^{ia}\mathrm{\Gamma }_{aB}^A\mathrm{\Psi }_i^B\mathrm{\Psi }_i^B\left\{\mathrm{\Gamma }_{abB}^A(R^{ijk}+{\displaystyle \frac{15}{3!}}p^{ijkl_3}a_{l_3})E_{jk}^{ab}ı\mathrm{\Gamma }_{a_5B}^Ap^{ik_5}E_{k_5}^{b_5}\right\},`$ (27)
$`G^a`$ $`=`$ $`D_k\pi ^{ka}\mathrm{\Psi }_i^A\mathrm{\Psi }_j^B\left\{\mathrm{\Gamma }_{bAB}^a(R^{ijk}+{\displaystyle \frac{15}{3!}}p^{ijkl_3}a_{l_3})e_k^b+{\displaystyle \frac{5ı}{2}}\mathrm{\Gamma }_{b_4AB}^ap^{ijk_4}E_{k_4}^{b_4}\right\},`$ (28)
$`J^{ab}`$ $`=`$ $`{\displaystyle \frac{1}{2}}D_i\pi ^{iab}\pi ^{i[a}e_i^{b]}{\displaystyle \frac{1}{4}}\mathrm{\Psi }_{iA}\mathrm{\Gamma }_{}^{abA}{}_{B}{}^{}\pi ^{Bi},`$ (29)
$`G_5^{i_5}`$ $`=`$ $`{\displaystyle \frac{1}{5!}}_ip^{ij_5},`$ (30)
$`G_2^{ij}`$ $`=`$ $`{\displaystyle \frac{1}{2}}[_kR^{ijk}{\displaystyle \frac{15}{3!}}(_ka_{l_3})p^{ijkl_3}].`$ (31)
It’s now straightforward to check the Poisson brackets of two supersymmetric constraint functional satisfies weakly the relation (1). More precisely, we find,
$`\{Q^A(x),Q^B(y)\}_+`$ $`=`$ $`\delta ^{10}(x,y)\{\mathrm{\Gamma }_a^{AB}G^a`$ (32)
$`\mathrm{\Gamma }_{ab}^{AB}\left[+E_{jk}^{ab}(G_2^{jk}+{\displaystyle \frac{15}{3!}}G_5^{jkl_3}a_{l_3})+{\displaystyle \frac{15}{3!}}p^{ijkl_3}(F_{il_3}^{\mathrm{}}E_{jk}^{ab}+a_{l_3}F_{ij}^ae_k^b)\right]`$
$`ı\mathrm{\Gamma }_{ab^4}^{AB}[E_{k_5}^{ab_4}G_5^{k_5}+p^{ikl_4}E_{l_4}^{b_4}F_{ik}^a]\}.`$
It is interesting to see that the constraints by which the curvatures vanish are needed to close the algebra. It is also interesting that in order to realize the central charge proportional to $`\mathrm{\Gamma }_{abcde}^{AB}`$ it is necessary to have the canonical momenta and the constraint associated with the six form field $`b`$.
Similarly, we find that
$`\{G^a(x),G^b(y)\}`$ $`=`$ $`\delta ^{10}(x,y)\{\mathrm{\Gamma }_{AB}^{ab}[2F_{ki}^A\mathrm{\Psi }_j^B[R^{ijk}+{\displaystyle \frac{15}{3!}}p^{ijkl_3}a_{l_3}]+\mathrm{\Psi }_i^A\mathrm{\Psi }_j^B({\displaystyle \frac{1}{4}}G_2^{ij}+G_5^{ijl_3}a_{l_3}+p^{ijl_4}F_{l_4}^{\mathrm{}})]`$ (33)
$`+\mathrm{\Gamma }_{defAB}^{ab}[2ıF_{ki}^A\mathrm{\Psi }_j^Bp^{ijkl_3}E_{l_3}^{def}+{\displaystyle \frac{ı}{4}}G_5^{ijl_3}E_{l_3}^{def}+ıp^{ijklm_2}F_{kl^d}E_{m_2}^{ef}]\}.`$
In both of these relations there are delicate cancellations involving two and four fermion terms. These involve careful application of the Fiertz identities for $`11`$ dimensions. It is also interesting that the central charges come into the commutator of the translation generators (33). This may be a clue as to how the whole structure may descend from some $`Osp(1|32)`$ invariant framework.
We also find that,
$$\{Q^A(x),G_2^{ij}(y)\}=\frac{5}{4}\delta ^{10}\mathrm{\Gamma }_{abB}^A\left(G_5^{lmnij}\mathrm{\Psi }_l^BE_{mn}^{ab}+p^{lmnkij}[F_{kl}^BE_{mn}^{ab}+2F_{kl}^a\mathrm{\Psi }_l^Be_n^b]\right),$$
(34)
$$\{G^a(x),G_2^{ij}(y)\}=\frac{5}{4}\delta ^{10}\mathrm{\Gamma }_{bAB}^a\left(G_5^{lmnij}\mathrm{\Psi }_l^A\mathrm{\Psi }_m^Be_n^b+p^{lmnkij}[2F_{kl}^A\mathrm{\Psi }_m^Be_n^b+\mathrm{\Psi }_l^A\mathrm{\Psi }_m^BF_{kn}^b]\right),$$
(35)
$$\{Q^A(x),G_5^{ijklm}(y)\}=\{G^a(x),G_5^{ijklm}(y)\}=0.$$
(36)
Thus, $`G_2^{ij}`$ and $`G_5^{ijklm}`$ form a supersymmetry multiplet. We find also
$$\{Q^A(x),G^a(y)\}=0.$$
(37)
The commutators of the $`J^{ab}`$ are defined by the transformation properties under the local lorentz group.
$$\{J^{ab},J_{cd}\}=\delta ^{10}(x,y)\delta _{[c}^{[a}J_{d]}^{b]},$$
(38)
$$\{J^{ab}(x),G^c(y)\}=\delta ^{10}(x,y)\delta _c^{[a}G^{b]},$$
(39)
$$\{J^{ab}(x),Q_A(y)\}=\delta ^{10}(x,y)\mathrm{\Gamma }_A^{abB}Q_B,$$
(40)
$$\{J^{ab}(x),G_2^{ij}(y)\}=0,$$
(41)
$$\{J^{ab}(x),G_5^{ijklm}\}=0.$$
(42)
Finally, we find that all the other communtators vanish. Thus the constraint algebra does in fact reproduce the super-Poincare algebra with central charges.
## 5 Conclusions
We have reported here the first step of a program to construct a non-perturbative formulation of $``$ theory by making a background independent quantization of $`11`$ dimensional supergravity. The next step is to construct the quantization of the $`11`$ dimensional topological quantum field theory, using the methods of loop quantum gravity. This can be done both canonically and through a path integral quantization using an extension of the methods of spin foam or evolving spin networks to $`p`$-form gauge fields. Some work in this direction already exists and this part of the program should go through directly. It is clear from the form of the theory that this will involve extended objects whose spacetime dimensions are $`2,3`$ and $`6`$. These will then give background independent objects corresponding to strings, membranes and five-branes.
In the topological quantum field theory these will have trivial dynamics and the states will be functionals only of homotopy classes of the ten dimensional spacial manifold (which is fixed in a canonical quantization). The problem will then be to reduce the gauge invariance of the topological quantum field theory so as to give rise to local degrees of freedom. There are two ways to accomplish this. The conservative, straightforward path will be to construct the canonical quantization of the eleven dimensional supergravity, by imposing the constraints in the above action classically. This will involve a lot of tedious calculation, to check the resulting algebra of constraints, but should be nonetheless straightforward. One will then impose the quantum constraints rather than the vanishing curvature conditions on the Hilbert space of states.
A second route to the theory will be to follow Barrett and Crane and impose the constraints directly in a path integral expression for the topological quantum field theory.
While these will involve a great deal of work, the key point is that at every stage one will be working with a background independent definition of the Hilbert space of the theory, whose degrees of freedom have the correct dimensionality and supersymmetry transformation properties to lead to strings, membranes and fivebranes in the classical limit. It is difficult to believe that something of value for the understanding of $``$ theory will not come out of such an investigation.
## ACKNOWLEDGEMENT
We are grateful to Clifford Johnson, Hermann Nicolai, Mike Reisenberger and Kelle Stelle for discussions during the course of this investigation as well as to the theoretical physics group at Imperial College for hospitality during the course of this work. This work was supported by the NSF through grant PHY95-14240 and a gift from the Jesse Phillips Foundation.
## Appendix A: Conventions and notations
Here we give the conventions and notations in this paper. We adopt the convention that if $`\omega `$ is a p-form, then:
$$\omega _p:=\omega _{\mu _1\mathrm{}\mu _p}dx^{\mu _1}dx^{\mu _2}\mathrm{}dx^{\mu _p},$$
(43)
or we could write it with abstract indices as,
$$\omega _{a_1a_2\mathrm{}a_p}=p!\omega _{\mu _1\mathrm{}\mu _p}(dx^{\mu _1})_{[a_1}(dx^{\mu _2})_{a_2}\mathrm{}(dx^{\mu _p})_{a_p]},$$
(44)
Where the antisymmetrization symbol is defined by,
$$[a_1\mathrm{}a_n]=\frac{1}{n!}\underset{p}{}()^{\delta _p}a_{p(1)}\mathrm{}a_{p(n)},$$
(45)
where $`_p`$ is the sum over permutations and $`\delta _p`$ is the parity of the permutation. Given any two forms such that one is p-form and the other one is q-form, we can define the wedge product of these two forms as,
$$\omega _p\omega _q^{^{}}=\frac{11!}{p!q!}\omega _{[a_1\mathrm{}a_p}\omega _{b_1\mathrm{}b_q]}^{^{}}.$$
(46)
It’s straightforward to show the wedge product has the following property,
$$\omega _p\omega _q^{^{}}=(1)^{pq}\omega _q^{^{}}\omega _p.$$
(47)
The exterior differential $`d`$ is a map from vector space of p-form to that of $`(p+1)`$-form,
$$d\omega _p=(d\omega )_{ab_1\mathrm{}b_p}=(p+1)_{[a}\omega _{b_1\mathrm{}b_p]}.$$
(48)
we can also show that
$$d(\omega _p\omega _q^{^{}})=(d\omega _p)\omega _q^{^{}}+(1)^p\omega _pd(\omega _q^{^{}}).$$
(49)
If $`\omega _{11}`$ is a 11-form on the manifold $``$, we define the integral of the form on the manifold as,
$$\omega _{11}=\frac{1}{11!}ϵ^{\mu _1\mathrm{}\mu _{11}}\omega _{\mu _1\mathrm{}\mu _{11}}d^{11}x,$$
(50)
where $`ϵ^{\mu _1\mathrm{}\mu _{11}}`$ is the volume element on $``$ such that
$$ϵ^{\mu _1\mathrm{}\mu _{11}}ϵ_{\mu _1\mathrm{}\mu _{11}}=11!,$$
(51)
and locally if the manifold splits into space $`\mathrm{\Sigma }_{10}`$ and time $`R`$ which is denoted by the coordinate $`0`$, then an induced volume on $`\mathrm{\Sigma }_{10}`$ is given by,
$$ϵ^{i_1\mathrm{}i_{10}}=ϵ^{0\mu _2\mathrm{}\mu _{11}}=\omega ^{\mu _10\mathrm{}\mu _{11}}=\mathrm{}=ϵ^{\mu _1\mu _2\mathrm{}\mu _{10}0}.$$
(52)
## Appendix B: Gamma matrix
Some important features of Gamma matrices in eleven dimensional space time are derived in this part. They are essential to show the closure of the constraint algebra in present paper. It’s well known that the $`\mathrm{\Gamma }`$-matrix plays an important role to describe the spinor fields in various dimensions. They form the Clifford algebra,
$$\mathrm{\Gamma }^a\mathrm{\Gamma }^c+\mathrm{\Gamma }^c\mathrm{\Gamma }^a=2\eta ^{ac},$$
(53)
we also introduce the notation
$$\mathrm{\Gamma }^a\mathrm{\Gamma }^c\mathrm{\Gamma }^c\mathrm{\Gamma }^a:=2\mathrm{\Gamma }^{ac},$$
(54)
Then in the case of eleven dimensional space time, we can derive the following identities
$$\mathrm{\Gamma }_a\mathrm{\Gamma }^a=11,$$
(55)
$$\mathrm{\Gamma }^a\mathrm{\Gamma }^c=\mathrm{\Gamma }^{ac}+\eta ^{ac}.$$
(56)
(55) and (56) are very useful when we try to simplify the expression or rearrange the $`\mathrm{\Gamma }`$-matrices into a new order. For instance,
$$\mathrm{\Gamma }_a\mathrm{\Gamma }^d\mathrm{\Gamma }^a=\mathrm{\Gamma }_a(2\eta ^{da}\mathrm{\Gamma }^a\mathrm{\Gamma }^d)=2\mathrm{\Gamma }^d\mathrm{\Gamma }_a\mathrm{\Gamma }^a\mathrm{\Gamma }^d=9\mathrm{\Gamma }^d,$$
(57)
$$\mathrm{\Gamma }_a\mathrm{\Gamma }^{d_1\mathrm{}d_n}\mathrm{\Gamma }^a=(1)^n(112n)\mathrm{\Gamma }^{d_1\mathrm{}d_n}.$$
(58)
In this paper we also often use the following formula which are given in ,
$$\mathrm{\Gamma }^{a_1\mathrm{}a_nb}=\mathrm{\Gamma }^{a_1\mathrm{}a_n}\mathrm{\Gamma }^bn\mathrm{\Gamma }^{[a_1\mathrm{}a_{n1}}\eta ^{a_n]b},$$
(59)
$$\mathrm{\Gamma }^{ba_1\mathrm{}a_n}=\mathrm{\Gamma }^b\mathrm{\Gamma }^{a_1\mathrm{}a_n}n\eta ^{b[a_1}\mathrm{\Gamma }^{a_2\mathrm{}a_n]}.$$
(60)
Next we give two important identities of Gamma matrices which are essential to show the closure of constraint algebra. Both of them involve the exchanging of spinor indices of different $`\mathrm{\Gamma }`$-matrices, therefore we write down the elements of these matrices labeled by the spinor indices explicitly. These identities are<sup>9</sup><sup>9</sup>9We ignore an important matrix in all the paper, namely, the charge conjugation matrix for a neat version of Gamma matrix. But we’d better keep it in mind and realize it appeared where it should be. In eleven dimensions it is also important to know that only $`\{\mathrm{\Gamma }^a,\mathrm{\Gamma }^{ab},\mathrm{\Gamma }^{a_1\mathrm{}a_5}\}`$ and their dual matrices are symmetric (under the action of charge conjugation matrix) while the others are antisymmetric. Since in the Poisson bracket of supersymmetry constraints, the charge conjugation matrix is involved and any term which is anti-symmetric will vanish. Therefore in the following equations we only write down the symmetric term explicitly.
$$\mathrm{\Gamma }_{a}^{}{}_{E}{}^{(A}\mathrm{\Gamma }_{}^{ab}{}_{F}{}^{B)}=\frac{1}{4}(\mathrm{\Gamma }_a^{AB}\mathrm{\Gamma }_{EF}^{ab}+\mathrm{\Gamma }^{abAB}\mathrm{\Gamma }_{aEF}),$$
(61)
$$\mathrm{\Gamma }_{}^{a}{}_{E}{}^{(A}\mathrm{\Gamma }_{ab_1\mathrm{}b_4}^{}{}_{F}{}^{B)}3\mathrm{\Gamma }_{[b_1b_2}^{}{}_{E}{}^{(A}\mathrm{\Gamma }_{b_3b_4]}^{}{}_{F}{}^{B)}=\frac{1}{4}(\mathrm{\Gamma }^{aAB}\mathrm{\Gamma }_{EF}^{ab_1\mathrm{}b_4}+\mathrm{\Gamma }_{ab_1\mathrm{}b_4}^{AB}\mathrm{\Gamma }_{EF}^a)+\frac{3}{2}\mathrm{\Gamma }_{[b_1b_2}^{}{}_{}{}^{AB}\mathrm{\Gamma }_{b_3b_4]EF},$$
(62)
It’s not difficult to see, they go back to the ordinary Fiertz identities respectively when four spinor fields are involved,
$$\mathrm{\Psi }_A\mathrm{\Gamma }_B^{aA}\mathrm{\Psi }^B\mathrm{\Psi }_E\mathrm{\Gamma }_{ab}^{}{}_{F}{}^{E}\mathrm{\Psi }^F=0,$$
(63)
$$\mathrm{\Psi }_A\mathrm{\Gamma }_B^{aA}\mathrm{\Psi }^B\mathrm{\Psi }_E\mathrm{\Gamma }_{ab_1..b_4}^{}{}_{F}{}^{E}\mathrm{\Psi }^F=3\mathrm{\Psi }_A\mathrm{\Gamma }_{[b_1b_2B}^A\mathrm{\Psi }^B\mathrm{\Psi }_E\mathrm{\Gamma }_{b_3b_4]}^{}{}_{F}{}^{E}\mathrm{\Psi }^F.$$
(64)
To prove identities $`(\text{61})`$ and $`(\text{62})`$, we need apply the Fiertz decomposition formula in eleven dimensional space time.
$$\mathrm{\Gamma }_E^{a(A}\mathrm{\Gamma }_{ab}^{}{}_{F}{}^{B)}=\frac{1}{32}\left[\mathrm{\Gamma }_d^{AB}(\mathrm{\Gamma }_a\mathrm{\Gamma }^d\mathrm{\Gamma }^{ab})_{EF}\frac{1}{2}\mathrm{\Gamma }_{de}^{AB}(\mathrm{\Gamma }_a\mathrm{\Gamma }^{de}\mathrm{\Gamma }^{ab})_{EF}+\frac{1}{5!}\mathrm{\Gamma }_{d_1\mathrm{}d_5}^{AB}(\mathrm{\Gamma }_a\mathrm{\Gamma }^{d_1\mathrm{}d_5}\mathrm{\Gamma }^{ab})_{EF}\right].$$
(65)
Then our task is just to simplify the terms involving the multiplication of several Gamma matrices in $`(\text{65})`$. Exploiting the formula given above, it’s straightforward to derive the following results,
$$\mathrm{\Gamma }_a\mathrm{\Gamma }^d\mathrm{\Gamma }^{ab}=8\mathrm{\Gamma }^{db}+(antisymmetricterms\mathrm{}),$$
(66)
$$\mathrm{\Gamma }_a\mathrm{\Gamma }^{de}\mathrm{\Gamma }^{ab}=16\mathrm{\Gamma }^{[d}\eta ^{e]b}+(antisymmetricterms\mathrm{}),$$
(67)
$$\mathrm{\Gamma }_a\mathrm{\Gamma }^{d_1\mathrm{}d_5}\mathrm{\Gamma }^{ab}=0+(antisymmetricterms\mathrm{}).$$
(68)
Substituting $`(\text{66})`$-$`(\text{68})`$ into $`(\text{65})`$, we easily arrive at the identity $`(\text{61})`$.
As far as the second identity $`(\text{62})`$ is concerned, we just need do more complicated but straightforward calculations as in the case of first identity.
$$\mathrm{\Gamma }_E^{a(A}\mathrm{\Gamma }_{ab_1\mathrm{}b_4}^{}{}_{F}{}^{B)}=\frac{1}{32}\left[\mathrm{\Gamma }_d^{AB}(\mathrm{\Gamma }_a\mathrm{\Gamma }^d\mathrm{\Gamma }^{ab_1\mathrm{}b_4})_{EF}\frac{1}{2}\mathrm{\Gamma }_{de}^{AB}(\mathrm{\Gamma }_a\mathrm{\Gamma }^{de}\mathrm{\Gamma }^{ab_1\mathrm{}b_4})_{EF}+\frac{1}{5!}\mathrm{\Gamma }_{d_1\mathrm{}d_5}^{AB}(\mathrm{\Gamma }_a\mathrm{\Gamma }^{d_1\mathrm{}d_5}\mathrm{\Gamma }^{ab_1\mathrm{}b_4})_{EF}\right],$$
(69)
and
$$\mathrm{\Gamma }_{[b_1b_2}^{}{}_{E}{}^{(A}\mathrm{\Gamma }_{b_3b_4]}^{}{}_{F}{}^{B)}=\frac{1}{32}\left[\mathrm{\Gamma }_d^{AB}(\mathrm{\Gamma }_{[b_1b_2}\mathrm{\Gamma }^d\mathrm{\Gamma }_{b_3b_4]})_{EF}\frac{1}{2}\mathrm{\Gamma }_{de}^{AB}(\mathrm{\Gamma }_{[b_1b_2}\mathrm{\Gamma }^{de}\mathrm{\Gamma }_{b_3b_4]})_{EF}+\frac{1}{5!}\mathrm{\Gamma }_{d_1\mathrm{}d_5}^{AB}(\mathrm{\Gamma }_{[b_1b_2}\mathrm{\Gamma }^{d_1\mathrm{}d_5}\mathrm{\Gamma }_{b_3b_4]})_{EF}\right].$$
(70)
In $`(\text{69})`$, the terms in brackets can be simplified respectively as,
$$\mathrm{\Gamma }_a\mathrm{\Gamma }^d\mathrm{\Gamma }^{ab_1\mathrm{}b_4}=5\mathrm{\Gamma }^{db_1\mathrm{}b_4}+antisym.terms,$$
(71)
$$\mathrm{\Gamma }_a\mathrm{\Gamma }^{de}\mathrm{\Gamma }^{ab_1\mathrm{}b_4}=3\mathrm{\Gamma }^{deb_1\mathrm{}b_4}+84\eta ^{e[b_1}\eta ^{|d|b_2}\mathrm{\Gamma }^{b_3b_4]}+antisym.terms,$$
(72)
$`\mathrm{\Gamma }_a\mathrm{\Gamma }^{d_1\mathrm{}d_5}\mathrm{\Gamma }^{ab_1\mathrm{}b_4}`$ $`=`$ $`3\mathrm{\Gamma }^{d_1\mathrm{}d_5}\mathrm{\Gamma }^{b_1\mathrm{}b_4}40\mathrm{\Gamma }^{[d_1\mathrm{}d_4}\eta ^{d_5][b_1}\mathrm{\Gamma }^{b_2\mathrm{}b_4]}`$ (73)
$`=`$ $`3\mathrm{\Gamma }_{b_1\mathrm{}b_4}^{d_1\mathrm{}d_5}5!\delta _{[b_2}^{[d_1}\delta _{b_1}^{d_2}\mathrm{\Gamma }_{b_3b_4]}^{d_3d_4d_5]}+55!\mathrm{\Gamma }^{[d_1}\delta _{[b_1}^{d_2}\delta _{b_2}^{d_3}\delta _{b_3}^{d_4}\delta _{b_4]}^{d_5]}`$
$`+antisym.terms,`$
and in (70), the terms in brackets can be simplified as
$$\mathrm{\Gamma }_{[b_1b_2}\mathrm{\Gamma }^d\mathrm{\Gamma }_{b_3b_4]}=\mathrm{\Gamma }_{b_1\mathrm{}b_4}^d+antisym.terms,$$
(74)
$$\mathrm{\Gamma }_{[b_1b_2}\mathrm{\Gamma }^{de}\mathrm{\Gamma }_{b_3b_4]}=\mathrm{\Gamma }_{b_1\mathrm{}b_4}^{de}4\delta _{[b_2}^{[d}\delta _{b_1}^{e]}\mathrm{\Gamma }_{b_3b_4]}+antisym.terms,$$
(75)
$$\mathrm{\Gamma }_{[b_1b_2}\mathrm{\Gamma }^{d_1\mathrm{}d_5}\mathrm{\Gamma }_{b_3b_4]}=\mathrm{\Gamma }_{b_1\mathrm{}b_4}^{d_1\mathrm{}d_5}40\delta _{[b_2}^{[d_1}\delta _{b_1}^{d_2}\mathrm{\Gamma }_{b_3b_4]}^{d_3d_4d_5]}+5!\delta _{[b_2}^{[d_1}\delta d_{2}^{}{}_{b_1}{}^{}\mathrm{\Gamma }^{d_3}\delta _{b_4}^{d_4}\delta _{b_3]}^{d_5]}.$$
(76)
Substituting all the terms into (69) and (70) respectively, we will find the identity (62) holds indeed.
## Appendix C: The proof of the closure of constraint algebra
In this section we only show the closure of two Poisson brackets. One involves the Gaussian constraint, and the other is the supersymmetric constraint. The other Poisson brackets are closed trivially. First we consider the Poisson bracket of Gaussian constraint. To make the calculation clear, we divide the constraint into two parts,
$$G^a=:G_1^a+G_2^a,$$
(77)
Where
$$G_1^a=D_k\pi ^{ka},$$
(78)
and
$$G_2^a=\mathrm{\Psi }_i^A\mathrm{\Psi }_j^B\left\{\mathrm{\Gamma }_{bAB}^a(R^{ijk}+\frac{15}{3!}p^{ijkl_3}a_{l_3})e_k^b+\frac{5ı}{2}\mathrm{\Gamma }_{b_4AB}^ap^{ijk_4}E_{k_4}^{b_4}\right\}.$$
(79)
It’s straightforward to compute the Poisson brackets of them,
$$\{G_1^a,G_2^b\}=0,$$
(80)
$$\{G_2^a,G_2^b\}=\delta ^{10}30\mathrm{\Psi }_{iA}\mathrm{\Gamma }_B^{acA}\mathrm{\Psi }_j^B\mathrm{\Psi }_{mE}\mathrm{\Gamma }_D^{bdE}\mathrm{\Psi }_n^De_{kc}e_{pd}p^{mnpijk},$$
(81)
$`\{G_1^a,G_2^b\}`$ $`+`$ $`\{G_2^a,G_1^b\}=\delta ^{10}\{[15\mathrm{\Psi }_{iA}\mathrm{\Gamma }_B^{abA}\mathrm{\Psi }_j^B\mathrm{\Psi }_{mE}\mathrm{\Gamma }_D^{cdE}\mathrm{\Psi }_n^De_{kc}e_{pd}p^{mnpijk}`$ (84)
$`15\mathrm{\Psi }_{iA}\mathrm{\Gamma }_B^{abcdeA}\mathrm{\Psi }_j^B\mathrm{\Psi }_{kC}\mathrm{\Gamma }_{cD}^C\mathrm{\Psi }_p^DE_{demn}p^{ijkmnp}]`$
$`+\mathrm{\Gamma }_{AB}^{ab}\left[2F_{ki}^A\mathrm{\Psi }_j^B(R^{ijk}+{\displaystyle \frac{15}{3!}}p^{ijkl_3}a_{l_3})+\mathrm{\Psi }_i^A\mathrm{\Psi }_j^B({\displaystyle \frac{1}{4}}G_2^{ij}+G_5^{ijl_3}a_{l_3}+p^{ijl_4}F_{l_4}^{\mathrm{}})\right]`$
$`+\mathrm{\Gamma }_{defAB}^{ab}[2ıF_{ki}^A\mathrm{\Psi }_j^Bp^{ijkl_3}E_{l_3}^{def}+{\displaystyle \frac{ı}{4}}G_5^{ijl_3}E_{l_3}^{def}+ıp^{ijklm_2}F_{kl^d}E_{m_2}^{ef}]\}.`$
Now add $`(\text{81})`$ and $`(\text{84})`$ together and notice that
$$\mathrm{\Gamma }^{[ab}\mathrm{\Gamma }^{cd]}=\frac{1}{3}(\mathrm{\Gamma }^{ab}\mathrm{\Gamma }^{cd}+\mathrm{\Gamma }^{ac}\mathrm{\Gamma }^{db}+\mathrm{\Gamma }^{ad}\mathrm{\Gamma }^{bc}).$$
(85)
we find the sum of three terms $`(\text{81})`$, $`(\text{84})`$ and $`(\text{84})`$ vanishes by employing the standard Fiertz identity $`(\text{64})`$. Making a collection of $`(\text{80})`$-$`(\text{84})`$, we show the closure of Poisson bracket of Gauss constraint which corresponds to (33) in the paper.
The Poisson bracket $`(\text{32})`$ can be derived in a similar way except that we need apply the identities (61) and (62) to cancel those extra terms.
The supersymmetric constraint is
$$Q^A=Q_1^A+Q_2^A,$$
(86)
where
$$Q_1^A=D_i\pi ^{iA}ı\pi ^{ia}\mathrm{\Gamma }_{aB}^A\mathrm{\Psi }_i^B,$$
(87)
and
$$Q_2^A=\mathrm{\Psi }_i^B\left\{\mathrm{\Gamma }_{abB}^A(R^{ijk}+\frac{15}{3!}p^{ijkl_3}a_{l_3})E_{jk}^{ab}ı\mathrm{\Gamma }_{a_5B}^Ap^{ik_5}E_{k_5}^{b_5}\right\}.$$
(88)
We find the Poisson brackets of them are
$$\{Q_1^A,Q_1^B\}=i\delta ^{10}\mathrm{\Gamma }_a^{AB}D_i\pi ^{ia},$$
(89)
$$\{Q_2^A,Q_2^B\}=30\delta ^{10}\mathrm{\Psi }_i^C\mathrm{\Psi }_j^D\mathrm{\Gamma }_{[abC}^{(A}\mathrm{\Gamma }_{cd]D}^{B)}E_{kmnp}^{abcd}p^{ijkmnp},$$
(90)
$`\{Q_1^A,Q_2^B\}`$ $`+`$ $`\{Q_2^A+Q_1^B\}=\delta ^{10}\{D_i[\mathrm{\Gamma }_{ab}^{AB}e_{jk}^{ab}(R^{ijk}+{\displaystyle \frac{15}{3!}}p^{ijkl_3}a_{l_3})+i\mathrm{\Gamma }_{a_5}^{AB}E_{j_5}^{a_5}p^{ij^5}]`$ (91)
$`+4i\mathrm{\Psi }_i^C\mathrm{\Psi }_j^D(\mathrm{\Gamma }_{aC}^{(A}\mathrm{\Gamma }_D^{B)ab}e_{bk}(R^{ijk}+{\displaystyle \frac{15}{3!}}p^{ijkl_3}a_{l_3})+{\displaystyle \frac{5i}{2}}\mathrm{\Gamma }_{aC}^{(A}\mathrm{\Gamma }^{B)ab_4}E^{b_4k_4}p^{ijk_4})\}.`$
In (91), it’s a little tedious to deal with the terms with covariant derivative. To express them as the sum of constraints, we use the fact that
$$D_{[\mu }e_{\nu ]}^a=\frac{1}{2}F_{\mu \nu }^a+\frac{i}{2}\mathrm{\Psi }_{A[\mu }\mathrm{\Gamma }_B^{aA}\mathrm{\Psi }_{\nu ]}^B,$$
(92)
and
$$_{[\mu }a_{\nu \rho \sigma ]}=\frac{1}{4}F^{\mathrm{}}+\frac{3!}{2}\mathrm{\Psi }_{A[\mu }\mathrm{\Gamma }_B^{abA}\mathrm{\Psi }_\nu ^Be_{|a|\rho }e_{|b|\sigma ]},$$
(93)
and then expand those terms as follows,
$`D_i[\mathrm{\Gamma }_{ab}^{AB}e^ab_{jk}(R^{ijk}+{\displaystyle \frac{15}{3!}}a_{mnp}p^{ijkmnp})]=(\mathrm{})F^a+(\mathrm{})G^{ij}+(\mathrm{})G^{i_5}+(\mathrm{})F_{imnp}^{\mathrm{}}`$
$`+i\mathrm{\Gamma }_{ab}^{AB}\mathrm{\Psi }_{ic}\mathrm{\Gamma }_D^{aC}\mathrm{\Psi }_j^De_k^b(R^{ijk}+{\displaystyle \frac{15}{3!}}a_{mnp}p^{ijkmnp})+15\mathrm{\Gamma }_{ab}^{AB}e_{jk}^{ab}\mathrm{\Psi }_{iC}\mathrm{\Gamma }_{cdD}^C\mathrm{\Psi }_m^De_{np}^{cd}p^{ijkmnp},`$ (94)
$$D_i(\mathrm{\Gamma }_{a_5}^{AB}E_{j_5}^{a_5}p^{ij^5})=(\mathrm{})G^{j_5}+(\mathrm{})F^a+\frac{5i}{2}\mathrm{\Gamma }_{ab_4}^{AB}\mathrm{\Psi }_{iC}\mathrm{\Gamma }_D^{aC}\mathrm{\Psi }_j^DE_{k_4}^{b_4}p^{ijk_4},$$
(95)
where we ignore the explicit expressions of terms involving curvatures and constraints, but the final results are given in (32). Making use of identities $`(\text{61})`$, we add the terms containing $`(R^{ijk}+\frac{15}{3!}p^{ijkl_3}a_{l_3})`$ together in the brackets of $`(\text{91})`$ and have
$$i\mathrm{\Psi }_i^C\mathrm{\Psi }_j^De_{bk}(R^{ijk}+\frac{15}{3!}p^{ijkl_3}a_{l_3})(4\mathrm{\Gamma }_{aC}^{(A}\mathrm{\Gamma }_D^{B)ab}+\mathrm{\Gamma }_{ab}^{AB}\mathrm{\Gamma }_{CD}^a)=i\mathrm{\Psi }_i^C\mathrm{\Psi }_j^De_{bk}(R^{ijk}+\frac{15}{3!}p^{ijkl_3}a_{l_3})\mathrm{\Gamma }_a^{AB}\mathrm{\Gamma }_{CD}^{ab},$$
(96)
and using $`(\text{62})`$, we pick out all the terms containing $`p^{ijkmnp}`$ in the algebra $`(\text{90})`$ and $`(\text{91})`$, and find
$`\mathrm{\Psi }_i^C\mathrm{\Psi }_j^DE_{kmnp}^{abcd}p^{ijkmnp}(30\mathrm{\Gamma }_{C[ab}^{(A}\mathrm{\Gamma }_{cd]D}^{B)}10\mathrm{\Gamma }_{Ce}^{(A}\mathrm{\Gamma }_D^{B)eabcd}{\displaystyle \frac{5}{2}}\mathrm{\Gamma }_{eabcd}^{AB}\mathrm{\Gamma }_{CD}^e+15\mathrm{\Gamma }_{[ab}^{AB}\mathrm{\Gamma }_{cd]CD})`$
$`={\displaystyle \frac{5}{2}}\mathrm{\Psi }_i^C\mathrm{\Psi }_j^DE_{k_4}^{b_4}p^{ijkmnp}\mathrm{\Gamma }_a^{AB}\mathrm{\Gamma }_{b_4CD}^a.`$ (97)
Next combining all the terms in (89), (96) and (97) together, we find it’s nothing but the Gauss constraint $`G^a`$! After collecting the other terms remaining which contain curvatures and constraints in the algebra we finally arrive at $`(\text{32})`$, which is what we need to show.
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# On the ‘Irreducible’ Freedman-Townsend Vertex
## Abstract
An irreducible cohomological derivation of the Freedman-Townsend vertex in four dimensions is given.
PACS number: 11.10.Ef
The problem of consistent interactions that can be introduced among gauge fields in such a way to preserve the number of gauge symmetries has been reformulated as a deformation problem of the master equation in the framework of the antifield BRST formalism . That deformation setting was applied to Chern-Simons models , Yang-Mills theories and two-form gauge fields . The deformation procedure for two-form gauge fields employs a reducible BRST background.
The purpose of this letter is to reanalyze the problem of constructing consistent interactions among two-form gauge fields (in four dimensions), but following an irreducible BRST line in spite of the reducibility present within the initial model. Our method contains two basic steps. First, starting with abelian two-form gauge fields in first-order version (obtained by means of adding some auxiliary fields), we construct an irreducible BRST symmetry associated with the reducible one and show that we can substitute the reducible symmetry by the irreducible one. Second, we consistently deform the solution to the master equation associated with the irreducible BRST symmetry. In this manner we obtain precisely the well-known Freedman-Townsend vertex , the deformed gauge symmetries and also the deformed solution to the master equation. However, our deformed solution to the master equation differs from that obtained in the literature by the fact that on the one hand our method introduces some additional fields (necessary at the construction of the irreducible BRST symmetry) which do not contribute to new type of couplings, and, on the other hand, the deformed BRST transformations resulting from our formalism do not involve the antifields due to the absence of terms quadratic in the antifields in the deformed solution to the master equation, so the gauge-fixed BRST symmetry is off-shell nilpotent. To our knowledge, such an irreducible procedure for the Freedman-Townsend model has not previously been published, our method establishing thus a new result.
The starting point is the Lagrangian action for abelian two-form gauge fields in first-order form (also known as the abelian Freedman-Townsend model)
$$S_0^L[A_\mu ^a,B_a^{\mu \nu }]=\frac{1}{2}d^4x\left(F_{\mu \nu }^aB_a^{\mu \nu }+g_{ab}A_\mu ^aA^{b\mu }\right),$$
(1)
where $`B_a^{\mu \nu }`$ is an antisymmetric tensor field, the field strength of $`A_\mu ^a`$ is defined by $`F_{\mu \nu }^a=_{[\mu }A_{\nu ]}^a`$, with $`\left[\mu \nu \right]`$ expressing antisymmetry with the indices between brackets, and $`g_{ab}`$ is an invertible, symmetric and constant matrix. It is simply to see that if we eliminate the auxiliary fields $`A_\mu ^a`$ on their equations of motion, we recover the action of free abelian two-form gauge fields. Action (1) is invariant under the gauge transformations $`\delta _ϵB_a^{\mu \nu }=\epsilon ^{\mu \nu \lambda \rho }_\lambda ϵ_{\rho a}`$, $`\delta _ϵA_\mu ^a=0`$, where $`\epsilon ^{\mu \nu \lambda \rho }`$ is the antisymmetric symbol in four dimensions. The above gauge transformations are off-shell first-stage reducible as if we take $`ϵ_{\rho a}=_\rho ϵ_a`$, then $`\delta _ϵB_a^{\mu \nu }=0`$.
The reducible Lagrangian BRST symmetry corresponding to the model described by action (1), $`s_R=\delta _R+\gamma _R`$, contains two main pieces, namely, the Koszul-Tate differential $`\delta _R`$ and a model of longitudinal derivative along the gauge orbits $`\gamma _R`$. In the case of our model, the generators of the Koszul-Tate complex are the fermionic antighost number one antifields $`B_{\mu \nu }^a`$ and $`A_a^\mu `$, the bosonic antighost number two antifields $`\eta ^{a\mu }`$ and the fermionic antighost number three antifields $`C^a`$. The definitions of $`\delta _R`$ read as
$$\delta _RB_a^{\mu \nu }=0,\delta _RA_\mu ^a=0,$$
(2)
$$\delta _RB_{\mu \nu }^a=\frac{1}{2}F_{\mu \nu }^a,\delta _RA_a^\mu =\left(g_{ab}A^{b\mu }+_\nu B_a^{\nu \mu }\right),$$
(3)
$$\delta _R\eta ^{a\mu }=\epsilon ^{\mu \nu \lambda \rho }_\nu B_{\lambda \rho }^a,$$
(4)
$$\delta _RC^a=_\mu \eta ^{a\mu }.$$
(5)
The introduction of the antifields $`C^a`$ is implied by the necessity to ‘kill’ the non trivial antighost number two co-cycles $`\mu ^a=_\mu \eta ^{a\mu }`$ in the homology of $`\delta _R`$. The longitudinal complex contains the pure ghost number one fermionic ghosts $`\eta _{a\mu }`$ and the pure ghost number two bosonic ghosts for ghosts $`C_a`$. The definitions of $`\gamma _R`$ read as
$$\gamma _RA_\mu ^a=0,\gamma _RB_a^{\mu \nu }=\epsilon ^{\mu \nu \lambda \rho }_\lambda \eta _{\rho a},\gamma _R\eta _{a\mu }=_\mu C_a,\gamma _RC_a=0.$$
(6)
Extending $`\delta _R`$ on the ghosts through $`\delta _R\eta _{a\mu }=0`$, $`\delta _RC_a=0`$, and $`\gamma _R`$ on the antifields by $`\gamma _RA_a^\mu =0`$, $`\gamma _RB_{\mu \nu }^a=0`$, $`\gamma _R\eta ^{a\mu }=0`$, $`\gamma _RC^a=0`$, we find that $`s_R^2=0`$, $`H^0\left(s_R\right)=\left\{\mathrm{physical}\mathrm{observables}\right\}`$, where $`H^0\left(s_R\right)`$ represents the zeroth order cohomological group of $`s_R`$.
The main idea underlying our construction is to redefine the antifields $`\eta ^{a\mu }`$ in such a way that the new co-cycles of the type $`\mu ^a`$ identically vanish. If this is done, then the antifields $`C^a`$ are useless as there are no longer any non trivial co-cycles at antighost number two. In this way, we infer an irreducible Koszul-Tate complex, which further leads to a longitudinal complex that contains no more the ghosts for ghosts $`C_a`$. Accordingly our idea, we redefine the antifields $`\eta ^{a\mu }`$ like
$$\eta ^{a\mu }\eta ^{a\mu }=M_{b\nu }^{a\mu }\eta ^{b\nu },$$
(7)
where $`M_{b\nu }^{a\mu }`$ are taken to satisfy the conditions
$$_\mu M_{b\nu }^{a\mu }=0,$$
(8)
$$M_{b\nu }^{a\mu }\epsilon ^{\nu \sigma \lambda \rho }_\sigma B_{\lambda \rho }^b=\epsilon ^{\mu \sigma \lambda \rho }_\sigma B_{\lambda \rho }^a.$$
(9)
With the help of (4), (7) and (9) we find that
$$\delta \eta ^{a\mu }=\epsilon ^{\mu \sigma \lambda \rho }_\sigma B_{\lambda \rho }^a.$$
(10)
The last equations do not further imply non trivial co-cycles because the new co-cycles of the type $`\mu ^a`$ identically vanish via (8), hence we passed to an irreducible situation. In (10) we employed the notation $`\delta `$ instead of $`\delta _R`$ in order to emphasize that the Koszul-Tate complex becomes irreducible. The solution to (78) is expressed by
$$M_{b\nu }^{a\mu }=\delta _b^a\left(\delta _\nu ^\mu \frac{^\mu _\nu }{\mathrm{}}\right),$$
(11)
where $`\mathrm{}=_\lambda ^\lambda `$. Substituting (11) in (10), we get
$$\delta \left(\eta ^{a\mu }\frac{^\mu _\nu }{\mathrm{}}\eta ^{a\nu }\right)=\epsilon ^{\mu \sigma \lambda \rho }_\sigma B_{\lambda \rho }^a.$$
(12)
At this stage we introduce some scalar fields $`\phi _a`$ whose antifields $`\phi ^a`$ are demanded to be the non vanishing solutions to the equations
$$\mathrm{}\phi ^a=\delta \left(_\mu \eta ^{a\mu }\right).$$
(13)
The non vanishing solutions $`\phi ^a`$ enforce the irreducibility as (13) possess non vanishing solutions if and only if $`\delta \left(_\mu \eta ^{a\mu }\right)0`$, therefore if and only if $`\mu ^a`$ are no longer co-cycles. Using (1213) we find that
$$\delta \eta ^{a\mu }=\epsilon ^{\mu \sigma \lambda \rho }_\sigma B_{\lambda \rho }^a^\mu \phi ^a.$$
(14)
In order to preserve the nilpotency of $`\delta `$ we set
$$\delta \phi ^a=0.$$
(15)
If we maintain the actions of $`\delta `$ like in the reducible case
$$\delta B_a^{\mu \nu }=0,\delta A_\mu ^a=0,$$
(16)
$$\delta B_{\mu \nu }^a=\frac{1}{2}F_{\mu \nu }^a,\delta A_a^\mu =\left(g_{ab}A^{b\mu }+_\nu B_a^{\nu \mu }\right),$$
(17)
and define
$$\delta \phi _a=0,$$
(18)
then the formulas (1418) describe an irreducible Koszul-Tate complex. We remark that the irreducibility was gained by introducing the supplementary fields $`\phi _a`$ and their antifields $`\phi ^a`$ in the theory. From (1418) we can derive the Lagrangian action and the gauge transformations of the irreducible theory. If we denote by $`\stackrel{~}{S}_0^L[A_\mu ^a,B_a^{\mu \nu },\phi _a]`$ the Lagrangian action of the irreducible model, then by means of the general relations $`\delta \phi ^a=\delta \stackrel{~}{S}_0^L/\delta \phi _a`$ and (15) we obtain that
$$\stackrel{~}{S}_0^L[A_\mu ^a,B_a^{\mu \nu },\phi _a]=S_0^L[A_\mu ^a,B_a^{\mu \nu }],$$
(19)
such that the dependence on $`\phi _a`$ is trivial. On the other hand, with the help of the original gauge transformations and (14), it results that the gauge transformations of the irreducible system are expressed by
$$\delta _ϵB_a^{\mu \nu }=\epsilon ^{\mu \nu \lambda \rho }_\lambda ϵ_{\rho a},\delta _ϵA_\mu ^a=0,\delta _ϵ\phi _a=^\mu ϵ_{a\mu }.$$
(20)
In this manner, we derived an irreducible theory based on action (19) and the irreducible gauge transformations (20) associated with the abelian Freedman-Townsend model. From (19) we notice that the newly added fields $`\phi _a`$ are not involved with the Lagrangian action of the irreducible theory, hence they are purely gauge. As a consequence, the physical observables (gauge invariant functions) of the irreducible model do not depend on the $`\phi _a`$’s and, in addition, are invariant under the gauge transformations $`\delta _ϵB_a^{\mu \nu }=\epsilon ^{\mu \nu \lambda \rho }_\lambda ϵ_{\rho a},\delta _ϵA_\mu ^a=0`$, so they coincide with the physical observables of the original redundant theory. The construction of the irreducible longitudinal differential along the gauge orbits, $`\gamma `$ is realized via the definitions
$$\gamma B_a^{\mu \nu }=\epsilon ^{\mu \nu \lambda \rho }_\lambda \eta _{\rho a},\gamma A_a^\mu =0,\gamma \phi _a=^\mu \eta _{a\mu },\gamma \eta _{a\mu }=0,$$
(21)
such that $`\gamma `$ is nilpotent, $`\gamma ^2=0`$, without introducing the ghosts for ghosts. If we extend $`\delta `$ to the ghosts through $`\delta \eta _{a\mu }=0`$ and $`\gamma `$ to the antifields by $`\gamma B_{\mu \nu }^a=0`$, $`\gamma A_a^\mu =0`$, $`\gamma \phi ^a=0`$, $`\gamma \eta ^{a\mu }=0`$, then the homological perturbation theory guarantees the existence of the irreducible BRST symmetry $`s_I=\delta +\gamma `$ that is nilpotent, $`s_I^2=0`$, and satisfies the property $`H^0\left(s_I\right)=\left\{\mathrm{physical}\mathrm{observables}\right\}`$, where ‘physical observables’ are referring to the irreducible system. As we previously mentioned, the physical observables corresponding to the reducible and irreducible formulations coincide, which leads to $`H^0\left(s_R\right)=H^0\left(s_I\right)`$, and moreover, the two Lagrangian BRST symmetries are nilpotent $`s_R^2=0=s_I^2`$. By virtue of the last two relations we conclude that the two symmetries are equivalent from the BRST point of view, i.e., from the point of view of the basic equations underlying the antifield-BRST formalism. In consequence, we can replace the reducible Lagrangian BRST symmetry with the irreducible one in the case of the model under study.
With the above conclusion at hand, we pass to the deformation procedure of the irreducible version in the context of the antifield formalism . A consistent deformation of the free action $`S_0^L[A_\mu ^a,B_a^{\mu \nu }]`$ and of its gauge invariances defines a deformation of the corresponding solution to the master equation that preserves both the master equation and the field/antifield spectra. So, if $`S_0^L[A_\mu ^a,B_a^{\mu \nu }]+gd^4x\alpha _0+O\left(g^2\right)`$ stands for a consistent deformation of the free action, with deformed gauge transformations $`\overline{\delta }_ϵB_a^{\mu \nu }=\epsilon ^{\mu \nu \lambda \rho }_\lambda ϵ_{\rho a}+g\beta _a^{\mu \nu }+O\left(g^2\right)`$, $`\overline{\delta }_ϵ\phi _a=^\mu ϵ_{a\mu }+g\beta _a+O\left(g^2\right)`$, then the deformed solution to the master equation
$$\overline{S}=S+gd^4x\alpha +O\left(g^2\right),$$
(22)
satisfies $`(\overline{S},\overline{S})=0`$, where
$$S=S_0^L[A_\mu ^a,B_a^{\mu \nu }]+d^4x\left(\epsilon ^{\mu \nu \lambda \rho }B_{\mu \nu }^a_\lambda \eta _{\rho a}+\phi ^a^\mu \eta _{a\mu }\right),$$
(23)
and $`\alpha =\alpha _0+B_{\mu \nu }^a\overline{\beta }_a^{\mu \nu }+\phi ^a\overline{\beta }_a+\mathrm{`}\mathrm{more}`$’. Here, ‘$`\mathrm{more}`$’ stands for terms ’of antighost number greater than one. The master equation $`(\overline{S},\overline{S})=0`$ holds to order $`g`$ if and only if
$$s_I\alpha =_\mu j^\mu ,$$
(24)
for some local $`j^\mu `$. This means that the non trivial first-order consistent interactions belong to $`H^0\left(s_I|d\right)`$, where $`d`$ is the exterior space-time derivative. In the case where $`\alpha `$ is a coboundary modulo $`d`$ ($`\alpha =s_I\rho +_\mu b^\mu `$), then the deformation is trivial (it can be eliminated by a redefinition of the fields). In order to investigate the solution to (24) we develop $`\alpha `$ accordingly the antighost number
$$\alpha =\alpha _0+\alpha _1+\mathrm{},antigh\left(\alpha _k\right)=k,$$
(25)
where the last term from the sum can be assumed to be annihilated by $`\gamma `$. Because the free theory is irreducible, we can assume that $`\alpha `$ stops at antighost number one, i.e., $`\alpha =\alpha _0+\alpha _1`$, with $`\alpha _1=\alpha ^{a\mu }\eta _{a\mu }`$, where $`\alpha ^{a\mu }`$ pertains to $`H_1\left(\delta |d\right)`$, hence is a solution of the equation $`\delta \alpha ^{a\mu }+_\rho \lambda ^{a\rho \mu }=0`$. Like in the reducible case , $`H_2\left(\delta |d\right)`$ does not vanish, but the term $`\alpha _2`$ can be shown to vanish. Indeed, on the one hand $`\alpha _2`$ is of the form $`\alpha _2=\alpha ^{ab\mu \nu }\eta _{a\mu }\eta _{b\nu }`$, where $`\alpha ^{ab\mu \nu }`$ belongs to $`H_2\left(\delta |d\right)`$. On the other hand, the most general element in $`H_2\left(\delta |d\right)`$ reads as
$$\alpha ^a=C_{bc}^a\left(\eta ^{b\mu }A_\mu ^c+\frac{1}{2}\epsilon ^{\mu \nu \rho \sigma }B_{\mu \nu }^bB_{\rho \sigma }^c+g^{cd}\phi ^b_\mu A_d^\mu \right),$$
(26)
with $`g^{cd}`$ the inverse of $`g_{cd}`$, which further gives that $`\alpha ^{ab\mu \nu }=\alpha ^ah^{b\mu \nu }`$, where $`h^{b\mu \nu }`$ are some constants. By Lorentz covariance $`\alpha ^{ab\mu \nu }`$ must vanish, therefore $`\alpha _2`$ also vanishes. Let us investigate now the term $`\alpha _1`$. The general form of an object from $`H_1\left(\delta |d\right)`$ that is annihilated by $`\gamma `$ reads as
$$\alpha ^{a\mu }=C_{bc}^a\left(\phi ^bf^{c\mu }\left(A\right)+\epsilon ^{\rho \nu \lambda \mu }B_{\rho \nu }^bA_\lambda ^c\right),$$
(27)
where $`f^{c\mu }\left(A\right)`$ is a function of $`A_\mu ^a`$ and $`C_{bc}^a`$ are some constants, with $`C_{bc}^a=C_{cb}^a`$. It is simple to see that $`\delta \alpha ^{a\mu }=_\rho \left(\frac{1}{2}C_{bc}^a\epsilon ^{\rho \nu \lambda \mu }A_\nu ^bA_\lambda ^c\right)`$, so $`\alpha ^{a\mu }`$ is in $`H_1\left(\delta |d\right)`$. On the other hand, we obtain
$$\delta \alpha _1+\gamma \left(\frac{1}{2}C_{bc}^aB_a^{\mu \nu }A_\mu ^bA_\nu ^c\right)=_\mu \left(\frac{1}{2}C_{bc}^a\epsilon ^{\mu \nu \lambda \rho }A_\nu ^bA_\lambda ^c\eta _{a\rho }\right).$$
(28)
If we compare the last equation with (24) at antighost number zero (i.e., with the equation $`\delta \alpha _1+\gamma \alpha _0=_\mu n^\mu `$), it follows that
$$\alpha _0=\frac{1}{2}C_{bc}^aB_a^{\mu \nu }A_\mu ^bA_\nu ^c.$$
(29)
Thus, the deformed solution to order $`g`$ reads as
$`\overline{S}=S+g{\displaystyle }d^4x({\displaystyle \frac{1}{2}}C_{bc}^aB_a^{\mu \nu }A_\mu ^bA_\nu ^c+`$
$`C_{bc}^a(\phi ^bf^{c\mu }\left(A\right)+\epsilon ^{\rho \nu \lambda \mu }B_{\rho \nu }^bA_\lambda ^c)\eta _{a\mu }).`$ (30)
If we compute the antibracket $`(\overline{S},\overline{S})`$ we obtain
$$(\overline{S},\overline{S})=\frac{1}{3}g^2C_{[bc}^eC_{d]e}^a\epsilon ^{\mu \nu \lambda \rho }d^4xA_\mu ^bA_\nu ^cA_\lambda ^d\eta _{a\rho }g^2d^4xu,$$
(31)
where $`\left[bcd\right]`$ expresses the antisymmetry with respect to the indices between brackets. If we denote the term in $`g^2`$ from (22) by $`g^2d^4xb`$, then the interaction is consistent to order $`g^2`$ if and only if $`u=s_Ib+_\mu k^\mu `$ . However, from (31) we see that $`u`$ cannot be of that form, and so it must vanish. This means that the constants $`C_{bc}^a`$ must fulfill the Jacobi identity
$$C_{[bc}^eC_{d]e}^a=0,$$
(32)
hence must define the structure constants of a Lie algebra. In this situation (31) vanishes, so $`\overline{S}`$ (which is only of order $`g`$) is a solution of the master equation without adding higher order terms in $`g`$ (the vanishing of $`u`$ implies that all the higher order terms vanish).
The terms from $`\overline{S}`$ that do not involve the antifields, $`S_0^L[A_\mu ^a,B_a^{\mu \nu }]\frac{1}{2}C_{bc}^agd^4xB_a^{\mu \nu }A_\mu ^bA_\nu ^c`$, give nothing but the well-known action of the non-abelian Freedman-Townsend model, while $`\overline{S}`$ itself represents the corresponding solution to the master equation deriving from our irreducible BRST approach to this model. The terms from (On the ‘Irreducible’ Freedman-Townsend Vertex) that are linear in the antifields show that the deformed gauge transformations read as $`\overline{\delta }_ϵB_a^{\mu \nu }=\epsilon ^{\mu \nu \lambda \rho }\left(D_\lambda \right)_a^bϵ_{\rho b}`$, $`\overline{\delta }_ϵA_\mu ^a=0`$, $`\overline{\delta }_ϵ\phi _a=^\mu ϵ_{a\mu }+gC_{ab}^cf^{b\mu }\left(A\right)ϵ_{c\mu }`$, such that the gauge transformations for $`B_a^{\mu \nu }`$ and $`A_\mu ^a`$ take the familiar form in the literature. In the above formulas, the covariant derivative is defined by $`\left(D_\lambda \right)_a^b=\delta _a^b_\lambda +gC_{ac}^bA_\lambda ^c`$. In addition, we have derived a class of gauge transformations for $`\phi _a`$. We remark that the functions $`f^{c\mu }\left(A\right)`$ are still undetermined. They must be in such a way that the deformed gauge transformations are irreducible. A choice that preserves the irreducibility and, in the meantime, makes manifest the nice structure represented by the covariant derivative is $`f^{c\mu }\left(A\right)=A^{c\mu }`$, so $`\overline{\delta }_ϵ\phi _a=\left(D^\mu \right)_a^bϵ_{b\mu }`$. The solution (On the ‘Irreducible’ Freedman-Townsend Vertex) with $`f^{c\mu }\left(A\right)`$ replaced by $`A^{c\mu }`$ differs from that obtained in the literature by many authors in the reducible framework. The solution (On the ‘Irreducible’ Freedman-Townsend Vertex) does not contain terms that are quadratic in the antifields (like in the reducible situation), so the irreducible BRST transformations $`\overline{s}_IF=(F,\overline{S})`$ do not involve the antifields, such that the gauge-fixed BRST symmetry does not depend on the gauge-fixing fermion, by contrast with the reducible setting. In consequence, our approach leads to a gauge-fixed BRST symmetry that is off-shell nilpotent. Indeed, we have that $`\overline{s}_IB_a^{\mu \nu }=\epsilon ^{\mu \nu \lambda \rho }\left(D_\lambda \right)_a^b\eta _{b\rho }`$, $`\overline{s}_IA_\mu ^a=0`$, $`\overline{s}_I\phi _a=\left(D^\mu \right)_a^b\eta _{b\mu }`$, $`\overline{s}_I\eta _{a\mu }=0`$. In the meantime, the absence of the term quadratic in the antifields will consequently imply the absence of the three-ghost coupling term in the gauge-fixed action, such that the gauge-fixed action in the context of our irreducible approach takes a simpler form. This completes our irreducible procedure for deriving the Freedman-Townsend vertex.
To conclude with, in this letter we have exposed a cohomological approach to the Freedman-Townsend model consisting in two basic steps, namely, the construction of an irreducible BRST symmetry for the abelian version and the subsequent deformation of the irreducible theory. The results arising in our irreducible procedure prove the uniqueness of the Freedman-Townsend vertex in four dimensions (which has also been derived in , but within the reducible background) and also lead to a deformed solution of the master equation that has not previously been derived in the literature. In this light, our irreducible approach represents an efficient alternative to the reducible version exposed in .
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# Star formation properties of UCM galaxies
## 1 Introduction
The study of the evolution of the Star Formation Rate (SFR) of individual galaxies and the SFR history of the Universe has experienced considerable progress recently (see, e.g., Madau, Dickinson & Pozzeti 1998 and references therein). These are key observables needed to extend our understanding of galaxy formation and evolution. In the last few years, the combination of very deep ground-based and HST multi-band imaging with deep spectroscopic surveys carried out with 4-m and 10-m class telescopes has allowed the sketching of the SFR history of the Universe up to $`z>4`$ (see, e.g., Madau et al. 1998 and references therein).
A great deal of effort has been devoted to both observational and theoretical studies of star-forming objects and their evolution with look-back-time. Deep imaging and spectroscopy of faint galaxies at intermediate and high redhsifts have yielded vast amounts of quantitative information in this field (Lilly et al. 1995; 1998 and references therein; Driver, Windhorst & Griffiths 1995; Steidel et al. 1996; Lowenthal et al. 1997; Hammer et al. 1997; Hu, Cowie & McMahon 1998; see Ellis 1997 for a recent comprehensive review). Although substantial uncertainties still exist, a reasonably coherent picture is emerging. The Star Formation Rate density of the universe was probably about an order of magnitude higher in the past than it is now, perhaps peaking at $`z1`$$`2`$ (e.g., Gallego et al. 1995; Madau et al. 1996; Connolly et al. 1997; Madau et al. 1998). These observational results seem to be in good agreement with the predictions of recent theoretical models of galaxy formation (Pei & Fall 1995; Baugh et al. 1998; Somerville, Primack & Faber 1999), although the question of whether the SFR density decreased beyond $`z2`$ is still a matter of intense debate (Hu et al. 1998; Somerville et al. 1999; Hughes et al. 1998; Barger et al. 1998; Steidel et al. 1999).
Given the large redshift range covered by these studies, different SFR indicators have perforce been used, all of which have different calibrations, selection effects and systematic uncertainties. These indicators include emission line luminosities (e.g., H$`\alpha `$, H$`\beta `$, \[OII\]$`\lambda 3727`$Å), blue and ultraviolet fluxes, far-infrared and sub-mm fluxes, etc (see, e.g., Gallego et al. 1995; Rowan-Robinson et al. 1997; Tresse & Maddox 1998; Glazebrook et al. 1999; Treyer et al. 1998; Madau et al. 1996; Connolly et al. 1997; Hughes et al. 1998; Barger et al. 1998; see also Charlot 1998 and Kennicutt 1992). It would be highly desirable to use the same SFR indicator at all redshifts, so that the problems related to different selection effects and systematics could be avoided. It is widely accepted that the H$`\alpha `$ is one of the most reliable measurements of the current star formation rate (modulo the IMF; see, e.g., Kennicutt 1992). Several groups have used the H$`\alpha `$ line to estimate SFRs at different redshifts, from the local universe to beyond $`z=1`$ (Gallego et al. 1995; Tresse & Maddox 1998; Glazebrook et al. 1999), albeit with very different sample selection methods. Nevertheless, it is clear that it is now necesary to build sizeable samples of H$`\alpha `$-selected star-forming galaxies at different redshifts and study their properties. One would like to know the preferred sites of star formation in the local universe and beyond, and the main propeties of the star-forming galaxies and their evolution. Some questions that need to be answered include: does star formation mainly occur in dwarf, starbursting galaxies or in more quiescent, normal L galaxies? how has that evolved with time? what fraction of the stellar mass of the galaxies is being built by their current star-formation episodes?
Progress towards answering questions such as these requires, as a first step, a comprehensive study of the properties of the star-forming galaxies in the local universe. The Universidad Complutense de Madrid survey (UCM hereafter; Zamorano et al. 1994, 1996) is currently the most complete local sample of galaxies selected by their H$`\alpha `$ emission (see section 2). It has been used to determine the local H$`\alpha `$ luminosity function, the SFR function and the SFR density (Gallego et al. 1995). It is also widely used as a benchmark for high redshift studies (e.g., Madau et al. 1998 and references therein). Thus, the UCM survey provides a suitable sample of local star-forming galaxies for detailed studies.
Both optical imaging (Gunn-$`r`$; Vitores et al. 1996a, 1996b) and spectroscopy of the whole UCM sample (Gallego et al. 1996; GAL96 hereafter; see also Gallego et al. 1997) are already available. The optical data provides information on the current star-formation activity, but is rather insensitive to the past star-formation history of the galaxies. In this paper we present new near-infrared imaging observations for a representative subsample of UCM galaxies. The near infrared luminosities are sensitive to the mass in older stars, and therefore provide a measurement of the integrated past star formation in the galaxies and their total stellar masses (see, e.g., Aragón-Salamanca et al. 1993; Alonso-Herrero et al. 1996; Charlot 1998). Alonso-Herrero et al. (1996; AH96 hereafter) carried out a pilot study of similar nature with a very small sample. We will now extend the work to a galaxy sample that is large enough for statistical studies, and that is expected to represent the properties of the complete UCM sample and thus those of the local star-forming galaxy population. We will also improve the work of AH96 in two fronts: first, we will use up-to-date population synthesis models, and second, we will use a more sofisticated statistical technique when comparing observational data and model predictions.
In section 2 we briefly introduce the UCM sample. In section 3 the observations, reduction procedures, and data analysis are described. The evolutionary synthesis models are presented in section 4, and the results are described in section 5. Finally, section 6 contains a summary of this work.
## 2 UCM survey
The UCM survey is a wide-field objective-prism search for star-forming galaxies which used the H$`\alpha `$ emission line as main selection criterium (Zamorano et al. 1994, 1996). This survey was carried out at the 80-120cm Schmidt Telescope of the Calar Alto German-Spanish Observatory (Almería, Spain), using IIIaF photographic plates. The identification of the emission-line objects was done by visual inspection of the plates over the 471.4 square degrees that the survey covers. An automatic procedure for the detection has also been developed by Alonso et al. (1995, 1999) which avoids possible human subjectivities in the selection. The number of emission-line candidates found was 264, about 44 per cent of them previously uncatalogued. This yielded a detection rate of about 0.6 objects per square degree \[Zamorano et al. 1994\]. A total of 191 of these objects were confirmed spectroscopically by GAL96 as emission-line galaxies.
The wavelength cut-off of the photographic emulsion limits the redshift range spanned by the survey to H$`\alpha `$-emitting objects below $`z`$=0.045$`\pm `$0.005. The completeness tests performed (Vitores 1994; Gallego 1995) ensure that the Gunn-$`r`$ limiting magnitude of the whole sample is 16.5<sup>m</sup> with an H$`\alpha `$ equivalent width detection limit of about 20Å.
Details about the observations, data reduction, reliability and accessibility of the complete data set are summarized in Zamorano et al. (1994, 1996).
## 3 Observational data
### 3.1 Optical imaging
The complete description of the optical Gunn-$`r`$ \[Thuan & Gunn 1976\] observations and image reduction is given in Vitores et al. (1996a, 1996b). Briefly, these images were acquired during a total of eight observing runs from December 1988 through January 1992 using different CCD detectors on the CAHA/MPIA 2.2-m and 3.5-m telescopes, both at Calar Alto (Almería, Spain).
### 3.2 Near-infrared images
Near-infrared (nIR hereafter) images in the $`J`$ (1.2$`\mu `$m) and $`K`$ (2.2$`\mu `$m) or $`K^{}`$-bands (2.1$`\mu `$m), were obtained for 67 galaxies from the UCM survey during three observing runs at the Lick Observatory and one run at the Calar Alto Observatory.
The three Lick observing runs took place in 1996 (January 9–14, May 4–7 and June 7–9). We used the Lick InfraRed Camera (LIRC-II) equipped with a nicmos3 256$`\times `$256 detector on the 1-m telescope of the Lick Observatory (California, USA). The instrumental setup provided a total field of view of 2.4$`\times `$2.4 square arc minutes with a spatial scale of 0.57″ per pixel. For details about the LIRC-II camera and the nicmos3 detector used see Misch, Gilmore & Rank \[Misch et al. 1995\]. The Calar Alto observations were carried out in 1996 (August 4–6). We used the MAGIC camera with a nicmos3 256$`\times `$256 detector attached to the CAHA/MPIA 2.2-m telescope at Calar Alto (Almería, Spain). The field of view was 2.70$`\times `$2.70 square arc minutes and the spatial scale 0.63″ per pixel. Details about the MAGIC camera can be found in Herbst et al. \[Herbst et al. 1993\]. Images were obtained in the $`J`$ and $`K^{}`$ bands in all the observing runs, except for the January 9–12 one, when a standard $`K`$ filter was used. A $`K^{}`$ filter \[Wainscoat & Cowie 1992\] was used in order to reduce the thermal background introduced by the red wing of the standard $`K`$ passband.
The observational procedure followed was extensively described in Aragón-Salamanca et al. \[Aragón-Salamanca et al. 1993\]. Briefly, we subdivided the total exposure time required for each object in a number of images, offset by a few arcseconds, in order to avoid saturation. Also, blank sky images were obtained between consecutive object images for sky-subtraction and flat-fielding purposes with offests larger than $`1`$ arc minute. This procedure allows to reduce the effect of pixel-to-pixel variations, bad pixels, cosmic rays, and faint star images in the sky frames.
The reduction was carried out using our own IRAF<sup>1</sup><sup>1</sup>1IRAF is distributed by the National Optical Astronomy Observatories, which is operated by the Association of Universities for Research in Astronomy, Inc. (AURA) under cooperative agreement with the National Science Foundation. procedures, following standard reductions steps also described in Aragón-Salamanca et al. \[Aragón-Salamanca et al. 1993\], these included bad pixel removal, dark substraction, flat-fielding, and sky substraction. Finally, all the object images were aligned, combined and flux calibrated. Flux calibration was performed using standard stars from the lists of Elias \[Elias et al. 1982\] and Courteau \[Courteau 1995\] observed at airmasses close to those of the objects. The atmospheric extinction coefficients used were $`\kappa _J`$=0.102 mag/airmass and $`\kappa _K`$=0.09 mag/airmass, while independent zero-points were derived for each night. The $`K^{}`$-band magnitudes were converted into $`K`$-band magnitudes using the empirical relation given by Wainscoat & Cowie \[Wainscoat & Cowie 1992\], $`K^{}K`$=0.22$`\times (HK)`$. Based on the zero-redshift SEDs of Aragón-Salamanca et al. \[Aragón-Salamanca et al. 1993\] and the mean nIR colours given in AH96 for a small sample of UCM galaxies, we used an $`HK`$ colour of 0.3$`\pm `$0.1<sup>m</sup>. Thus, the $`K^{}K`$ correction was 0.07$`\pm `$0.02<sup>m</sup>.
Aperture photometry was carried out on the Gunn-$`r`$ and nIR images using the IRAF/APPHOT routines. We measured $`rJK`$ magnitudes and optical-nIR colours ($`rJ`$, $`JK`$) through several physical apertures, including three disk-scale lengths<sup>2</sup><sup>2</sup>2$`H_0`$=50 km s<sup>-1</sup> Mpc<sup>-1</sup> and $`q_0`$=0.5 have been assumed throughout this paper (see Vitores et al. 1996a). We also measured total $`K`$-band magnitudes using physical apertures large enough to ensure that all the light from the galaxy was included. Finally, we corrected for contamination from field stars by replacing affected pixels by adjacent sky counts. Circular aperture $`rJ`$ and $`JK`$ colours measured at three disk-scale lengths and integrated $`K`$-band magnitudes are given in Table LABEL:data. Magnitude and colour errors include both calibration and photometric uncertainties. In all the objects analyzed, except UCM0014$`+`$1748 and UCM1432$`+`$2645 (which was observed off-centre), the field covered by the detector was large enough to include three disk-scale length apertures. For these two galaxies we obtained the three disk-scale colours from their extrapolated growth curves. The differences between the larger measurable aperture and the extrapolated values were 0.05<sup>m</sup> for UCM0014$`+`$1748 and 0.1<sup>m</sup> for UCM1432$`+`$2645.
In our analysis, the $`rJK`$ magnitudes and optical-nIR colours are corrected for Galactic and internal extinction. Since most of the luminosity of these galaxies in the $`rJK`$ passbands comes from the stellar continuum, we have estimated the colour excesses on the continuum, $`E(BV)_{\mathrm{continuum}}`$, with the expression given by Calzetti, Kinney & Storchi-Bergmann (1996, see also Calzetti 1997a; Storchi-Bergmann, Calzetti & Kinney 1994),
$$E(BV)_{\mathrm{continuum}}=0.44\times E(BV)_{\mathrm{gas}}.$$
(1)
where the $`E(BV)_{\mathrm{gas}}`$ were obtained from the spectroscopic Balmer decrements measured by GAL96 (see below). We assumed a diffuse dust model that implies a total-to-selective extinction ratio of $`R_V`$=3.1 (see Mathis 1990; Cardelli, Clayton & Mathis 1989). Assuming a Galactic extinction curve, we obtained that $`A_r/A_V`$, $`A_J/A_V`$ and $`A_K/A_V`$ are 0.83, 0.28 and 0.11 respectively \[Mathis 1990\]. Note than when correcting the H$`\alpha `$ fluxes and equivalent widths we assume that the line emission comes from the gaseous component, but the continuum is mainly stellar. In table LABEL:data we present the observational data before correcting for extinction, together with the gas colour excesses needed for the correction.
### 3.3 Optical spectroscopy
Optical long-slit spectroscopy for the UCM survey was obtained by Gallego \[Gallego 1995\] at the 2.5-m Isaac Newton Telescope (INT) at Roque de los Muchachos Observatory, La Palma (Spain), and the 2.2-m and 3.5-m telescopes at Calar Alto (Almería, Spain), during a total of 10 observing runs. Details about the instrumental setups, slit widths, spatial scales and dispersions achieved are given in Table 1 of GAL96.
The line fluxes and equivalent widths of different emission lines are given in GAL96. Gas colour excesses, $`E(BV)_{\mathrm{gas}}`$, were obtained assuming a Galactic extinction curve and intrinsic ratios I$`(\mathrm{H}\alpha )/I(\mathrm{H}\beta )`$=2.86 and I$`(\mathrm{H}\gamma )/I(\mathrm{H}\beta )`$=0.468, which are the theoretical values expected for a low density gas with $`T_\mathrm{e}`$=10<sup>4</sup> K in Case B recombination \[Osterbrock 1989\]. We estimate that the effect of differential atmospheric refraction on the H$`\alpha /`$H$`\beta `$ ratio is in most cases (82% of the galaxies) below 5%. Only in four of the galaxies studied here the uncertainty in $`E(BV)_{\mathrm{gas}}`$ due to differential refraction is larger than 0.1<sup>m</sup>. The observation and analysis procedures followed by Gallego \[Gallego 1995\] —slit widths, position angles, spectrum extraction— ensure good integrated spectroscopic information for these objects. Errors in the EW(H$`\alpha `$+\[NII\]) have been estimated from the signal-to-noise and spectral resolution data given by GAL96. We have assumed a 100Å interval for the continuum fit range, $`\mathrm{\Delta }\lambda _{\mathrm{cont}}`$ \[Gallego 1995\], and a reciprocal dispersion of $`\mathrm{\Delta }\lambda `$3Å/pixel. Thus,
$$\mathrm{\Delta }\mathrm{EW}=\frac{1}{\mathrm{SNR}\sqrt{\mathrm{N}}}\sqrt{\mathrm{EW}^2+\mathrm{FWZI}^2}$$
(2)
where SNR is the signal-to-noise ratio of the continuum, N is the number of points used to determine the mean continuum flux, i.e. N=$`\mathrm{\Delta }\lambda _{\mathrm{cont}}/\mathrm{\Delta }\lambda `$30, and FWZI is the Full Width at Zero Intensity. The FWZI was computed as two times the Full Width at Half Maximum (FWHM) of the comparison arc lines. Typical FWHMs are about 12.5Å (for a 3″ wide slit with the R300V$`+`$IDS$`+`$Tek3 configuration). Equivalent widths (EWs hereafter) of H$`\alpha `$+\[NII\] and their corresponding errors are given in Table LABEL:data. The H$`\alpha `$ equivalent width data used in this work were corrected for contamination of the \[NII\]$`\lambda `$6548Å and \[NII\]$`\lambda `$6584Å emission lines using the \[NII\]$`\lambda `$$`\lambda `$6548,6584Å/H$`\alpha `$ line ratios given by GAL96 (see also Gallego 1995). Finally, we consider a correction of 2 Å in the EW(H$`\alpha `$) due to the H$`\alpha `$ underlying absorption in G-K giants (see, e.g., Kennicutt 1983).
## 4 Evolutionary synthesis models
Although the $`K`$-band luminosity can provide a very good estimate of the stellar mass in galaxies (Aragón-Salamanca et al. 1993; Charlot 1998), the contribution from red supergiants associated with recent star forming events may lead to the overestimate of the stellar mass when standard mass-luminosity relations are used. Thus, in order to estimate the relative contribution of the old underlying and young stellar populations to the magnitudes and colours measured, we have developed a complete set of evolutionary synthesis models. These models are based on those developed by AH96, but use the new population synthesis models of Bruzual & Charlot (private communication, BC96 hereafter), instead of the old Bruzual & Charlot (1993; BC93 hereafter) models. From the number of ionizing photons supplied by the BC96 models, we have also calculated the contribution of the hydrogen and helium emission-lines and nebular continuum to the optical and nIR passbands.
AH96 demonstrated that the properties of most of the UCM star-forming galaxies are better reproduced with instantaneous burst models rather that models with constant star formation. Therefore, we have computed the evolution with time of the optical-nIR colours and EW(H$`\alpha `$) of an instantaneous burst superimposed on a 15 Gyr old evolving population. A Scalo Initial Mass Function (IMF; Scalo 1986) with lower and upper mass cutoffs of $`M_{\mathrm{low}}`$=0.1 M and $`M_{\mathrm{up}}`$=125 M was adopted. The Cousins-$`R`$ magnitudes given by the BC96 models have been converted into Gunn-$`r`$ magnitudes using the relation $`r`$=$`R_\mathrm{C}+`$0.383$``$0.083$`\times (VR_\mathrm{C})`$ (Fernie 1983; Kent 1985).
In order to match the colours predicted by the BC96 models for a 15 Gyr old Single Stellar Population (SSP hereafter; $`rJ`$=2.09, $`JK`$=0.85) to those measured in the bulges of local relaxed spiral galaxies, assuming a negligible dust reddening (see Peletier & Balcells 1996; Fioc & Rocca-Volmerange 1999), we applied a small correction to our models
$`(rJ)^{\mathrm{obs}}`$=$`(rJ)_{15\mathrm{G}\mathrm{y}\mathrm{r}}^{\mathrm{mod}}0.03`$
$`(JK)^{\mathrm{obs}}`$=$`(JK)_{15\mathrm{G}\mathrm{y}\mathrm{r}}^{\mathrm{mod}}+0.06`$.
In addition, the stellar mass-to-light ratio predicted by the model for a 15 Gyr old stellar population, 1.34 M/L<sub>K,☉</sub>, was corrected to match that measured in local relaxed spiral galaxies, $``$1 M/L<sub>K,☉</sub> (see Héraudeau & Simien 1997, and references therein), using M<sub>K,☉</sub>=3.33 \[Worthey 1994\]. Since most UCM galaxies (about 83 per cent) are morphologically classified as Sa–Sc$`+`$ \[Vitores 1994\], we are confident that the corrections applied to the models are reasonable. In any case, these small corrections, intended to provide a good agreement between the model predictions and observations for the underlying stellar populations of the galaxies, do not affect significantly any of the conclusions of this paper.
The main parameters of our models are those inherent to the BC96 models (age, metallicity, IMF, etc), together with the strength of the current star-forming burst. The burst strength, $`b`$, is defined as the ratio of the mass of the newly formed stars to the total stellar mass of the galaxy \[Krüger, Fritze-v. Alvensleben & Loose 1995\]. We have explored models with metallicities between 1/50 Z and 2 Z, and burst strengths in the range 1–$`10^4`$ (in 0.04 dex steps). The new BC96 models (Scalo IMF) produce slightly redder $`rJ`$ and $`JK`$ colours and fewer Lyman photons than the BC93 ones (Salpeter IMF), for the same burst strength and solar metallicity.
## 5 Results
The main goal of this work is the characterization of the star formation activity of a representative sample of local galaxies. The properties of the current star-formation events and the host galaxies will be studied. In particular, we are interested in linking the properties of the local star-forming galaxies with those of galaxies forming stars at higher redshifts.
First, we study the completeness and representativeness of our sample in relation to the local star-forming galaxy population (see section 5.1). In section 5.2 we analyze the measured magnitudes and colours of the galaxies. Then, and in order to obtain the burst strengths, burst ages, stellar masses and specific star formation rates (SFR per unit mass; Guzmán et al. 1997; Lowenthal et al. 1997), we compare our data with evolutionary synthesis models (section 5.3). The comparison between data and models, and the determination of the best-fitting set of parameters, are not straightforward tasks, and some details of the method we have applied are described in the Appendix. Finally, we discuss the derived burst strengths, ages, metallicities, galaxy stellar masses and star formation rates derived for our sample (section 5.4, 5.5 and 5.6).
### 5.1 Sample completeness
Our nIR sample will suffer, first, from the intrinsic selection effects of the objective-prism$`+`$photographic plate technique used in the UCM survey. Those were discussed in detail in Zamorano et al. (1994, 1996), Vitores \[Vitores 1994\] and Gallego \[Gallego 1995\]. Briefly, the observational procedure employed limits the UCM sample to local galaxies with redshift lower than 0.045$`\pm `$0.005 and H$`\alpha `$ equivalent width larger than 20Å. The Gunn-$`r`$ limiting magnitude is 16.5<sup>m</sup> with a brigth end cut-off, due to saturation of the photographic plates, placed at $``$14.2<sup>m</sup>. However, additional selection effects may be present in our work due to the limited size of our nIR sample ($``$35 per cent of the UCM whole sample), thus it is necessary to ensure that the properties of this subsample are representative of those of the complete UCM survey.
In Figure 1a we compare the Gunn-$`r`$ apparent magnitude histogram of the whole UCM sample with that of the galaxies observed in the nIR. Although the apparent magnitude distributions match reasonably well, the objects in the nIR subsample tend to be marginally brighter than those in the UCM complete sample. The median $`m_r`$ for the UCM complete sample is 15.5<sup>m</sup>, while that of the nIR sample is 15.2<sup>m</sup> (see Vitores et al. 1996b). This may imply a small deficiency of low-luminosity and/or higher redshift galaxies. However, that does not seem to be a strong effect (see figures 1b and d).
From Figures 1b, 1c and 1e it is clear that the nIR sample represents about 35 per cent of the UCM complete sample \[Gallego 1995\] in every redshift, $`E(BV)`$ and EW(H$`\alpha `$) bin. A Kolmogorov-Smirnov test indicates that both samples show similar distributions in $`z`$, $`E(BV)`$ and EW(H$`\alpha `$), with probabilities of 45, 93 and 87 per cent, respectively. For the $`r`$-band absolute magnitudes and H$`\alpha `$ luminosities the probabilities are 25 and 47 per cent, respectively (see Figures 1d and 1f). Thus the galaxies in the nIR sample seem to be a fair subsample of the UCM complete sample in their global properties.
The only small difference arises when comparing in detail the spectroscopic type distributions of the star-forming galaxies in the nIR and UCM complete samples. There is a small deficiency of HII-like galaxies (see ahead) in the nIR subsample relative to the whole UCM sample. About 30 per cent of the UCM whole sample are HII-like galaxies, whereas only 19 per cent are present in the nIR subsample. Consequences of such a limitation will be taken into account in further discussions.
### 5.2 Aperture magnitudes and colours
#### 5.2.1 Mean colours
Global colours (obtained inside three disk-scale lengths), together with the integrated $`K`$-band magnitudes (not corrected for extinction) and the $`E(BV)_{\mathrm{gas}}`$ colour excesses, are given in Table LABEL:data. In GAL96, the galaxies in the UCM sample were classified in different morphological and spectroscopic classes (listed in table LABEL:data). We will briefly describe them here (see GAL96 for details):
SBNStarburst Nuclei— Originally defined by Balzano \[Balzano 1983\], they show high extinction values, with very low \[NII\]/H$`\alpha `$ ratios and faint \[OIII\]$`\lambda `$5007 emission. Their H$`\alpha `$ luminosities are always higher than 10<sup>8</sup> L.
DANSDwarf Amorphous Nuclear Starburst— Introduced by Salzer, MacAlpine & Boroson \[Salzer et al. 1989\], they show very similar spectroscopic properties to SBN objects, but with H$`\alpha `$ luminosities lower than 5$`\times `$10<sup>7</sup> L.
HIIHHII Hotspot— The HII Hotspot class shows (see GAL96) similar H$`\alpha `$ luminosities to those measured in SBN galaxies but with large \[OIII\]$`\lambda `$5007/H$`\beta `$ ratios, that is, higher ionization.
DHIIHDwarf HII Hotspot—- This is an HIIH subclass with identical spectroscopic properties but H$`\alpha `$ luminosities lower than 5$`\times `$10<sup>7</sup> L.
BCDBlue Compact Dwarf— Finally, the lowest luminosity and highest ionization objects have been classified as Blue Compact Dwarf galaxies, showing in all cases H$`\alpha `$ luminosities lower than 5$`\times `$10<sup>7</sup> L. They also show large \[OIII\]$`\lambda `$5007/H$`\beta `$ and H$`\alpha `$/\[NII\]$`\lambda `$6584 line ratios and intense \[OII\]$`\lambda `$3727 emission.
In our analysis, we separate the galaxies in two main categories: starburst disk-like (SB hereafter) and HII-like galaxies (see Guzmán et al. 1997; Gallego 1998). The SB-like class includes SBN and DANS spectroscopic types, whereas the HII-like includes HIIH, DHIIH and BCD type galaxies.
In order to determine representative mean optical-nIR colours for each galaxy group, we have assumed Gaussian probability distributions for the $`rJ`$ and $`JK`$ colours and EW(H$`\alpha `$) with the centres and widths ($`\sigma `$) given in Table LABEL:data. We have weighted the data for each galaxy with their corresponding errors when determining the mean values.
The HII-like objects seem to be on average 0.2<sup>m</sup> bluer in $`rJ`$ and 0.1<sup>m</sup> in $`JK`$ than the SB galaxies (see Table LABEL:datamean). Since the mean colour excess of the SB population ($`\overline{E(BV)}`$=0.7<sup>m</sup>) is 0.2<sup>m</sup> higher than that of the HII-like galaxies, these colour differences are even more significant when data not corrected for extinction are used: the differences for the un-corrected colours are 0.35<sup>m</sup> in $`rJ`$ and 0.15<sup>m</sup> in $`JK`$. K-S tests performed on the SB-like and HII-like objects indicate that both subsamples arise from independent distributions with a probability of 99.7 and 99.9 per cent respectively for the un-corrected $`rJ`$ and $`JK`$ colours.
Finally, whereas more than 60 per cent of the HII-like objects show not corrected for extinction equivalent widths of H$`\alpha `$ higher than 120 Å, only 3 per cent of the SB galaxies do. The relatively low EW(H$`\alpha `$) detection limit estimated for the UCM survey ($``$ 20 Å; Gallego 1995) ensures that the difference in EW(H$`\alpha `$) between SB and HII-like galaxies is not due to selection effects. A K-S test gives a probability of 99.9 per cent that these samples have independent EW(H$`\alpha `$) distributions. These differences in colours and H$`\alpha `$ equivalent widths are probably related to differences in their evolutionary properties (typical starburst age, starburst strength, specific star formation rate, etc) between both galaxy types (see section 5.4 and section 5.6).
#### 5.2.2 Colour-colour and colour-EW(H$`\alpha `$) diagrams
In Figure 2 we show colour-colour ($`rJ`$ vs. $`JK`$) and colour-EW(H$`\alpha `$) plots for the nIR sample. The offset between the position of the star forming galaxies (filled circles) in the $`rJ`$$`JK`$ plane (Figure 2a) and the bulges and disks of relaxed nearby spirals \[Peletier & Balcells 1996\] indicates the existence of ongoing star formation. Error bars in Figure 2a represent $`\pm `$1$`\sigma `$ errors. In Figures 2a and 2b we plot solar metallicity models with burst strengths, 10<sup>-3</sup>, 10<sup>-2</sup>, 10<sup>-1</sup>, and 1. In the case of Figures 2c and 2d, models with 10<sup>-1</sup> burst strength and different metallicities are displayed (cf. section 4).
Figures 2a and 2c show that changes in the optical-nIR colours due to changes in burst strength and age are much more significant than those produced by changes in metallicity. This fact is also observed in the colour-EW(H$`\alpha `$) diagrams (Figures 2b and 2d), especially in the case of sub-solar metallicity models. It is thus clear that it is in principle possible to infer burst strengths and ages from these diagrams, but the determination of metallicities would be very uncertain.
### 5.3 Determination of the physical properties of the galaxies
For each individual galaxy we have information on its $`rJ`$ and $`JK`$ colours and H$`\alpha `$ equivalent width. Thus, each galaxy has a point associated in the $`rJ`$, $`JK`$, 2.5$`\times `$log EW(H$`\alpha `$) three-dimensional space. However, due to the calibration and photometric errors, the uncertainty in these measurements transforms these points into probability distributions. As in section 5.2.1, we will assume Gaussian probability distributions for the $`rJ`$, $`JK`$ colours and 2.5$`\times `$log EW(H$`\alpha `$) with the centres and widths ($`\sigma `$) given in Table LABEL:data. The evolutionary synthesis models that we will associate with each galaxy probability distribution, will also follow different tracks in this three-dimensional space.
The three-dimensional probability distributions ($`rJ`$,$`JK`$,2.5$`\times `$log EW(H$`\alpha `$)) have been reproduced using a Monte Carlo simulation method. A total of 10<sup>3</sup> data points were generated in order to reproduce this distribution for each galaxy. No significant differences were observed using a larger number (e.g., 10<sup>4</sup>) of input particles. We estimated the model that better reproduces the colours and EW(H$`\alpha `$) for each of the 10<sup>3</sup> test particles applying a maximum likelihood method. The maximum likelihood estimator used, $``$, includes two colour terms and an EW(H$`\alpha `$) term. Thus,
$$(t,b,Z)=\underset{n=1}{\overset{3}{}}\frac{1}{\sqrt{2\pi }\mathrm{\Delta }C_n}\mathrm{exp}\left(\frac{(c_nC_n)^2}{2\mathrm{\Delta }C_n^2}\right)$$
(3)
where $`C_1`$, $`C_2`$ and $`C_3`$ are the $`rJ`$ and $`JK`$ colours and 2.5$`\times `$log EW(H$`\alpha `$) and $`\mathrm{\Delta }C_1`$, $`\mathrm{\Delta }C_2`$ and $`\mathrm{\Delta }C_3`$ are their corresponding errors. The $`c_n`$ coefficients are the $`rJ`$, $`JK`$ and 2.5$`\times `$log EW(H$`\alpha `$) values predicted by a given model. A similar estimator was employed by Abraham et al. \[Abraham et al. 1999\] for a sample of intermediate-$`z`$ HDF \[Williams et al. 1996\] galaxies.
Finally, we obtained the age, $`t`$, burst strength, $`b`$, and metallicity, $`Z`$, of the model that maximizes $``$ for each test particle of the ($`rJ`$,$`JK`$,2.5$`\times `$log EW(H$`\alpha `$)) probability distribution. Therefore, this procedure effectively provides the ($`t`$,$`b`$,$`Z`$) probability distribution for each input galaxy.
The resulting ($`t`$,$`b`$,$`Z`$) probability distributions are in many cases multi-peaked. Instead of analyzing these probability distributions as a whole, we have studied the clustering pattern present in the ($`t`$,$`b`$,$`Z`$) solution space. We have used a single linkage hierarchical clustering method (see Murtagh & Heck 1987; see also Appendix), which allows to isolate different solutions in the ($`t`$,$`b`$,$`Z`$) space. We have recovered the three most representative solution clusters for each galaxy. In Table LABEL:tablafin we show the mean properties of those solution clusters with probability higher than 20 per cent. This probability is computed as the number of test particles in a given cluster over the total number of test particles (10<sup>3</sup>). The errors shown in Table LABEL:tablafin correspond to the standard deviation of the data for each solution cluster. In those cases where all the solutions within a cluster yield the same age, burst strength, metallicity or mass, no errors were given.
The subsequent statistical analysis of each of the ($`t`$,$`b`$,$`Z`$) clusters indicates that significant correlations between $`t`$, $`b`$ and $`Z`$ are present. We have performed a principal component analysis (PCA hereafter; see Morrison 1976; see also Appendix) of the individual clusters given in Table LABEL:tablafin. The orientation of the first PCA component and the contribution of this component to the total variance within the solution cluster are also given.
After applying this procedure to the observed sample, only three galaxies (UCM1440$`+`$2521S, UCM1506$`+`$1922 and UCM1513$`+`$2012; see Figure 3a and 3b) show $`_{\mathrm{max}}`$$`<`$10.0. Note that a value $`_{\mathrm{max}}=10.0`$ corresponds to a model where the differences between the observed data and the model predictions equal the measurement errors, assuming mean errors of 0.12<sup>m</sup> in $`(rJ)`$ and $`(JK)`$ and 10 per cent in EW(H$`\alpha `$). This indicates that the range of model predictions covers reasonably well the observed properties of the galaxies. In Figure 3a and 3b we plot the differences between the colours and EW(H$`\alpha `$) measured and those predicted by the best-fit model. These differences have been calculated for the central values of the $`rJ`$, $`JK`$ and 2.5$`\times `$log EW(H$`\alpha `$) probability distributions.
Figure 3a shows that, in some cases, the models predict bluer $`JK`$ colours than those measured. At first sight, these discrepancies could be explained if we were underestimating the extinction correction factors applied to these objects. However, since the extinction correction affects $`rJ`$ more than $`JK`$, applying a higher extinction correction would destroy the good agreement between the observed and model $`rJ`$ colours. We have quantified the effect of a change in the correction for extinction assumed for our galaxies by comparing data de-reddened using $`E(BV)_{\mathrm{continuum}}`$=$`E(BV)_{\mathrm{gas}}`$ with our models. Figure 3c shows that data corrected using the relation given by Calzetti et al. \[Calzetti et al. 1996\] fit the models better than data corrected assuming $`E(BV)_{\mathrm{continuum}}`$=$`E(BV)_{\mathrm{gas}}`$. In addition, none of the three galaxies given above show higher $`_{\mathrm{max}}`$ values after using $`E(BV)_{\mathrm{continuum}}`$=$`E(BV)_{\mathrm{gas}}`$. Thus, we are confident that the extinction correction applied provides a reasonable fit to the models, and the discrepancies in $`JK`$ between the data and the models are probably due to inherent uncertainties in the modelling of the nIR continuum by the Bruzual and Charlot code.
The cluster analysis performed indicates that the clustering in the solution space is basically produced by the discretization in metallicity of the models. Fortunately, in many of the objects ($``$30 per cent; see Table LABEL:tablafin), only one ($`t`$,$`b`$,$`Z`$) solution cluster is able to reproduce the observables. About 33 per cent need two solutions and three solutions are needed for the remaining 37 per cent. The goodness of this comparison method, given by the number and size of statistically significant ($`t`$,$`b`$,$`Z`$) solution clusters, basically depends on the particular position of the object in the ($`rJ`$,$`JK`$,2.5$`\times `$log EW(H$`\alpha `$)) space and on its measurement errors. Fortunately, in those cases where the ($`t`$,$`b`$,$`Z`$) probability distribution is multi-valuated, the different solution clusters give similar burst strengths and total stellar masses. This is another manifestation of the fact that, as we saw in section 5.2.2, our data is not very sensitive to metallicity, and we will not attempt to derive it. Nonetheless, the burst ages are affected somewhat by small changes in metallicity, and frequently show wider distributions than the burst strenghts.
Finally, the PCA performed on each solution cluster suggests that the best-axis, given by the vector ($`u_t`$,$`u_{\mathrm{log}(b)}`$,$`u_{\mathrm{log}(Z)}`$)=($`u_x`$,$`u_y`$,$`u_z`$) shown in Table LABEL:tablafin, is commonly placed in the $`u_x`$-$`u_y`$ ($`t`$-$`b`$) plane and obeys $`u_x`$$``$$`u_y`$. This implies that age and burst strength are in many cases degenerated and, therefore, the properties of an individual object can be reproduced both with a young, low burst strength or an old, high burst strength model, within the ranges given in Table LABEL:tablafin.
### 5.4 Burst strengths and ages
In Table LABEL:tablafin we give mean burst strengths and ages for the individual solutions with probability higher than 20 per cent. Errors given are the standard deviations of the data points in each solution. For the stellar masses, the error related with the uncertainty in the $`K`$-band absolute magnitude determination is also given (in parenthesis). Using these probability distributions we have derived the burst strength, age, mass and metallicity frequency histograms for the whole sample as well as for the SB-like and HII-like galaxies (see Figures 4a–d). The number of points in the $`y`$ axis of these figures corresponds to the number of Monte Carlo test particles with a given burst strength, age, mass or metallicity within the accepted high-probability solutions.
This analysis yields a typical burst strength of 2$`\times `$10<sup>-2</sup> with approximately 90 per cent of the sample having burst strengths between 10<sup>-3</sup> and 10<sup>-1</sup>. Only seven objects in the sample show burst strengths higher than 10<sup>-1</sup> with a probability larger than 50 per cent, UCM0003$`+`$2200, UCM0145$`+`$2519, UCM1257$`+`$2808, UCM1259$`+`$3011, UCM1308$`+`$2958, UCM1432$`+`$2645 and UCM1440$`+`$2511.
Although the properties of the local star-forming galaxies seem to be well reproduced with an episodic star formation history (see also AH96), some of these objects may have evolved under more constant star formation rates (Glazebrook et al. 1999; Coziol 1996). In those objects the instantaneous burst assumption could yield very high burst strengths.
The burst strength histograms shown in Figure 4a give typically larger burst strengths for the HII-like objects, especially for the DHIIH and BCD type galaxies, than for the SB-like (see also Table LABEL:resmean). This segregation in burst strength is probably related to the difference in mean EW(H$`\alpha `$) pointed out in section 5.2.1 (see Table LABEL:datamean). In Table LABEL:resmean we also show the burst strengths and ages derived under the unrealistic assumption that the continuum extinction is as high as that measured for the ionized gas (see section 5.3).
The distribution of the burst ages is shown in Figure 4b. Since the probability of detection increases with EW(H$`\alpha `$) in objective-prism surveys (see García-Dabó et al. 1999), and the EW(H$`\alpha `$) continuously decreases with the the burst age, the number of objects detected with old burst ages is expected to be lower than with young ages, as observed. This behaviour is observed at ages older than 4 Myr, both for the SB and HII-like galaxies.
However, one would expect a reasonably flat distribution in the number of objects with young ages if the sample selection depended only on the H$`\alpha `$ equivalent width. But other factors such as the H$`\alpha `$ flux and continuum luminosity play an important role (see García-Dabó et al. 1999). Moreover, in our models we have estimated the H$`\alpha `$ luminosity ($`L_{\mathrm{H}\alpha }`$) from the number of Lyman continuum photons (Brocklehurst 1971) assuming that no ionizing photons escape from the galaxies. If some Lyman photons escape, the predicted H$`\alpha `$ luminosity would be lower and the derived ages could be significantly younger. Recent studies estimate the fraction of Lyman photons escaping from starburst galaxies to be about 3 per cent \[Leitherer et al. 1995\]. Bland-Hawthorn & Maloney \[Bland-Hawthorn & Maloney 1997\] estimated this quantity to be about 5 per cent for the Milky Way. Another feasible explanation could be that a significant fraction of these Lyman photons is absorbed by dust within the ionized gas (see, e.g., Armand et al. 1996). Both mechanisms would produce lower H$`\alpha `$ equivalent widths than those predicted by the standard super-ionizing models, and could explain the paucity of young star-forming bursts in Figure 4b. In this figure (dotted line) we also show the age distribution obtained assuming that 25 per cent of the Lyman continuum photons are missing. This distribution yields a larger number of objects at ages younger than 3 Myr, and a very steep decay at ages older than 4–5 Myr.
Finally, Bernasconi & Maeder \[Bernasconi & Maeder 1996\] have argued that, during the first 2–3 Myr in the main-sequence, stars more massive than 40 M, are still accreting mass embedded in the molecular cloud, and do not contribute to the ionizing radiation. Therefore, due to this reduction in the number of Lyman photons, the predicted H$`\alpha `$ equivalent widths below 2-3 Myr will be significantly lower and the ages deduced for the bursts should be younger.
### 5.5 Total stellar masses
In order to determine the total galaxy stellar mass we have assumed that the burst strengths and mass-to-light ratios derived from our models at three disk scale-length apertures are representative of the galaxy global properties. Thus, using these $`K`$-band mass-to-light ratios and the total $`K`$-band absolute magnitudes we have obtained stellar masses for the whole sample.
The inferred galaxy stellar masses derived depend, in principle, on four quantities: the galaxy $`K`$-band absolute magnitude, burst strength, and the mass-to-light ratios of the burst and the old underlying population. Since the derived burst strengths are very low ($``$10<sup>-2</sup>), the total mass-to-light ratios are dominated by the old stellar component. In fact, the ratio of the $`K`$-band luminosity of the old and young stellar populations is $``$20 for $`t`$=4 Myr, 4 for $`t`$=8 Myr and 7 for $`t`$=15 Myr (for $`Z`$=$`Z_{\mathrm{}}`$). Moreover, the absolute age of the old stellar component has a very small effect on the $`K`$-band mass-to-light ratio: there is only a 0.1 dex difference between 10 Gyr and 15 Gyr for solar metallicity, and the difference is even lower in the case of sub-solar metallicity models. In Figure 5 we show that the derived $`K`$-band stellar mass-to-light ratios span a very narrow range. Although statistically the SB- and HII-like mass-to-light ratio distributions are different with a probability of 95.3 per cent (from a K-S test), the difference in absolute value is only minor: the median mass-to-light ratios are 0.93 M/L<sub>K,☉</sub> for the whole sample, and 0.93 M/L<sub>K,☉</sub> and 0.91 M/L<sub>K,☉</sub>, for the SB and HII-like objects respectively. Consequently, the derived mass values mainly depend on the $`K`$-band absolute magnitude. In the upper-panel of Figure 5 we also show the range in $`K`$-band mass-to-light ratios given by Worthey \[Worthey 1994\] for 12 Gyr old modeled ellipticals. Thus, we can conclude that the $`K`$-band luminosity is a very good tracer of the stellar mass for both old stellar populations and local star-forming objects.
The distribution of stellar masses is shown in Figure 4c. This frequency histogram indicates that a typical star forming galaxy in our Local Universe has a stellar mass of about 5$`\times `$10<sup>10</sup> M. This value is somewhat lower than the stellar mass expected for a local L galaxy. Assuming M$`{}_{K}{}^{}{}_{}{}^{}`$=$``$25.1 (for H<sub>0</sub>=50 km s<sup>-1</sup> Mpc<sup>-1</sup>; Mobasher, Sharples and Ellis, 1993) and a $`K`$-band mass-to-light ratio of 1 M/L<sub>K,☉</sub> \[Héraudeau & Simien 1997\], the stellar mass inferred for an L galaxy is about 2$`\times `$10<sup>11</sup> M. Thus, star-forming galaxies in the local universe are typically a factor 4 less massive than L galaxies.
In addition, a clear offset between the stellar mass histograms of the SB and HII-like objects is seen in Figure 4c. The distributions of their stellar masses are centred at 7$`\times `$10<sup>10</sup> and 2$`\times `$10<sup>10</sup> M respectively (Figure 4c). This difference is even more significant, about 1 dex, when DHIIH and BCD spectroscopic types (Dwarfs in Table LABEL:resmean) and SBN galaxies are compared. A K-S test analysis of the SB-like and HII-like objects indicates that these two samples come from different age, burst strength and stellar mass distributions with probabilities 98.8, 77.1 and 99.9 per cent respectively.
In Table LABEL:resmean we also present the mean properties that would be obtained using $`E(BV)_{\mathrm{continuum}}=E(BV)_{\mathrm{gas}}`$. In this case, although we obtain important differences in age and burst strength, very similar stellar masses are derived since the stellar mass depends mainly on the $`K`$-band magnitude, only weakly affected by extinction.
### 5.6 Star formation rates
Since the star formation activity in the UCM galaxies is better described as a succession of episodic star formation events rather than continuous star formation (see also AH96), the current star formation rate (SFR) is not a meaningful quantity: the latest star formation event might have finished in many of the galaxies, and their current SFR would be zero. However, these galaxies have substantial H$`\alpha `$ luminosities, and it is accepted that the H$`\alpha `$ luminosity is a good measurement of the current SFR. In AH96 we showed that this is true, in a statistical sense, for a population of galaxies undergoing a series of star formation events, and we defined an ‘effective’ present-day SFR which coincides with the SFR we would derive if the galaxies were forming stars at a constant rate, producing the same mass in new stars as the ensemble of all the star-formation episodes (see AH96 for details). Here we will follow the same approach.
When estimating star formation rates (SFRs) in AH96 we used the BC93 models and a Salpeter IMF. In the present work, we have used the updated BC96 models with a Scalo IMF. Since the number of Lyman photons (N<sub>Ly</sub>) predicted by the old models is about 0.94 dex higher than that predicted by the new ones (for solar metallicity and ages lower than 16 Myr), we need to re-compute the relation between the H$`\alpha `$ luminosity, $`L_{\mathrm{H}\alpha }`$, and star formation rate. In addition, we will investigate the change produced in this relation using different metallicity models.
In order to compute the $`L_{\mathrm{H}\alpha }`$/SFR ratio, we have used a very similar procedure to that employed by AH96: we simulated a population of galaxies undergoing random bursts of star-formation and computed their total H$`\alpha `$ luminosities and the mass in newly-formed stars. However, instead of considering a uniform age and burst strength probability distribution, we have considered the burst strength, age and metallicity distributions for our galaxy sample. We used 67$`\times `$1000 points in order to reproduce this distribution in our Monte Carlo simulations. The SFR was computed as the ratio between the stellar mass produced in the burst, i.e. $`b`$$`\times `$M, and the maximum age for which we could have detected the galaxy in the UCM sample, that is, the time while EW(H$`\alpha `$)$`>`$20Å \[Gallego 1995\]. The $`L_{\mathrm{H}\alpha }`$/SFR ratios obtained are shown in Figure 6 for different metallicities and for the total ($`t`$,$`b`$,$`Z`$) distribution. The mean, median, and standard deviation values are given in Table LABEL:fig6t.
Since the changes in this ratio for different metallicity models are quite small, we have adopted the median value of the whole distribution in order to determine the ’effective’ SFR of the galaxies from our sample. The difference between the value adopted here and that of AH96 is about 1 dex, which is very close to the difference in the number of Lyman photons predicted by the BC93 and BC96 models, as expected.
Therefore, we have evaluated the current ’effective’ SFR using the expression
$$\mathrm{SFR}=\frac{L_{\mathrm{H}\alpha }}{1.7\times 10^{40}\mathrm{erg}\mathrm{s}^1}\mathrm{M}_{\mathrm{}}\mathrm{yr}^1$$
(4)
This expression assumes that every Lyman photon emitted effectively ionizes the surrounding gas. If, however, as is suggested in Section 5.4, we consider a fraction of non-ionizing Lyman photons of 25 per cent, the star formation rates computed should be 0.1 dex higher.
Specific star formation rates (SFR per unit mass; Guzmán et al. 1997) have been obtained using these SFR values and the stellar masses given by the highest probability solution cluster in Table LABEL:tablafin. The mean SFR and specific SFR for SB-like, HII-like, and whole sample are given in Table LABEL:tablasfrs (see also Figure 4d).
The specific SFR vs. stellar mass diagram is shown in Figure 7 (see Guzmán et al. 1997). In panel-a we show the stellar masses and star formation rates per unit mass for three reference samples. We have included the sample of Kennicutt (1983, K83 hereafter), taking the H$`\alpha `$ and $`B`$-band luminosities given by K83 and the stellar mass-to-light ratios of Faber & Gallagher \[Faber & Gallagher 1979\]. In addition, the sample of HII-galaxies of Telles \[Telles 1995\] was included, after converting virial masses to stellar masses using a correction of 0.6 dex (Gallego et al. 1999, in preparation) and assuming the H$`\beta `$-to-H$`\alpha `$ luminosity ratios used by Guzmán et al. \[Guzmán et al. 1997\] for this sample. Masses and specific SFRs for the Calzetti \[Calzetti 1997b\] sample are also shown. In this case, stellar masses were inferred subtracting the HI mass from the dynamical mass measured. The SFR values for the Calzetti \[Calzetti 1997b\] sample were obtained from their Br$`\gamma `$ luminosities assuming $`L_{\mathrm{H}\alpha }`$=102.8$`\times `$$`L_{\mathrm{Br}\gamma }`$ (Osterbrook 1989, for $`T_e`$=10<sup>4</sup> K and $`n_e`$=100 cm<sup>-3</sup>). Finally, the dwarf irregular galaxy GR8 \[Reaves 1956\] was included. Its H$`\alpha `$ luminosity was obtained from the H$`\beta `$ luminosity given by Gallagher, Hunter & Bushouse \[Gallagher, Hunter & Bushouse 1989\], assuming $`L_{\mathrm{H}\alpha }`$/$`L_{\mathrm{H}\beta }`$=2.86, and its stellar mass, 3.2$`\times `$10<sup>6</sup> M, from Carignan, Beaulieu & Freeman \[Carignan, Beaulieu & Freeman 1990\]. The limits in the H$`\alpha `$ luminosity function of Gallego et al. \[Gallego et al. 1995\], 10<sup>40.4</sup>–10<sup>42.8</sup> erg s<sup>-1</sup>, are also drawn.
Figure 7 shows that the UCM sample clearly represents a bridge between relaxed spiral galaxies and the most extreme HII galaxies from Telles \[Telles 1995\], that is, Sp$``$SB-like$``$HII-like$``$HII galaxies. In fact, some of the HII galaxies from Telles \[Telles 1995\] have very similar properties to those of the less massive HII-like UCM galaxies, mainly DHIIH and BCD spectroscopic types, very rare in our sample (see section 5.1). In addition, most of the SBN type UCM galaxies seem to be normal late-type spirals with enhanced star formation. This star formation enhancement is about a factor of three, and is due to the ongoing nuclear starburst. Thus, the range in specific SFR spanned by the population of the star-forming galaxies that dominate the SFR in the local universe is (10–10<sup>3</sup>)$`\times `$10<sup>-11</sup> yr<sup>-1</sup>, from the local relaxed spirals to the most extreme HII galaxies. In fact, this range is not very different from that obtained by Guzmán et al. \[Guzmán et al. 1997\] for a sample of intermediate/high-$`z`$ compact galaxies from the HDF. The high specific SFR region, where the HII galaxies from Telles \[Telles 1995\] are placed, is not very well covered by our sample due to the scarcity of very low-luminosity objects, basically DHIIH and BCD spectroscopic type galaxies, relative to the UCM whole sample.
## 6 Summary
Using new nIR observations and published optical data for 67 galaxies from the Universidad Complutense de Madrid (UCM) survey, we have derived the main properties of their star-forming events and underlying stellar populations. This sample represents about 35 per cent of the UCM galaxies covering the whole range of absolute magnitudes, H$`\alpha `$ luminosities and equivalent widths spanned by the survey. Burst strengths and ages, stellar masses, stellar mass-to-light ratios and, to a certain extent, metallicities, have been obtained by comparing the observed $`rJ`$ and $`JK`$ colours, $`K`$-band magnitudes, and H$`\alpha `$ equivalent widths and luminosities with those predicted by evolutionary synthesis models. The comparison of the observations with the model predictions was carried out using a maximum-likelihood estimator in combination with Monte Carlo simulations which take into account the observational uncertainties. Our main results are:
1. The star-forming galaxies in the UCM sample (used to determine the SFR density of the local universe), show typical burst strengths of about 2 per cent and stellar masses of 5$`\times `$10<sup>10</sup> M. The current star formation in these galaxies is taking place in discrete star formation events rather than in a continuous fashion. If this is typical of the past star-formation history in the galaxies, many of such star formation events would be necessary to build up their stellar mass. However, our observations provide very little information on star-formation episodes that took place before the current one.
2. We have identified two separate classes of star-forming galaxies in the UCM sample: SB-like and HII-like galaxies. Within the HII-like class the DHIIH and BCD spectroscopic type galaxies, i.e. dwarfs, constitute the most extreme case. The mean burst strength deduced for the SB-like galaxies is about a 25 per cent lower than for the dwarf HII-like galaxies. The average stellar mass is an order of magnitude larger in the former than in the latter. The SB-like galaxies are relatively massive galaxies where the current star formation episode is a minor event in the build up of their stellar masses, while HII-like galaxies are less massive systems in which the present star formation could dominate in some cases their observed properties and contributes to a greater extent to their stellar population.
3. Because of the low burst strengths inferred, $`K`$-band luminosity is dominated by the old stellar populations, and the $`K`$-band stellar mass-to-light ratio is almost the same (within $``$20 per cent) for all the galaxies. Thus, the $`K`$-band luminosity is a very good estimator of the stellar mass for typical star-forming galaxies.
4. The average SFR of the galaxies is log(SFR)$``$1.5, with the SFR expressed in M yr<sup>-1</sup>, and it is similar for the SB-like and the HII-like galaxies. However, since the latter are typically less massive, their specific SFR (SFR per unit stellar mass) is significantly larger, in a factor 2.3, than that of the former.
5. The UCM galaxies represent a bridge in specific SFR between relaxed spirals and extreme HII galaxies. The range in specific star formation rate spanned by the local star-forming galaxies, (10–10<sup>3</sup>)$`\times `$10<sup>-11</sup> yr<sup>-1</sup>, is very similar to that observed in higher redshift objects.
## APPENDIX: Analysis of the space of solutions
In this appendix we briefly describe the analysis performed onto the ($`t`$,log $`b`$,log $`Z`$) space of solutions. For each galaxy we have 10<sup>3</sup> ($`t`$,log $`b`$,log $`Z`$) points that correspond to the 10<sup>3</sup> points generated in the ($`rJ`$,$`JK`$,2.5$`\times `$log EW(H$`\alpha `$)) space using a Monte Carlo simulation method.
The mean values ($`<t>`$,$`<`$log $`b`$$`>`$,$`<`$log $`Z`$$`>`$) are primary indicators of the best ($`t`$,log $`b`$,log $`Z`$) solution, and their standard deviations ($`\sigma _t`$,$`\sigma _{\mathrm{log}b}`$,$`\sigma _{\mathrm{log}Z}`$) could be taken as estimators of the deviation of the data. However, due to the well known age-metallicity and age-burst strength degeneracies these standard deviations are not representative of the distribution of these solutions in the ($`t`$,log $`b`$,log $`Z`$) space. Fortunately, these degenaracies do not span the whole range in age, burst strength and metallicity given by the models, being relatively well constrained by the ($`rJ`$,$`JK`$,2.5$`\times `$log EW(H$`\alpha `$)) data.
Therefore, we have studied the clustering of the ($`t`$,log $`b`$,log $`Z`$) solutions for each individual galaxy. We have used for this analysis a single linkage hierarchical clustering method \[Murtagh & Heck 1987\]. First, (1) we determine the distances between every couple of solutions, which represents a total of $`N\times (N1)/2`$ dissimilarities (=distances), being N the number of solutions. The dissimilarity between the elements $`j`$ and $`k`$, $`d_{j,k}`$, is defined as
$$d_{j,k}^2=\underset{i=1}{\overset{n}{}}(x_{ij}x_{ik})^2$$
(5)
The matrix of dissimilarities is known as dendogram. Then, (2) we find the smallest dissimilarity, $`d_{ik}`$. These points, $`i`$ and $`k`$, (3) are therefore agglomerated and replaced with a new point, $`ik`$, and the dissimilarities updated such that, for all objects $`ji,k`$,
$$d_{ik,j}=min\{d_{i,j},d_{k,j}\}$$
(6)
Then, (4) the dissimilarities $`d_{i,j}`$ and $`d_{k,j}`$, for all $`j`$, are deleted, as these are no longer used. Finally, we return to step (1) after reducing the dimension of the dissimilarities matrix and the number of clusters. Finally, we recover the last three clusters of solutions.
The clustering pattern obtained is basically produced by the discretization in metallicity of the original BC96 evolutionary synthesis models. Now, we analyze the solutions within each solution cluster. In this case, the discretizations in burst strength and age are comparable and a Principal Component Analysis (PCA hereafter) is the most suitable choice (0.04 dex in burst strength and $``$0.05 dex in age).
The PCA basically determines, in a R<sup>n</sup> data array, the set of $`n`$ orthogonal axes that better reproduces our data distribution. The first new axis, i.e. the principal component, will try to go as close as possible through all the data points, describing the larger fraction of the data variance. Figure 8 shows the principal component for each of the three solution clusters of a hypothetical ($`t`$,log $`b`$,log $`Z`$) distribution.
Formally, following the PCA (see, e.g., Morrison 1976), (1) we construct the variance-covariance and the correlation matrix of the sample, being the $`(j,k)^{th}`$ term of these matrixes, respectively,
$`c_{jk}={\displaystyle \frac{1}{n}}{\displaystyle \underset{i=1}{\overset{n}{}}}(r_{ij}\overline{r}_j)(r_{ik}\overline{r}_k)`$ (7)
$`\rho _{jk}={\displaystyle \frac{1}{n}}{\displaystyle \underset{i=1}{\overset{n}{}}}{\displaystyle \frac{(r_{ij}\overline{r}_j)(r_{ik}\overline{r}_k)}{s_js_k}}`$ (8)
where
$$s_j^2=\frac{1}{n}\underset{i=1}{\overset{n}{}}(r_{ij}\overline{r}_j)^2.$$
(9)
Then, (2) solving the eigenvector equation, $`\rho u`$=$`\lambda u`$, we obtain the eigenvalues and eigenvectors of the correlation matrix. The ratio between an eigenvalue and the sum of all the eigenvalues, $`\lambda _i/_{i=1}^n\lambda _i`$, gives us the contribution of the new axis, determined by the corresponding eigenvector, to the total data variance. Therefore, the eigenvector with higher eigenvalue is the principal component and will indicate which is the dominant degeneracy inside each solution cluster.
## ACKNOWLEDGMENTS
This work is based on observations obtained at the Lick Observatory, operated by the University of California and on observations collected at the German-Spanish Astronomical Centre, Calar Alto, Spain, operated by the Max-Planck Institute fur Astronomie (MPIE), Heidelberg, jointly with the Spanish Commission for Astronomy. It is also partly based on observations made with the Isaac Newton Telescope operated on the island of La Palma by the Royal Greenwich Observatory in the Spanish Observatorio del Roque de los Muchachos of the Instituto de Astrofísica de Canarias.
A. Gil de Paz thanks the Institute of Astronomy of the University of Cambridge for all the facilities and support during his stay there. J. Gallego and A. Gil de Paz acknowledge the invitation, hospitality and facilities provided during the 3rd Guillermo Haro Workshop, included in the Guillermo Haro Programme at the INAOE (Mexico). We also thank C.E. García-Dabó, C. Sánchez Contreras and R. Guzmán for stimulating conversations and the referee Dr. M. Edmunds for his useful comments and suggestions. A. Gil de Paz acknowledges the receipt of a Formación del Profesorado Universitario fellowship from the Spanish Ministry of Education. A. Aragón-Salamanca acknowledges generous financial support from the Royal Society. A. Alonso-Herrero was supported by NASA on grant NAG 5-3042. This research was also supported by the Spanish Programa Sectorial de Promoción General del Conocimiento under grants PB96-0610 and PB96-0645. This work was partially carried out under the auspices of EARA, a European Association for Research in Astronomy, and the TMR Network on Galaxy Formation and Evolution funded by the European Commission.
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# Magnetic-field-induced singularities in spin dependent tunneling through InAs quantum dots
## Abstract
Current steps attributed to resonant tunneling through individual InAs quantum dots embedded in a GaAs-AlAs-GaAs tunneling device are investigated experimentally in magnetic fields up to 28 T. The steps evolve into strongly enhanced current peaks in high fields. This can be understood as a field-induced Fermi-edge singularity due to the Coulomb interaction between the tunneling electron on the quantum dot and the partly spin polarized Fermi sea in the Landau quantized three-dimensional emitter.
PACS numbers: 73.40.Gk, 73.23.Hk, 72.10.Fk 85.30.Vw
The interaction of the Fermi sea of a metallic system with a local potential can lead to strong singularities close to the Fermi edge. Such effects have been predicted more than thirty years ago for the X-ray absorption and emission of metals and observed subsequently. Similar singularities as a consequence of many body effects are also known from the luminescence of quantum wells. Matveev and Larkin were the first to predict interaction-induced singularities in the tunneling current via a localized state which were measured experimentally in several resonant tunneling experiments from two-dimensional electrodes through a zero-dimensional system.
Here we report on singularities observed in the resonant tunneling from highly doped three-dimensional (3D) GaAs electrodes through an InAs quantum dot (QD) embedded in an AlAs barrier. These Fermi-edge singularities (FES) show a considerable magnetic field dependence and a strong enhancement in high magnetic fields where the 3D electrons occupy the lowest Landau level in the emitter. We observe an asymmetry in the enhancement for electrons of different spins with an extremely strong FES for electrons carrying the majority spin of the emitter. The experimental observations are explained by a theoretical model taking into account the electrostatic potential experienced by the conduction electrons in the emitter due to the charged QD. We will show that the partial spin polarisation of the emitter causes extreme values of the edge exponent $`\gamma `$ not observed until present and going beyond the standard theory valid for $`\gamma 1`$ .
The active part of our samples are self-organized InAs QDs with 3-4 nm height and 10-15 nm diameter embedded in the middle of a 10 nm-thick AlAs barrier and
sandwiched between two 3D electrodes. They consist of a 15 nm undoped GaAs spacer layer and a GaAs-buffer with graded doping. A typical InAs dot is sketched in inset (a) of Fig. 1, the vertical band structure across a dot is schematically shown in inset (b).
Current voltage ($`I`$-$`V`$) characteristics were measured in large area vertical diodes ($`40\times 40\mu `$m<sup>2</sup>) patterned on the wafer. In Fig. 1 we show a part of a typical $`I`$-$`V`$-curve with several discrete steps. We have demonstrated previously that such steps can be assigned to single electron tunneling from 3D electrodes through individual InAs QDs consistent with other resonant tunneling experiments through self-organized InAs QDs .
For the positive bias voltages shown in Fig. 1 the electrons tunnel from the bottom electrode into the base of an InAs QD and leave the dot via the top. The tunneling current is mainly determined by the tunneling rate through the effectively thicker barrier below the dot (single electron tunneling regime). A step in the current occurs at bias voltages where the energy level of a dot, $`E_D`$, coincides with the Fermi level of the emitter, $`E_F`$.
In the following we will concentrate on the step labeled (\*) in Fig. 1. Other steps in the same structure as well as steps observed in the $`I`$-$`V`$-characteristics of other structures show a very similar behavior.
After the step edge a slight overshoot in the tunneling current occurs consistent with other tunneling experiments through a localized impurity or through InAs dots . This effect is caused by the Coulomb interaction between a localized electron on the dot and the electrons at the Fermi edge of the emitter. The decrease of the current $`I(V)`$ towards higher voltages $`V>V_0`$ follows a power law $`I(VV_0)^\gamma `$ ($`V_0`$ is the voltage at the step edge) with an edge exponent $`\gamma =0.02\pm 0.01`$.
The evolution of step (\*) in a magnetic field applied parallel to the current direction is shown in Fig. 2a. The step develops into two separate peaks with onset voltages marked as $`V_{}`$ and $`V_{}`$. The Landau quantization of the emitter leads to an oscillation of $`V_{}`$ and $`V_{}`$ and a shift to smaller voltages as a function of magnetic field, see Fig. 2b. This reflects the magneto-quantum-oscillation of the Fermi energy in the emitter . From the period and the amplitude of the oscillation we can extract a Fermi energy (at $`B=0`$) $`E_0=13.6`$meV and a Landau level broadening $`\mathrm{\Gamma }=1.3`$meV in the 3D emitter. The measured $`E_0=13.6`$meV agrees well with the expected electron concentration at the barrier derived from the doping profile in the electrodes.
For $`B>6`$ T only the lowest Landau level remains occupied. With a level broadening $`\mathrm{\Gamma }=1.3`$meV the Fermi level $`E_F`$ for 15 T $`<`$ B $`<`$ 30 T is within less than $`2`$meV pinned to the bottom of the lowest Landau band, $`E_L=\mathrm{}\omega _c/2`$. As a consequence the onset voltage shifts to lower values as $`\alpha e\mathrm{\Delta }V\mathrm{}\omega _c/2`$ with $`\alpha =0.34`$. The diamagnetic shift of the energy level in the dot can be neglected compared to this shift of the Fermi energy in the emitter. For the dot investigated in with $`r_0=3.7`$ nm the diamagnetic shift at 30 T is $`\mathrm{\Delta }E_D=3.5`$ meV negligible compared to $`E_L=26`$meV.
The two distinct steps with onset voltages $`V_{}`$ and $`V_{}`$ originate from the spin-splitting of the energy level $`E_D`$ in the dot. Their distance $`\mathrm{\Delta }V_p`$ is given by the Zeeman splitting $`\mathrm{\Delta }E_z=g_D\mu _BB=\alpha e\mathrm{\Delta }V_p`$ with an energy-to-voltage conversion factor $`\alpha =0.34`$ . As shown in Fig. 2c $`\mathrm{\Delta }V_p`$ is indeed linear in B, with a Landé factor $`g_D=0.8`$ in agreement with other experiments on InAs dots .
For low magnetic fields ($`B9`$T in our case, see graph for $`B=9T`$ in Fig. 2a) the size of the steps is very similar for both spins and about half of the size at zero field. Also the slight overshoot in the current as the signature of a Fermi edge singularity is similar for both spin orientations and comparable to the zero field case with an edge exponent $`\gamma <0.05`$ for all magnetic fields $`B<10`$T.
The form of the current steps changes drastically in high magnetic fields where only the lowest Landau level of the emitter remains occupied. In particular, the second current step at higher voltage evolves into a strongly enhanced peak with a peak current of one order of magnitude higher compared to the zero-field case.
Following we assume that $`g_D`$ is positive whereas the Landé factor in bulk GaAs is negative. This assumption is verified by the fact that the energetically lower lying state (first peak in Fig. 3) is thermally occupied at higher temperatures and can therefore be identified with the minority spin in the emitter.
The strongly enhanced current peak at higher energies is due to tunneling through the spin state corresponding to the majority spin (spin up) in the emitter. The resulting spin configuration is scetched in the inset of Fig. 2a and will also be confirmed below by our theoretical results.
The shape of this current peak can be described by a steep ascent and a more moderate decrease of the current towards higher voltages. Down to temperatures $`T<100`$ mK the steepness of the ascent is only limited by thermal broadening. The decrease of the current for $`V>V_0`$ is again described with the characteristic behavior for a Fermi-edge singularity, $`I(VV_0)^\gamma `$, where $`V_0`$ here is the voltage at the maximum peak current.
However, along with the drastic increase of the peak current the edge exponent $`\gamma `$ increases dramatically reaching a value $`\gamma >0.5`$ for the highest fields.
A different way to visualize the signature of a FES is a temperature dependent experiment. As an example we have plotted the $`I`$-$`V`$-curve at $`B=22`$ T for different temperatures in Fig. 3. As shown in the inset the peak maximum $`I_0`$ for the spin-up electrons decreases according to a power law $`I_0T^\gamma `$ with an edge exponent $`\gamma =0.43\pm 0.05`$. Such a strong temperature dependence is characteristic for a FES and allows us to exclude that pure density of states effects in the 3D emitter are responsible for the current peaks in high magnetic fields.
As shown in Fig. 3 an edge exponent $`\gamma =0.43`$ also fits within experimental accuracy the observed decrease of the current for $`V>V_0`$.
It is not possible to extract the edge exponent for the minority spin directly from temperature dependent experiments. At high magnetic fields the observed increase of the current with increasing temperature is mainly caused by an additional thermal population of the minority spin in the emitter. The general form of the curve is merely affected by temperature. Therefore, the edge exponent can only be gained from fitting the shape of the current peaks.
A compilation of the edge exponents $`\gamma `$ for various magnetic fields and both spin orientations is shown in Fig. 4.
For the data related to the majority spin two independent methods were used to extract $`\gamma `$. For the minority spin only fitting of the shape of the $`I`$-$`V`$-curves was used.
For a theoretical description of these effects we consider a 3D electron gas in the half space $`z<0`$. In a sufficiently strong magnetic field $`B||\widehat{z}`$ all electrons are in the lowest Landau level. This defines a set of one-dimensional channels with momentum $`\mathrm{}k`$ perpendicular to the boundary. This situation is different from the cases considered for scattering off point defects as in Refs. or for a 2D electron gas where the current is carried by edge states . The single particle wave functions in channel $`m0`$ are $`\psi _m(\rho ,\varphi )\mathrm{sin}kz`$ with $`\psi _m(\rho ,\varphi )\rho ^m\mathrm{exp}(im\varphi \rho ^2/4\mathrm{}_0^2)`$. In the experiments the magnetic length $`\mathrm{}_0=\sqrt{\mathrm{}/eB}`$ ($`\mathrm{}_0=5.6`$ nm at 20 T) is comparable to the lateral size of the QD $`2r_07`$ nm. Hence the effect of the electrostatic potential of a charged dot on the electrons in a given channel of the emitter decreases rapidly with $`m`$, and the observed FES are mainly due to tunneling of electrons from the $`m=0`$ channel into the dot. Following tunneling processes of spin $`\sigma `$ electrons from the $`m=0`$ state in the emitter give rise to a FES with edge exponent
$$\gamma _\sigma =\frac{2}{\pi }\delta _0(k_{F\sigma })\frac{1}{\pi ^2}\underset{m}{}\underset{\tau =,}{}\left(\delta _m(k_{F\tau })\right)^2$$
(1)
where $`\delta _m(k)`$ is the Fermi phase shift experienced by the electrons in the $`m`$-th channel due to the potential of the quantum dot . From (1) the observed field dependence of the edge exponents is a consequence of the variation of the Fermi momenta for spin-$`\sigma `$ electrons with magnetic field *and* the field dependence of the effective potential in the one-dimensional channels. The former can be computed from the one-dimensional density of states (DOS) of the lowest Landau band
$$D(E,B)=\frac{e\sqrt{m^{}}}{(2\pi \mathrm{})^2}B\left(d(ϵ_{})+d(ϵ_{})\right).$$
(2)
Here $`ϵ_\sigma =E(\mathrm{}\omega _c\pm g^{}\mu _BB)/2`$ is the energy of electrons with spin-$`\sigma `$ measured from the bottom of the Landau band. $`g^{}0.33`$ is the effective Landé factor of the electrons in the emitter. The DOS for the spin-subbands is $`d(ϵ)=\sqrt{2}\mathrm{Re}(ϵ+i\mathrm{\Gamma })^{1/2}`$. Without broadening, $`\mathrm{\Gamma }=0`$, one has $`k_{F\sigma }=\pi ^2n\mathrm{}_0^2(1\pm b^3)`$ where $`n`$ is the 3D density of electrons and $`b`$ is the magnetic field measured in units of the field necessary for complete spin polarization of the 3D emitter. Using a Fermi energy $`E_0=13.6`$ meV and neglecting level broadening we find that only the lowest Landau level (*both spin states!*) is occupied for $`B_1>5.2`$ T. Including level broadening changes $`B_1`$ to a slightly higher value. With the known field dependence of the Fermi energy in the quantum limit we can calculate the field for total spin polarisation
$$B_{pol}=\left(\frac{16}{9\xi }\right)^{1/3}\frac{m^{}E_0}{\mathrm{}e}43\text{T}$$
(3)
with $`g^{}=0.33`$ and $`m^{}=0.067m_0`$. $`\xi =\frac{1}{2}|g^{}|m^{}/m_0`$ is the ratio between spin splitting and Landau level splitting.
To make contact to the experimental observations we have to specify the interaction of the screened charge on the QD and the conduction band electrons. A Thomas-Fermi calculation gives $`U(\rho ,z)=(2e^2\mathrm{exp}(\kappa z)/\kappa )(d/(\rho ^2+d^2)^{(3/2)})`$ . Here $`d=5\mathrm{nm}`$ is the width of the insulating layer and $`\kappa ^1=7\mathrm{nm}`$ is the Debye radius. The effective potential seen by electrons in channel $`m`$ is $`V_m\mathrm{exp}(\kappa z)/\kappa `$ with $`V_m=2e^2d𝑑\rho ^2|\psi _m(\rho ,\varphi )|^2/(\rho ^2+d^2)^{(3/2)}`$. For large $`\kappa `$ we obtain for the phase shift in the $`m=0`$ channel $`\delta _0(k)v_0f(B)k/\kappa `$ where
$$f(B)=\left(\frac{d}{\mathrm{}_0}\right)^2\left\{1\sqrt{\frac{\pi }{2}}\frac{d}{\mathrm{}_0}\mathrm{e}^{\frac{d^2}{2\mathrm{}_0^2}}\mathrm{erfc}\left(\frac{d}{\sqrt{2}\mathrm{}_0}\right)\right\}$$
(4)
and $`v_0(m^{}e^2/\mathrm{}^2\kappa )(\kappa d)^2`$ up to a numerical factor. Similarly we obtain the integrated effect of the channels $`m>0`$ in (1). In Fig. 4 the resulting exponents $`\gamma _\sigma `$ obtained for $`\sigma =,`$ are shown for $`v_0=6.75`$ and a broadening $`\mathrm{\Gamma }=0`$ and $`\mathrm{\Gamma }=1.3`$ meV, respectively. The value used for $`\mathrm{\Gamma }`$ reflects its realistic experimental value. $`v_0`$ is the only fit parameter.
Already the simple model with no level broadening ($`\mathrm{\Gamma }=0`$) is in good agreement with the experimentally measured edge exponents for both spin directions, especially in high magnetic fields where possible admixtures of higher Landau levels play a minor rule. Including level broadening leads to a less dramatic spin polarisation in the emitter and as a consequence smears out the field dependence of $`\gamma `$ for the minority spin. The basic features, however, remain unchanged. In particular, the edge exponent for the minority spin retains moderate values for high magnetic fields, whereas the edge exponent related to the majority spin shows a strong field dependence with very high values in high magnetic fields.
In conclusion we have evaluated experimental data concerning magnetic-field-induced FES in resonant tunneling experiments through InAs QDs. We have shown that the interaction between a localized charge and the electrons in the Landau quantized emitter leads to dramatic Fermi phase shifts if only the lowest Landau level in the 3D emitter is occupied. This results in edge exponents $`\gamma >0.5`$ which were measured and described theoretically.
We would like to thank H. Marx for sample growing, P. König for experimental support and F. J. Ahlers for valuable discussions. Part of this work has been supported by the TMR Programme of the European Union under contract no. ERBFMGECT950077. We acknowledge partial support from the Deutsche Forschungsgemeinschaft under Grants HA 1826/5-1 and Fr 737/3.
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# Is the Polarized Antiquark Sea in the Nucleon Flavor Symmetric?
## Abstract
We show that the model which naturally explains the $`\overline{u}\overline{d}`$ asymmetry in the nucleon and is in quantitative agreement with the Gottfried sum rule data, also predicts that in the proton $`\mathrm{\Delta }\overline{u}>0>\mathrm{\Delta }\overline{s}>\mathrm{\Delta }\overline{d}`$ and $`\mathrm{\Delta }\overline{u}\mathrm{\Delta }\overline{d}>\overline{d}\overline{u}>0`$. At the input scale, these results can be derived even analytically. Thus the violation of the flavor symmetry is more serious in the polarized case than in the unpolarized case. In contrast, many recent analyses of the polarized data have made a simplifying assumption that all the three $`\mathrm{\Delta }\overline{q}`$’s have the same sign and magnitude. We point out the need to redo these analyses, allowing for the alternate scenario as described above. We present predictions of the model for the $`W^{}`$ asymmetry in polarized $`pp`$ scattering, which can be tested at RHIC; these are quite different from those available in the literature.
PACS numbers: 14.20.Dh, 13.60.Hb, 13.88.+e
Keywords: polarized deep inelastic scattering, polarized pp collisions, polarized parton densities, antiquark flavor asymmetry, spin asymmetries, statistical model of the nucleon
Several comprehensive analyses of the polarized deep inelastic scattering (DIS) data, based on next-to-leading-order quantum chromodynamics (QCD), have appeared recently . In these analyses the polarized parton density functions (PDFs) are either written in terms of the well-known parameterizations of the unpolarized PDFs or parameterized independently, and the unknown parameters are determined by fitting the polarized DIS data. Additional simplifying assumptions are often made; the one that has been widely used in the literature is
$$\mathrm{\Delta }\overline{u}=\mathrm{\Delta }\overline{d}=\lambda \mathrm{\Delta }\overline{s},$$
(1)
with a positive $`\lambda `$ which is usually set equal to unity. Recently the HERMES and SMC collaborations too analyzed their inclusive and semi-inclusive DIS data assuming all $`\mathrm{\Delta }\overline{q}`$’s to be of the same sign. The same assumption has also been used to make predictions for future accelerators; see e.g. . In some analyses , the nonsinglet PDFs $`\mathrm{\Delta }q_3`$ and $`\mathrm{\Delta }q_8`$ are assumed to differ only by a constant multiplicative factor (see (21) below).
In this paper, we examine these simplifying assumptions made in the literature. This is important because a similar ad hoc assumption about the flavor decomposition of the unpolarized antiquark sea, $`\overline{u}=\overline{d}`$, turned out to be wrong when accurate data on muon DIS became available , and the global analyses of the unpolarized data had to be redone. Here we derive a series of inequalities satisfied by the PDFs and point out the need to redo the global analyses of the polarized DIS data in the light of these inequalities, allowing in particular for the violation of the flavor symmetry in the polarized antiquark sea; see (20) below. We then present predictions of our model, which can be tested in polarized $`pp`$ scattering at RHIC, BNL. Finally, we describe other recent works on flavor asymmetry of polarized sea distributions.
We use the framework of the statistical model for polarized and unpolarized structure functions and PDFs of the proton and the neutron, which was presented recently . This model provided a natural explanation of the $`\overline{u}\overline{d}`$ asymmetry in the nucleon and was in quantitative agreement with the Gottfried sum rule data. Additionally, it reproduced the data on $`F_2^p(x,Q^2)`$ for $`0.00001<x<1`$ and $`2.5<Q^2<5000`$ GeV<sup>2</sup>, $`F_2^p(x)F_2^n(x),F_2^n(x)/F_2^p(x),xg(x),\overline{d}(x)\overline{u}(x),d(x)/u(x)`$, the fractional momentum of charged partons and the polarized structure functions $`g_1^{p,n}(x)`$, at various $`Q^2`$. Out of those, only the $`F_2^p`$ and $`(F_2^pF_2^n)`$ data, both at $`Q^2=4`$ GeV<sup>2</sup> only, were used as an input to fix the model parameters, and all other results served as model predictions. In particular, the $`d(x)/u(x)`$ ratio in the limit $`x1`$ turned out to be 0.22 in good agreement with the QCD prediction 0.2 . At the input scale ($`Q^2=Q_0^2=M^2`$, where $`M`$ is the nucleon mass), all $`xq(x)`$ and $`x\overline{q}(x)`$ distributions were found to be valence-like, and $`xg(x)`$ was found to be constant in the limit $`x0`$. Thus the total number of gluons was logarithmically divergent providing a strong a posteriori justification for the statistical model ansatz . Contrary to the common practice, the polarized and unpolarized data were reproduced in a single framework and the simplifying assumption of charge symmetry was not made. Here we further explore the predictive power of the model.
If $`n_{\alpha (\overline{\alpha })()}`$ denotes the number of quarks (antiquarks) of flavor $`\alpha `$ and spin parallel (antiparallel) to the proton spin, then any model of PDFs in the proton has to satisfy the following constraints:
$`n_u+n_un_{\overline{u}}n_{\overline{u}}`$ $`=`$ $`2,`$ (2)
$`n_d+n_dn_{\overline{d}}n_{\overline{d}}`$ $`=`$ $`1,`$ (3)
$`n_s+n_sn_{\overline{s}}n_{\overline{s}}`$ $`=`$ $`0,`$ (4)
$`n_un_u+n_{\overline{u}}n_{\overline{u}}`$ $`=`$ $`\mathrm{\Delta }u+\mathrm{\Delta }\overline{u},`$ (5)
$`n_dn_d+n_{\overline{d}}n_{\overline{d}}`$ $`=`$ $`\mathrm{\Delta }d+\mathrm{\Delta }\overline{d},`$ (6)
$`n_sn_s+n_{\overline{s}}n_{\overline{s}}`$ $`=`$ $`\mathrm{\Delta }s+\mathrm{\Delta }\overline{s}.`$ (7)
The RHSs of (5)-(7) have been measured by several groups. We use $`(\mathrm{\Delta }u+\mathrm{\Delta }\overline{u})=0.83\pm 0.03,(\mathrm{\Delta }d+\mathrm{\Delta }\overline{d})=0.43\pm 0.03,(\mathrm{\Delta }s+\mathrm{\Delta }\overline{s})=0.10\pm 0.03`$; see . The parton numbers $`n_{\alpha (\overline{\alpha })()}`$ in (2)-(7) are obtained by integrating the appropriate number density $`dn/dx`$ over $`x`$. The various $`\mathrm{\Delta }`$’s are also $`x`$-integrated quantities.
The RHSs of (2)-(4) are clearly $`Q^2`$-independent. The RHSs of (5)-(7) are also $`Q^2`$-independent in the jet and Adler-Bardeen (AB) schemes: Recall that the nonsinglets $`\mathrm{\Delta }q_3=(\mathrm{\Delta }u+\mathrm{\Delta }\overline{u})(\mathrm{\Delta }d+\mathrm{\Delta }\overline{d})`$ and $`\mathrm{\Delta }q_8=(\mathrm{\Delta }u+\mathrm{\Delta }\overline{u})+(\mathrm{\Delta }d+\mathrm{\Delta }\overline{d})2(\mathrm{\Delta }s+\mathrm{\Delta }\overline{s})`$ are $`Q^2`$-independent in all renormalization schemes because of the conservation of the nonsinglet axial vector current, and the singlet $`\mathrm{\Delta }\mathrm{\Sigma }=(\mathrm{\Delta }u+\mathrm{\Delta }\overline{u})+(\mathrm{\Delta }d+\mathrm{\Delta }\overline{d})+(\mathrm{\Delta }s+\mathrm{\Delta }\overline{s})`$ is $`Q^2`$-independent in the jet and AB schemes because of the Adler-Bardeen theorem . As a result, $`(\mathrm{\Delta }u+\mathrm{\Delta }\overline{u}),(\mathrm{\Delta }d+\mathrm{\Delta }\overline{d})`$ and $`(\mathrm{\Delta }s+\mathrm{\Delta }\overline{s})`$ which can be expressed as linear combinations of $`\mathrm{\Delta }q_3,\mathrm{\Delta }q_8`$ and $`\mathrm{\Delta }\mathrm{\Sigma }`$, are also $`Q^2`$-independent in these two schemes. In the $`\overline{MS}`$ scheme, on the other hand, $`\mathrm{\Delta }\mathrm{\Sigma }`$ is $`Q^2`$-independent at the leading order and only weakly $`Q^2`$-dependent at the next-to-leading order. Empirically too $`\mathrm{\Delta }\mathrm{\Sigma }`$ is found to be almost $`Q^2`$-independent; see e.g. Fig. 5 of . Hence in the $`\overline{MS}`$ scheme the RHSs of (5)-(7) are expected to be nearly $`Q^2`$-independent.
We now show how the statistical model naturally leads to a violation of the flavor symmetry in the unpolarized and polarized seas in the nucleon. Consider the following 6 equations:
$`2n_u2n_{\overline{u}}`$ $`=`$ $`2.83,`$ ()
$`2n_u2n_{\overline{u}}`$ $`=`$ $`1.17,`$ ()
$`2n_d2n_{\overline{d}}`$ $`=`$ $`0.57,`$ ()
$`2n_d2n_{\overline{d}}`$ $`=`$ $`1.43,`$ ()
$`2n_s2n_{\overline{s}}`$ $`=`$ $`0.10,`$ ()
$`2n_s2n_{\overline{s}}`$ $`=`$ $`0.10.`$ ()
These are obtained from (2)-(7) by linearly combining the latter set of equations in pairs. E.g. (8) and (9) are obtained by adding or subtracting (2) and (5).
It was shown in that the parton number density $`dn/dx`$ in the infinite-momentum frame, at the input scale, is given by
$$\frac{dn}{dx}=\frac{M^2x}{2}_{xM/2}^{M/2}\frac{dE}{E^2}\frac{dn}{dE},$$
(8)
where
$$dn/dE=gf(E)(VE^2/2\pi ^2+aR^2E+bR),$$
(9)
is the density in the nucleon rest frame. Here $`M`$ is the nucleon mass, $`E`$ is the parton energy in the nucleon rest frame, $`g`$ is the spin-color degeneracy factor, $`f(E)`$ is the usual Fermi or Bose distribution function $`f(E)=\{\mathrm{exp}[(E\mu )/T]\pm 1\}^1`$, $`V`$ is the nucleon volume and $`R`$ is the radius of a sphere with volume $`V`$. The three terms in (14b) are the volume, surface and curvature terms, respectively; in the thermodynamic limit only the first survives. The two free parameters $`a`$ and $`b`$ in (14b) were determined in by fitting the structure function $`F_2(x,Q^2)`$ data at $`Q^2=4`$ GeV<sup>2</sup>. Their values as well as the values of the temperature $`(T)`$ and chemical potential $`(\mu )`$ which get determined due to (2)-(7), were given in .
At the input scale, with the help of (14), (8) can be written in a full form as
$$_0^1𝑑x\frac{M^2x}{2}_{xM/2}^{M/2}\frac{dE}{E^2}g(VE^2/2\pi ^2+aR^2E+bR)\left[\frac{2}{e^{\beta (E\mu _u)}+1}\frac{2}{e^{\beta (E\mu _{\overline{u}})}+1}\right]=2.83.$$
(10)
It is straightforward to show that the chemical potentials for quarks and antiquarks satisfy the relations
$`\mu _{\overline{q}}`$ $`=`$ $`\mu _q,`$ ()
$`\mu _{\overline{q}}`$ $`=`$ $`\mu _q.`$ ()
So it follows from (15) and (16b) that $`\mu _u>0`$. Similar arguments show that $`\mu _u,\mu _d,\mu _d`$ and $`\mu _s`$ are positive and $`\mu _s`$ is negative. Moreover, since the RHSs of (12) and (13) differ only in sign, we have $`\mu _s=\mu _s`$. Since RHSs of (8)-(13) can be arranged as $`2.83>1.43>1.17>0.57>0.10>0.10`$, the corresponding chemical potentials satisfy
$$\mu _u>\mu _d>\mu _u>\mu _d>(\mu _s=\mu _{\overline{s}})>0>(\mu _s=\mu _{\overline{s}})>\mu _{\overline{d}}>\mu _{\overline{u}}>\mu _{\overline{d}}>\mu _{\overline{u}}.$$
(11)
It will be useful to recall the actual values of the $`\mu `$’s given in . They are (in MeV) $`\mu _u=210,\mu _d=106,\mu _u=86,\mu _d=42,\mu _s=7,\mu _s=7.`$ $`\mu `$’s for the antiquarks follow from (16). \[The RHSs of (8)-(13) are sufficiently different from each other so that the experimental errors in $`(\mathrm{\Delta }q+\mathrm{\Delta }\overline{q})`$, quoted above, will not alter the ordering in (17).\] (17) together with (14) yields, at the input scale $`Q_0^2(=M^2=0.88`$ GeV<sup>2</sup>):
$$n_u>n_d>n_u>n_d>(n_s=n_{\overline{s}})>(n_s=n_{\overline{s}})>n_{\overline{d}}>n_{\overline{u}}>n_{\overline{d}}>n_{\overline{u}}>0.$$
(12)
As a check, it is easy to verify that (18) reproduces the correct signs of the RHSs of (2)-(7). Notice the symmetric arrangement of the $`\mu `$’s in (17) and the consequent arrangement of the $`n`$’s in (18).
To recapitulate, the statistical model provides a quantitative method to incorporate the effects of the Pauli exclusion principle into the PDFs: the RHSs of the number constraints (2)-(7) or equivalently (8)-(13), force the various chemical potentials and hence the parton distributions to be arranged as in (17) and (18), respectively, at the input scale.
Further consequences of (18) are easy to derive: (Note $`n_q=n_q+n_q`$ and $`\mathrm{\Delta }q=n_qn_q`$.)
(a) The general positivity constraints on the polarized and unpolarized PDFs: $`|\mathrm{\Delta }q|n_q`$ are satisfied trivially.
(b) $`\mathrm{\Delta }u>0,\mathrm{\Delta }d<0,\mathrm{\Delta }s<0`$.
(c) $`\mathrm{\Delta }\overline{u}>0,\mathrm{\Delta }\overline{d}<0,\mathrm{\Delta }\overline{s}<0`$. This is in contrast to the assumption (1) made in the literature that all the three $`\mathrm{\Delta }\overline{q}`$’s have the same sign.
(d) $`\mathrm{\Delta }u_v=\mathrm{\Delta }u\mathrm{\Delta }\overline{u}=n_un_un_{\overline{u}}+n_{\overline{u}}>0`$, because the two $`n_{\overline{u}}`$ terms are too small compared to the two $`n_u`$ terms (see (18)) to change the sign of the RHS.
(e) $`\mathrm{\Delta }d_v=\mathrm{\Delta }d\mathrm{\Delta }\overline{d}=n_dn_dn_{\overline{d}}+n_{\overline{d}}<0`$, because the two $`n_{\overline{d}}`$ terms are too small compared to the two $`n_d`$ terms (see (18)) to change the sign of the RHS.
(f) $`\mathrm{\Delta }s_v=\mathrm{\Delta }s\mathrm{\Delta }\overline{s}=0`$.
(g) $`\mathrm{\Delta }q_3=(\mathrm{\Delta }u+\mathrm{\Delta }\overline{u})(\mathrm{\Delta }d+\mathrm{\Delta }\overline{d})>0`$; see (b)-(c).
(h) $`n_{\overline{d}}>n_{\overline{u}}`$ which leads to the Gottfried sum rule violation. Thus the statistical model naturally leads to the $`\overline{u}\overline{d}`$ asymmetry in the unpolarized sea . Moreover, it was shown in that the model is in quantitative agreement with the data on $`(F_2^pF_2^n)`$ vs $`x`$ and the Gottfried sum $`S_G`$.
(i) $`\mathrm{\Delta }\overline{u}\mathrm{\Delta }\overline{d}>n_{\overline{d}}n_{\overline{u}}>0`$. Thus the violation of the flavor symmetry is more serious in the polarized case than in the unpolarized case.
(j) $`\mathrm{\Delta }d\mathrm{\Delta }s=n_dn_dn_s+n_s<0`$, because $`n_s`$ and $`n_s`$ tend to cancel each other, unlike $`n_d`$ and $`n_d`$. Combining this result with (b) above, one gets $`|\mathrm{\Delta }d|>|\mathrm{\Delta }s|`$, and
$$\mathrm{\Delta }u>0>\mathrm{\Delta }s>\mathrm{\Delta }d.$$
(13)
(k) $`\mathrm{\Delta }\overline{d}\mathrm{\Delta }\overline{s}=n_{\overline{d}}n_{\overline{d}}n_{\overline{s}}+n_{\overline{s}}<0`$, because $`n_{\overline{s}}`$ and $`n_{\overline{s}}`$ tend to cancel each other, unlike $`n_{\overline{d}}`$ and $`n_{\overline{d}}`$. Combining this result with (c) above, one gets $`|\mathrm{\Delta }\overline{d}|>|\mathrm{\Delta }\overline{s}|`$, and
$$\mathrm{\Delta }\overline{u}>0>\mathrm{\Delta }\overline{s}>\mathrm{\Delta }\overline{d}.$$
(14)
We have derived the results (a)-(k) analytically, at the input scale. They are borne out by actual numerical calculations; see Fig. 1 which shows our polarized PDFs at the input scale $`Q_0^2=M^2=0.88`$ GeV<sup>2</sup>. We have evolved our polarized PDFs in the next-to-leading-order QCD, in the $`\overline{MS}`$ scheme, in the range $`Q_0^2<Q^2<6500`$ GeV<sup>2</sup>. We find that the results (a)-(k) are valid throughout this range. Figure 2 shows that the violation of the flavor symmetry is more serious in the polarized case than in the unpolarized case, throughout this range.
Incidentally, we have examined another simplifying assumption made e.g. in , namely
$$\mathrm{\Delta }q_3(x,Q^2)=C\mathrm{\Delta }q_8(x,Q^2),$$
(15)
where C is a constant independent of $`x`$ and $`Q^2`$. The present model predicts that (21) is not justified (Fig. 1).
The statistical model makes concrete predictions for various asymmetries in polarized $`pp`$ scattering, which can be tested at RHIC. For example, parity-violating single- and double-spin asymmetries for $`W`$ production in the reactions $`\stackrel{}{p}pW^\pm X`$ and $`\stackrel{}{p}\stackrel{}{p}W^\pm X`$ respectively, are given by
$`A_L^{PV}(W^+)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }u(x_a,M_W^2)\overline{d}(x_b,M_W^2)\mathrm{\Delta }\overline{d}(x_a,M_W^2)u(x_b,M_W^2)}{u(x_a,M_W^2)\overline{d}(x_b,M_W^2)+\overline{d}(x_a,M_W^2)u(x_b,M_W^2)}},`$ ()
$`A_L^{PV}(W^{})`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }\overline{u}d+\mathrm{\Delta }d\overline{u}}{\overline{u}d+d\overline{u}}},`$ ()
$`A_{LL}^{PV}(W^+)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }u\overline{d}u\mathrm{\Delta }\overline{d}\mathrm{\Delta }\overline{d}u+\overline{d}\mathrm{\Delta }u}{u\overline{d}\mathrm{\Delta }u\mathrm{\Delta }\overline{d}+\overline{d}u\mathrm{\Delta }\overline{d}\mathrm{\Delta }u}},`$ ()
$`A_{LL}^{PV}(W^{})`$ $`=`$ $`{\displaystyle \frac{\overline{u}\mathrm{\Delta }d\mathrm{\Delta }\overline{u}dd\mathrm{\Delta }\overline{u}+\mathrm{\Delta }d\overline{u}}{\overline{u}d\mathrm{\Delta }\overline{u}\mathrm{\Delta }d+d\overline{u}\mathrm{\Delta }d\mathrm{\Delta }\overline{u}}},`$ ()
where $`x_a=\sqrt{\tau }e^y,x_b=\sqrt{\tau }e^y,\tau =M_W^2/s`$, $`y`$ is the rapidity of $`W`$ and $`\sqrt{s}`$ is the $`pp`$ center-of-mass energy. The arguments $`x_a,x_b`$ and $`M_W^2`$ are suppressed in (23)-(25) for brevity of notation.
In the present model, $`\mathrm{\Delta }u`$ and $`\mathrm{\Delta }\overline{u}`$ are positive and $`\mathrm{\Delta }d`$ and $`\mathrm{\Delta }\overline{d}`$ are negative (see (b), (c) and Fig. 1). Also note that $`\mathrm{\Delta }uu`$ and $`|\mathrm{\Delta }\overline{d}|\overline{d}`$ (see (a) above). Hence it is straightforward to show that $`0<A_L^{PV}(W^+)<1`$. Similarly, $`1<A_L^{PV}(W^{})<0,A_{LL}^{PV}(W^+)>0,A_{LL}^{PV}(W^{})<0.`$ It is somewhat tedious but again straightforward to show using (22)-(25) that $`A_{LL}^{PV}(W^+)>A_L^{PV}(W^+)`$ and $`|A_{LL}^{PV}(W^{})|>|A_L^{PV}(W^{})|`$. A quick and crude way to convince oneself that $`A_{LL}^{PV}(W^+)>A_L^{PV}(W^+)`$ is to ignore the (small) “$`\mathrm{\Delta }\mathrm{\Delta }`$” terms in the denominator of (24), which makes the denominators of (22) and (24) identical, and then to compare their numerators. In fact, at $`y=0`$ (or $`x_a=x_b`$), $`A_{LL}^{PV}(W^+)`$ is seen to be almost twice as big as $`A_L^{PV}(W^+)`$.
Figure 3 shows our predictions for $`A_L^{PV}`$ and $`A_{LL}^{PV}`$ for $`W^{}`$ production in polarized $`pp`$ scattering at $`\sqrt{s}=500`$ GeV, as a function of the rapidity $`y`$. The above inequalities for $`A_L^{PV}(W^{})`$ and $`A_{LL}^{PV}(W^{})`$, which we derived analytically here are borne out by the actual numerical results in Fig. 3. Also shown for comparison are results reported in . These are based on the parameterizations of polarized PDFs given in . Asymmetries for $`W^{}`$ production are sensitive to the sign of $`\mathrm{\Delta }\overline{u}`$ which is positive in the present model, negative in and $`x`$-dependent in . The recent work of de Florian and Sassot has yielded a clear preference for a positive $`\mathrm{\Delta }\overline{u}`$ distribution.
As stated earlier, the HERMES and SMC collaborations analyzed their inclusive and semi-inclusive DIS data assuming all $`\mathrm{\Delta }\overline{q}`$’s to be of the same sign. Recently, Morii and Yamanishi have reanalyzed these data and have estimated $`\mathrm{\Delta }\overline{d}(x)\mathrm{\Delta }\overline{u}(x)`$ at $`Q^2=4`$ GeV<sup>2</sup>. It is evident from their Fig. 1 that $`\mathrm{\Delta }\overline{u}(x)\mathrm{\Delta }\overline{d}(x)`$ is positive and has a peak at $`x0.06`$ where $`x(\mathrm{\Delta }\overline{u}(x)\mathrm{\Delta }\overline{d}(x))`$ is $`0.05`$. All these observations are consistent with our Fig. 2.
Another model which is able to generate flavor asymmetric polarized antiquark sea is the chiral quark soliton model (CQSM) . Results in our Fig. 2 are strikingly similar to those in . This is remarkable because the physics inputs of the two models are quite different. It is also noteworthy that the origin of the $`\overline{u}\overline{d}`$ and $`\mathrm{\Delta }\overline{u}\mathrm{\Delta }\overline{d}`$ asymmetries is quite simple in the statistical model. While the role of gluons is yet to be understood in CQSM, the statistical model predicts a positive $`\mathrm{\Delta }g(x,Q^2)`$. The pion cloud model also gives rise to the $`\overline{u}\overline{d}`$ asymmetry, \[for a recent review, see \], and there have been some attempts to generate polarization by including spin-1 resonances in that model. These attempts have been commented upon in . Recently, Glück and Reya have discussed the issue of flavor asymmetry, in a phenomenological way making use of the Pauli exclusion principle. We recall that the statistical model provides a quantitative method to incorporate the effects of the Pauli exclusion principle into the PDFs.
We have treated all partons as massless: $`m_u=m_d=m_s=0`$. If $`m_s`$ is taken to be nonzero, then (14) will have to be generalized, but the parton densities still have to satisfy (2)-(7) and equivalently (8)-(13). So it is not obvious how that will affect the symmetric arrangement of the $`\mu `$’s in (17) and the consequent arrangement of the $`n`$’s in (18), at the input scale. This is a nontrivial problem which needs to be investigated further.
In conclusion, we have derived, on rather general grounds, a series of inequalities for the polarized PDFs; see (a)-(k) above. This points to the need to redo the analyses of polarized data, allowing for the alternate scenario as in (19)-(20). Some of the inequalities can be tested in the forthcoming spin-physics program at RHIC, BNL. To illustrate, we have given our predictions for the $`W^{}`$ asymmetries; these are quite different from those available in the literature.
I would like to acknowledge the hospitality of the Nuclear Theory Center, Indiana University where this work was initiated.
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# Spin-triplet superconductivity due to antiferromagnetic spin-fluctuation in Sr2RuO4
## Abstract
A mechanism leading to the spin-triplet superconductivity is proposed based on the antiferromagnetic spin fluctuation. The effects of anisotropy in spin fluctuation on the Cooper pairing and on the direction of $`d`$ vector are examined in the one-band Hubbard model with RPA approximation. The gap equations for the anisotropic case are derived and applied to Sr<sub>2</sub>RuO<sub>4</sub>. It is found that a nesting property of the Fermi surface together with the anisotropy leads to the triplet superconductivity with the $`d=\widehat{z}(\mathrm{sin}k_x\pm i\mathrm{sin}k_y)`$, which is consistent with experiments.
74.20-z, 74.20Mn, 74.25Dw
Since the discovery of superconducting phase in Sr<sub>2</sub>RuO<sub>4</sub>, much effort has been paid for understanding its exotic properties. Among several interesting natures, the most fascinating one is that it is a spin-triplet superconductor confirmed by NMR experiment. While most superconductors found during several decades are singlet, the only exceptions were <sup>3</sup>He and UPt<sub>3</sub>. Therefore the fact that the triplet pairing is realized in Sr<sub>2</sub>RuO<sub>4</sub> has attracted much attention. While UPt<sub>3</sub>, the second example of spin-triplet superconductor, has a complicated electronic structure, Sr<sub>2</sub>RuO<sub>4</sub> has a rather simple electronic state. Thus clarifying the microscopic mechanism of superconductivity in Sr<sub>2</sub>RuO<sub>4</sub> is very important for understanding the triplet superconductors in general.
In <sup>3</sup>He, Cooper pairs are formed due to ferromagnetic spin fluctuations peaked at $`𝒒=\mathrm{𝟎}`$. Therefore it is natural to expect the origin of the triplet pairing in Sr<sub>2</sub>RuO<sub>4</sub> is also ferromagnetic spin fluctuation. This assumption has been believed to be justified by NMR experiments. However the recent neutron scattering experiment has shown that there exists a significant peak near $`𝒒_0=(\pm 2\pi /3,\pm 2\pi /3)`$ and no sizable ferromagnetic spin fluctuation. Thus it is difficult to assume that the spin fluctuation near $`𝒒_0`$ plays no role in the Cooper pairing in Sr<sub>2</sub>RuO<sub>4</sub>. (In the following discussion we call this fluctuation as antiferromagnetic (AF) spin fluctuation, for simplicity.) However this AF fluctuation leads to the singlet superconductivity rather than the triplet superconductivity as expected in analogy to high-$`T_c`$ cuprates.
In this paper we propose a mechanism which gives the triplet pairing even if the spin fluctuation is AF. We find that the characteristic features of Sr<sub>2</sub>RuO<sub>4</sub> are twofold: One is the anisotropy of the spin fluctuation found in NMR experiments, and the other is a nesting property with momentum $`𝒒_0`$ of the two-dimensional Fermi surface. We show that these two features explain the pairing in Sr<sub>2</sub>RuO<sub>4</sub>.
In addition to the competition between singlet and triplet pairing, the direction of the $`𝒅`$ vector, which is the order parameter of triplet superconductivity, is another interesting problem. We show that the anisotropy of the spin fluctuation also explains the experimental fact that the $`𝒅`$ vector is parallel to the $`z`$-direction. First we extend the RPA formulation to the case of anisotropic spin fluctuation. Using the obtained effective interactions, we investigate the most stable pairing based on the weak-coupling gap equations. When the spin fluctuation is isotropic, the so-called d$`_{x^2y^2}`$-wave pairing is the most stable. However when the anisotropy is increased, the state corresponding to $`\widehat{𝒛}(\mathrm{sin}k_x\pm i\mathrm{sin}k_y)`$, which is the prime candidate of Sr<sub>2</sub>RuO<sub>4</sub>, becomes the most stable.
For the $`\gamma `$ band which is one of the three bands in Sr<sub>2</sub>RuO<sub>4</sub>, we assume a two-dimensional effective Hamiltonian
$$H=H_0+\frac{I}{2N}\underset{kk^{}q\sigma }{}c_{k\sigma }^{}c_{k^{}\sigma }^{}c_{k^{}q\sigma }c_{k+q\sigma },$$
(1)
where $`c_{k\sigma }`$ is the annihilation operator of an electron with momentum $`𝒌`$ and spin $`\sigma `$. We consider only the on-site Coulomb repulsion, $`I`$, as in the previous studies of spin-fluctuation mechanism. Among the three bands, we consider the $`\gamma `$ band consisting of the antibonding band of Ru 4d<sub>xy</sub> and O 2p<sub>π</sub> orbitals in this paper, because it has the largest density of states at Fermi energy and the superconductivity is considered to be realized predominantly in the $`\gamma `$ band. Although the spin fluctuation near $`𝒒_0`$ is understood from the nesting effect of $`\alpha `$ and $`\beta `$ bands, we assume that the wave-number dependence of spin fluctuation is common in the three bands due to some interactions, such as spin-orbit couplings and/or Hund couplings. In the following calculation, we use the $`\gamma `$ band. However the same discussions can be also applied to the $`\alpha `$ and $`\beta `$ bands.
The anisotropy of spin fluctuation observed experimentally is implicitly included in the two-body Hamiltonian, $`H_0`$. Our purpose is not to investigate the origin of anisotropy in details but to examine the role of the anisotropy to Cooper pairing. Therefore we introduce a phenomenological parameter $`\alpha `$ by
$$\chi _{(+,0)}(𝒒)=\alpha \chi _{(,0)}(𝒒),$$
(2)
where $`\chi _{(,0)}(𝒒)`$ ($`\chi _{(+,0)}(𝒒)`$) is the unperturbed static susceptibility of $`z`$ axis ($`xy`$ plane), which originates from $`H_0`$. The parameter $`\alpha `$ represents the anisotropy of spin fluctuation and we take $`\alpha 1`$ since NMR experiments show that $`\chi _{(xx)}<\chi _{(zz)}`$.
Using this one-band model, we discuss the effective interactions between Cooper pairs due to spin fluctuations. Summation of bubble and ladder diagrams (i.e., RPA approximation) gives
$`H_{\mathrm{int}}=`$ $`{\displaystyle \underset{kk^{}s}{}}V_{\mathrm{b}.\mathrm{o}}(𝒌𝒌^{})c_{ks}^{}c_{ks}^{}c_{k^{}s}c_{k^{}s}`$ (5)
$`+{\displaystyle \underset{kk^{}s}{}}V_{\mathrm{b}.\mathrm{e}}(𝒌𝒌^{})c_{ks}^{}c_{ks}^{}c_{k^{}s}c_{k^{}s}`$
$`{\displaystyle \underset{kk^{}s}{}}V_{\mathrm{lad}}(𝒌𝒌^{})c_{ks}^{}c_{ks}^{}c_{k^{}s}c_{k^{}s},`$
with
$`V_{\mathrm{b}.\mathrm{o}}(𝒌𝒌^{})`$ $`=`$ $`{\displaystyle \frac{I}{N}}{\displaystyle \frac{(I/N)\chi _{(,0)}(𝒌𝒌^{})}{1(I/N)^2\chi _{(,0)}^2(𝒌𝒌^{})}},`$ (6)
$`V_{\mathrm{b}.\mathrm{e}}(𝒌𝒌^{})`$ $`=`$ $`{\displaystyle \frac{I}{N}}{\displaystyle \frac{(I/N)^2\chi _{(,0)}^2(𝒌𝒌^{})}{1(I/N)^2\chi _{(,0)}^2(𝒌𝒌^{})}},`$ (7)
$`V_{\mathrm{lad}}(𝒌𝒌^{})`$ $`=`$ $`{\displaystyle \frac{I}{N}}{\displaystyle \frac{(I/N)\chi _{(+,0)}(𝒌𝒌^{})}{1(I/N)\chi _{(+,0)}(𝒌𝒌^{})}}.`$ (8)
Here $`V_{\mathrm{b}.\mathrm{o}}`$ ($`V_{\mathrm{b}.\mathrm{e}}`$) comes from the summation of diagrams with odd (even) number of bubbles, and $`V_{\mathrm{lad}}`$ from the ladder diagrams. It is apparent that $`V_{\mathrm{b}.\mathrm{o}}`$ is between the electrons with equal spins while $`V_{\mathrm{b}.\mathrm{e}}`$ and $`V_{\mathrm{lad}}`$ are between those with the opposite spins.
It is straightforward to derive the gap equations in the anisotropic case, using the method developed by Leggett. First we introduce the operators
$`t_k^{(0)}`$ $`=`$ $`{\displaystyle \underset{ss^{}}{}}(\sigma _2)_{ss^{}}c_{ks}c_{ks^{}},`$ (9)
$`t_k^{(a)}`$ $`=`$ $`{\displaystyle \underset{ss^{}}{}}(\sigma _2\sigma _a)_{ss^{}}c_{ks}c_{ks^{}},\mathrm{for}a=1,2,3,`$ (10)
where $`\sigma _a(a=1,2,3)`$ are Pauli matrices. The operator $`t_k^{(0)}`$ ($`t_k^{(a)}`$) corresponds to spin singlet (triplet) Cooper pairs. In terms of these operators, the effective interaction (5) can be rewritten as
$`H_{\mathrm{int}}`$ $`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle \underset{kk^{}}{}}V_{\mathrm{sin}}(𝒌𝒌^{\mathbf{}})t_k^{(0)}t_k^{}^{(0)}`$ (12)
$`+{\displaystyle \frac{1}{4}}{\displaystyle \underset{kk^{}}{}}{\displaystyle \underset{a=1}{\overset{3}{}}}V_{\mathrm{tri}}^{(a)}(𝒌𝒌^{\mathbf{}})t_k^{(a)}t_k^{}^{(a)},`$
where
$`V_{\mathrm{sin}}(𝒌𝒌^{})`$ $``$ $`2[V_{\mathrm{b}.\mathrm{e}}(𝒌𝒌^{})+V_{\mathrm{lad}}(𝒌𝒌^{})],`$ (13)
$`V_{\mathrm{tri}}^{(1)}(𝒌𝒌^{})`$ $`=`$ $`V_{\mathrm{tri}}^{(2)}(𝒌𝒌^{})2V_{\mathrm{b}.\mathrm{o}}(𝒌𝒌^{}),`$ (14)
$`V_{\mathrm{tri}}^{(3)}(𝒌𝒌^{})`$ $``$ $`2[V_{\mathrm{b}.\mathrm{e}}(𝒌𝒌^{})V_{\mathrm{lad}}(𝒌𝒌^{})].`$ (15)
Since Sr<sub>2</sub>RuO<sub>4</sub> has a long coherence length in $`ab`$ plane, $`\xi _{ab}660`$ Å, we use mean-field approximation to $`H_{\mathrm{int}}`$. We restrict the discussion to unitary states because it is unrealistic to assume non-unitary states in Sr<sub>2</sub>RuO<sub>4</sub>. Requiring that there is no coexistence of singlet and triplet pairs, we obtain the gap equations
$`\mathrm{\Delta }(𝒌)`$ $`=`$ $`{\displaystyle \underset{k^{}}{}}V_{\mathrm{sin}}(𝒌𝒌^{})\mathrm{\Delta }(𝒌^{})\mathrm{\Theta }(E_{\mathrm{sin}}(𝒌^{})),`$ (16)
$`d^{(a)}(𝒌)`$ $`=`$ $`{\displaystyle \underset{k^{}}{}}V_{\mathrm{tri}}^{(a)}(𝒌𝒌^{})d^{(a)}(𝒌^{})\mathrm{\Theta }(E_{\mathrm{tri}}(𝒌^{})),`$ (17)
where $`\mathrm{\Theta }(E)\frac{1}{2E}\text{tanh}\frac{\beta E}{2}`$, $`E_{\mathrm{sin}}^2(𝒌)=\xi _k^2+\mathrm{\Delta }(𝒌)\mathrm{\Delta }^{}(𝒌)`$, and $`E_{\mathrm{tri}}^2(𝒌)=\xi _k^2+𝒅(𝒌)𝒅^{}(𝒌)`$ with $`\xi _k=\epsilon _k\mu `$. The singlet and triplet order parameters are defined as $`\mathrm{\Delta }(𝒌)=\frac{1}{2}_k^{}V_{\mathrm{sin}}(𝒌𝒌^{})t_k^{}^{(0)}`$ and $`d^{(a)}(𝒌)=\frac{1}{2}_k^{}V_{\mathrm{tri}}^{(a)}(𝒌𝒌^{})t_k^{}^{(a)}`$, respectively. Here $`𝒅(𝒌)`$ is the so-called $`𝒅`$-vector for the triplet superconductivity.
In the system with the rotational symmetry in spin space, $`\chi _{(,0)}(𝒒)=\chi _{(+,0)}(𝒒)`$ is satisfied and thus the relation $`V_{\mathrm{b}.\mathrm{o}}+V_{\mathrm{b}.\mathrm{e}}=V_{\mathrm{lad}}`$ holds. In this case, it is easy to see $`V_{\mathrm{tri}}^{(1)}(𝒌𝒌^{})=V_{\mathrm{tri}}^{(2)}(𝒌𝒌^{})=V_{\mathrm{tri}}^{(3)}(𝒌𝒌^{})`$.
On the other hand, the gap equation in Eq. (17) for the triplet pairing becomes dependent on the direction of the $`𝒅`$ vector in the anisotropic case. It means that $`𝐝`$ vector has some preferred direction if the triplet pairs are formed by anisotropic spin fluctuations. This is naturally understood because the $`𝒅`$ vector is orthogonal to the spin direction of triplet Cooper pairs. For the present case with $`\chi _{(+,0)}(𝒒)<\chi _{(,0)}(𝒒)`$ (i.e., $`\alpha <1`$) which is applied to the Sr<sub>2</sub>RuO<sub>4</sub>, we can see from Eq. (8) that $`V_{\mathrm{lad}}(𝒌𝒌^{})`$ is suppressed and the effective interaction $`V_{\mathrm{tri}}^{(3)}(𝒌𝒌^{})`$ approaches $`V_{\mathrm{sin}}(𝒌𝒌^{})`$. Consequently the triplet superconductivity with $`d^{(3)}(𝒌)`$ (i.e., $`𝒅\widehat{𝒛}`$) can be stabilized even due to the AF spin fluctuations.
In order to determine the symmetry of the superconducting order parameter, we have to take account of their sign change along the Fermi surface. For the high-$`T_\mathrm{c}`$ superconductors, the AF spin fluctuation with momentum $`(\pi ,\pi )`$ stabilizes the singlet d$`_{x^2y^2}`$-wave superconductivity. In that case, the singlet order parameters $`\mathrm{\Delta }(𝒌^{})`$ with $`𝒌^{}=(\pi ,0)`$ and $`\mathrm{\Delta }(𝒌)`$ with $`𝒌=(0,\pi )`$ have the opposite sign, so that the gap equation in (17) is satisfied with $`V_{\mathrm{sin}}(\pi ,\pi )>0`$.
For Sr<sub>2</sub>RuO<sub>4</sub> we consider that a kind of nesting property of the Fermi surface plays an important role. This is the second point of our mechanism. Figure 1 shows a schematic Fermi surface for the $`\gamma `$ band. Since the AF fluctuation in Sr<sub>2</sub>RuO<sub>4</sub> has momentum $`𝒒_0`$, the Fermi surface is also shifted by $`(2\pi /3,2\pi /3)`$ in Fig. 1. It is apparent that some part of the shifted Fermi surface overlaps with the original Fermi surface with modulo $`2\pi `$. In analogy to the case of high-$`T_\mathrm{c}`$ superconductivity, if the superconducting order parameters have the opposite sign on these overlapping portions of the Fermi surface, the gap equation is satisfied with $`V_{\mathrm{tri}}^{(a)}(2\pi /3,2\pi /3)>0`$. From Fig. 1, it is natural to consider the p-wave pairing instead of the singlet d$`_{x^2y^2}`$-wave pairing.
In order to clarify this point quantitatively, we compare various kinds of anisotropic superconductivity using the effective interaction and the simplified Fermi surface. Near the transition temperature $`T_\mathrm{c}`$, we rewrite the gap equations as
$$\varphi (𝒌)=\underset{k^{}}{}V_\varphi (𝒌𝒌^{})\varphi (𝒌^{})\frac{1}{2\xi _k^{}}\mathrm{tanh}\frac{\beta _c\xi _k^{}}{2},$$
(18)
where $`\varphi (𝒌)`$ represents one of the order parameters $`\mathrm{\Delta }(𝒌)`$ or $`d^{(a)}(𝒌)`$, and $`V_\varphi `$ is determined from Eqs. (15) depending on $`\varphi `$. In the weak coupling approximation, $`T_\mathrm{c}`$ is obtained as
$$k_BT_\mathrm{c}=1.13\mathrm{}v_Fk_c\mathrm{exp}\left[\frac{1}{N(0)V_\varphi _{FS}}\right],$$
(19)
where $`v_F`$, $`k_c`$ and $`N(0)`$ are the Fermi velocity, cut-off of the wave number, and the density of states at the Fermi energy, respectively. $`V_\varphi _{FS}`$ means the average over the Fermi surface,
$$V_\varphi _{FS}\frac{_{FS}𝑑𝒌_{FS}𝑑𝒌^{}V_\varphi (𝒌𝒌^{})\varphi (𝒌)\varphi (𝒌^{})}{\left[_{FS}𝑑𝒌^{}\right]_{FS}𝑑𝒌\varphi ^2(𝒌)}.$$
(20)
We identify that the order parameter which gives the largest $`N(0)V_\varphi _{FS}`$ is realized.
For Sr<sub>2</sub>RuO<sub>4</sub> we choose order parameters $`\varphi (𝒌)`$ as follows
$`\varphi _1(𝒌)`$ $`=`$ $`\mathrm{cos}k_x+\mathrm{cos}k_y,`$ (21)
$`\varphi _2(𝒌)`$ $`=`$ $`\mathrm{cos}k_x\mathrm{cos}k_y,`$ (22)
$`\varphi _3(𝒌)`$ $`=`$ $`\mathrm{sin}k_x\mathrm{sin}k_y,`$ (23)
$`\varphi _4(𝒌)`$ $`=`$ $`\mathrm{sin}k_x,(\widehat{𝒅}\widehat{𝒛}),`$ (24)
$`\varphi _5(𝒌)`$ $`=`$ $`\mathrm{sin}k_x,(\widehat{𝒅}\widehat{𝒛}),`$ (25)
where $`\varphi _1\varphi _3`$ correspond to singlet pairings, and $`\varphi _4`$, $`\varphi _5`$ to triplet pairings, respectively. The most probable candidate for Sr<sub>2</sub>RuO<sub>4</sub> is $`\widehat{𝒛}(\mathrm{sin}k_x\pm i\mathrm{sin}k_y)`$ which is equivalent to $`\varphi _5`$ just below $`T_\mathrm{c}`$, because the gap equation (18) for $`\mathrm{sin}k_x\pm i\mathrm{sin}k_y`$ is exactly same as that for $`\varphi _5`$. If $`N(0)V_{\varphi _5}_{FS}`$ is the largest, we expect that the order parameter $`𝒅(𝒌)=\widehat{𝒛}(\mathrm{sin}k_x\pm i\mathrm{sin}k_y)`$ is realized, because near zero temperature it acquires a larger energy gap than $`\varphi _5`$.
To emphasize the characteristic feature of the nesting, we assume the simplified Fermi surface as shown in Fig. 1. For the $`𝒒`$ dependence of $`\chi _{(,0)}(𝒒)`$ with a maximum at $`q_0`$, we use the susceptibility obtained in the LDA calculation, and fix $`S(\mathrm{𝟎})=0.8`$ with $`S(𝒒)\frac{I}{N}\chi _{(,0)}(𝒒)`$. We regard $`S(𝒒_0)`$ as a phenomenological parameter.
Figure 2 shows the $`\alpha `$ dependence of $`N(0)V_{\varphi _n}_{FS}`$ $`(n=15)`$ for $`S(𝒒_0)=0.95`$. We examined various choices of $`S(𝒒_0)`$ from $`0.90`$ to $`0.99`$ to find that the results do not change qualitatively. When the anisotropy is weak ($`\alpha 1`$), the singlet d$`_{x^2y^2}`$-wave superconductivity, $`\varphi _2`$, is stabilized. On the other hand, when $`\alpha `$ is small, the order parameter $`\varphi _5`$ is stabilized which is consistent with experiments.
The phase diagram as a function of $`\alpha `$ and $`S(𝒒_0)`$ is determined by examining various values of $`S(𝒒_0)`$. Because it is unphysical to assume that $`S(𝒒_0)`$ is very close to 1, we show the results up to $`S(𝒒_0)=0.99`$ in Fig. 3. When the spin fluctuation is isotropic (i.e., $`\alpha =1`$), the singlet d$`_{x^2y^2}`$-wave superconductivity, $`\mathrm{cos}k_x\mathrm{cos}k_y`$, is the most stable pairing. This is consistent with the previous studies. However we find a fairly large parameter region where the state corresponding to $`\widehat{𝒛}(\mathrm{sin}k_x\pm i\mathrm{sin}k_y)`$ is realized.
Finally we discuss the competition between the singlet $`\mathrm{cos}k_x\mathrm{cos}k_y`$ pairing and the triplet $`\widehat{𝒛}(\mathrm{sin}k_x\pm i\mathrm{sin}k_y)`$ pairing in terms of the effective interaction and the nesting property. From the explicit form of $`V_{\varphi _n}`$ for $`n=2`$ and 5, we can see that a relation $`V_{\varphi _2}V_{\varphi _5}`$ is satisfied. Therefore if we consider only the magnitude of the effective interaction, the singlet pairing is favorable. However the nesting property favors the triplet pairing. Let us assume that $`V_{\varphi _n}`$ is enhanced very strongly by the AF fluctuation and approximated as
$$V_{\varphi _n}(𝒒)=\frac{I}{N}A_n\widehat{\delta }(q_x\pm 2\pi /3)\widehat{\delta }(q_y\pm 2\pi /3),$$
(26)
with $`\widehat{\delta }`$ being the $`\delta `$ function with modulo $`2\pi `$. Using this approximated form of $`V_{\varphi _n}(n=2\mathrm{and}5)`$, we obtain
$`N(0)V_{\varphi _2}_{FS}`$ $`=`$ $`[2.79\times 10^2\delta (0)+4.91\times 10^2]A_2,`$
$`N(0)V_{\varphi _5}_{FS}`$ $`=`$ $`[4.24\times 10^2\delta (0)+5.06\times 10^2]A_5.`$
This estimation shows that the $`\widehat{𝒛}(\mathrm{sin}k_x\pm i\mathrm{sin}k_y)`$ pairing utilizes the peak of $`\chi _{(,0)}(𝒒)`$ at $`𝒒_0`$ more effectively than $`\mathrm{cos}k_x\mathrm{cos}k_y`$ pairing does. Therefore, even if $`A_2>A_5`$, the triplet pair can be stabilized.
In determining the phase diagram in Fig. 3, we have assumed simple functional forms of the order parameters, $`\varphi _n(𝒌)`$. For the detailed calculations, it will be necessary to optimize the $`𝒌`$-dependence of $`\varphi _n(𝒌)`$. However the global feature of the phase diagram will not change.
In summary, we have generalized the RPA formulation of the effective interaction due to the spin fluctuations and derived gap equations including the anisotropic case. We have shown that the state corresponding to $`\widehat{𝒛}(\mathrm{sin}k_x\pm i\mathrm{sin}k_y)`$ becomes the most stable even if the AF spin fluctuation is dominant, when the anisotropy is strong enough and the nesting property of the Fermi surface is present. Although the nesting for the actual Fermi surface will be weaker than what we assumed here, it is reasonable to think that our mechanism is the most promising one as far as the AF fluctuation is dominant.
In this paper we have investigated the pairing in the $`\gamma `$ band. However it is straightforward to consider the other bands ($`\alpha `$ and $`\beta `$ bands) in Sr<sub>2</sub>RuO<sub>4</sub>. Since the nesting property will be comparable or even stronger for these bands than for the $`\gamma `$ band, we expect the same mechanism for triplet superconductivity for $`\alpha `$ and $`\beta `$ bands even if the $`\gamma `$ band does not have the peak near $`𝒒_0`$.
It is reported that Sr<sub>2</sub>RuO<sub>4</sub> has exotic property called as 3K phase when Ru metal is embedded in the single crystal. We speculate that the enhancement of $`T_\mathrm{c}`$ is due to the increase of the anisotropy (i.e., decrease of $`\alpha `$) near the interface region between Sr<sub>2</sub>RuO<sub>4</sub> and Ru metal. We consider that to investigate the origin of the anisotropy is very important both for understanding the superconductivity and for finding the new exotic phenomena.
We are grateful for useful discussions with M. Sigrist, Y. Maeno and Y. Matsuda. One of the authors (T.K.) thanks to H. Namaizawa for useful instructions.
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# 1 Introduction
## 1 Introduction
Recently there has been a considerable effort in studying the AdS/CFT duality when conformal invariance is broken (see and references there in). The general philosophy is that conformal invariance can be broken by considering a non-conformal vacuum of the conformal theory, or otherwise by considering a deformation of its Lagrangian. In both cases the dual gravitational background will have an AdS boundary but it will differ from this space in the inside. Identifying the radial coordinate with the field theory energy scale the geometry will describe the RG flow starting from the UV conformal theory at the AdS boundary.
In this note we shall consider the case of a non-conformal vacuum of $`SU(N)`$ $`𝒩=4`$ Super Yang-Mills (SYM) theory. The moduli space of this theory is given by $`\text{}^{6(N1)}/S_N`$, parameterizing the relative position of $`N`$ identical D3-branes in the transverse space $`\text{𝔼}^6`$. Separating the D3-branes in two stacks of $`N_1`$ and $`N_2`$ branes by a distance $`2\mathrm{\Delta }`$ is equivalent to Higgs the theory by giving an expectation value to the scalar fields according to
$$\frac{\stackrel{}{Y}}{2\pi \alpha ^{}}\stackrel{}{\varphi }=\frac{\stackrel{}{\mathrm{\Delta }}}{2\pi \alpha ^{}}\left(\begin{array}{cc}I_1& 0\\ 0& I_2\end{array}\right),$$
(1)
where $`I_i`$ is the $`N_i\times N_i`$ unit matrix. The gauge symmetry is broken to $`S(U(N_1)\times U(N_2))`$ and the theory is left with 16 conserved supercharges because conformal symmetry is broken. The symmetry breaking process gives rise to massive W particles with mass $`m_W=\mathrm{\Delta }/(\pi \alpha ^{})`$. This construction provides an example of maximally supersymmetric RG flow: starting from the UV $`SU(N)`$ conformal theory we flow to the IR conformal fixed points with $`SU(N_i)`$ gauge symmetry. Related aspects of the AdS/CFT duality can be found in references \[5-23\].
To derive the dual gravitational background we start with the double-centered D3-brane geometry with metric
$$ds^2=H^{1/2}ds^2(\text{𝕄}^4)+H^{1/2}ds^2(\text{𝔼}^6),$$
(2)
where the harmonic function $`H`$ reads
$$H=1+\frac{R_1^{4}}{|\stackrel{}{y}\stackrel{}{\mathrm{\Delta }}|^4}+\frac{R_2^{4}}{|\stackrel{}{y}+\stackrel{}{\mathrm{\Delta }}|^4},$$
(3)
with $`R_i^{4}=4\pi g_s\alpha ^2N_i`$ and $`\stackrel{}{y}`$ the coordinates in the transverse space $`\text{𝔼}^6`$. Next, we take the decoupling limit of which corresponds to drop the factor 1 in the harmonic function $`H`$ and to consider energy scales in the brane theory much smaller than the string scale, i.e. $`ϵ,m_W1/\sqrt{\alpha ^{}}`$. We further require the string coupling $`g_s`$ to be small and $`g_sN_i`$ to be large for the supergravity approximation to hold.
In a recent paper we studied the classical absorption for a minimally coupled scalar, e.g. the dilaton field, in the double-centered D3-brane geometry. The analysis was valid for low energies such that $`ϵm_gm_W`$, where $`m_g=m_W/\sqrt{gN}`$ is the gravity mass gap. In the above decoupling limit the cross-section for absorption by the $`i`$-th hole was seen to be
$$\sigma _i=\frac{\pi ^4\omega ^3R_i^{8}}{8}\left[1\frac{(\omega R_i)^4}{12}\left(\frac{R_j}{2\mathrm{\Delta }}\right)^4\mathrm{log}\left(\frac{\omega R_i^{2}}{\mathrm{\Delta }}\right)+\mathrm{}\right],(ji).$$
(4)
In the field theory side, the total absorption cross section (i.e. $`\sigma =\sigma _1+\sigma _2`$) for a scalar field $`\varphi `$ is related to the two-point function of the field theory dual operator $`𝒪_\varphi `$
$$\mathrm{\Pi }(x)=𝒪_\varphi ^{}(x)𝒪_\varphi (0),$$
(5)
calculated in the non-conformal vacuum described above. In the case of the dilaton field the exact relation is
$$\sigma =\frac{2\kappa ^2}{2i\omega }\mathrm{Disc}^{}\mathrm{\Pi }(s)|_{\genfrac{}{}{0pt}{}{p^0=\omega }{\stackrel{}{p}=\stackrel{}{0}}},$$
(6)
where $`\kappa `$ is the gravitational coupling, $`s=p^2`$ and $`\mathrm{\Pi }(s)`$ is the Fourier transform of $`\mathrm{\Pi }(x)`$.
The low energy expansion of the classical cross section (4) corresponds to consider the dual field theory near the $`SU(N_i)`$ IR fixed point. The purpose of this letter is to determine the exact form of the effective action near the IR fixed points that results from integrating out the W particles and therefore to reproduce exactly the logarithmic correction to the cross section using the dual field theory. We shall see that the effective action is given by a deformation of the $`𝒩=4`$ SYM theory by a non-renormalized, irrelevant, dimension 8 operator as anticipated in . Then, the correction to the cross section will be related to a three-point function of chiral primary operators calculated at the IR fixed points . To obtain exact agreement with the classical cross section calculation it was essential to use the symmetrized trace of the Yang Mills fields in the deformed Lagrangian and to keep only the planar diagrams. This result is further evidence for a non-renormalization theorem for three-point functions of chiral primary operators \[27-32\] and provides a test of the AdS/CFT duality on the Coulomb branch.
## 2 Effective Action in the Infrared
We start by writing the Lagrangian for the bosonic sector of $`𝒩=4`$ $`SU(N)`$ SYM theory in the following form
$$_0=\frac{1}{4}\mathrm{tr}\left[F_{AB}F^{AB}\right],$$
(7)
where $`A,B,\mathrm{}`$ are ten-dimensional indices and $`F_{AB}`$ is short for $`F_{\mu \nu }=D_\mu A_\nu D_\nu A_\mu `$, $`F_{\mu m}=D_\mu \varphi _m`$ with $`D_\mu _\mu +ig_{YM}[A_\mu ,]`$ and $`F_{mn}=ig_{YM}[\varphi _m,\varphi _n]`$. We removed the gauge coupling from the front of the action by rescaling the fields as $`(A_\mu ,\varphi ^m)g_{YM}(A_\mu ,\varphi ^m)`$, which also rescales the expectation value for the scalars in equation (1).
We want to integrate the W’s in order to obtain an effective action for the light $`SU(N_1)`$ and $`SU(N_2)`$ coloured fields at energy scales $`ϵm_W`$ as explained in . This is similar to the probe calculations extensively considered in the literature \[33-36\], where $`N_1=N`$ and $`N_2=1`$ . In the latter case the first non-vanishing one-loop diagram involves 4 $`SU(N)`$ coloured legs, with the following contribution to the effective bosonic Lagrangian
$$_1=\frac{\pi ^2g_{YM}^4}{(2\mathrm{\Delta }/\alpha ^{})^4}\mathrm{Str}\left[F_{AB}F^{BC}F_{CD}F^{DA}\frac{1}{4}\left(F_{AB}F^{AB}\right)^2\right].$$
(8)
Notice that we are using the symmetrized trace as argued in . This fact will be essential to obtain agreement with the dual classical calculation of the absorption cross-section. Also, these $`F^4`$ terms are protected , i.e. they are not renormalized by higher loop diagrams and therefore comparison with the strongly coupled supergravity regime is allowed.
Now we generalize the above result to arbitrary large $`N_1`$ and $`N_2`$. The only difference is that the term in the effective action involving 4 $`SU(N_i)`$ fields will be multiplied by a factor $`N_j`$ $`(ij)`$ because we have a $`SU(N_j)`$ colour index to sum over around the loop (due to the W’s exchange). Also, for large $`N`$ graphs with both $`SU(N_1)`$ and $`SU(N_2)`$ coloured external legs are subleading since they are associated with non-planar graphs. Hence, in the IR and for large $`N`$ there are no terms in the effective action of the type $`\mathrm{tr}_1(F^2)\mathrm{tr}_2(F^2)`$, where the subscript in $`\mathrm{tr}_i`$ means that the trace is over $`SU(N_i)`$ fields. We conclude that the $`F^4`$ terms in the effective action for the IR $`SU(N_i)`$ theory read
$$_1=\frac{\pi ^2g_{YM}^4}{(2\mathrm{\Delta }/\alpha ^{})^4}N_j\mathrm{Str}_i\left[F_{AB}F^{BC}F_{CD}F^{DA}\frac{1}{4}\left(F_{AB}F^{AB}\right)^2\right],(ji).$$
(9)
To check this result consider the action for a probe of $`N_i`$ D3-branes in the AdS near-horizon geometry of $`N_j`$ D3-branes, i.e. we assume that $`N_jN_i`$. The probe dynamics is determined by the non-abelian DBI action which describes the effect of integrating all the massive strings stretching between the probe and the branes. The AdS background is described by the metric (2) with the harmonic function $`Hf=(R_j/r)^4`$. If the center of mass for the $`N_i`$ probes is placed at $`r`$ we have
$$S_{probe}=T_3d^4xf^1\mathrm{Str}_i\left[\sqrt{\mathrm{det}\left(\eta _{\alpha \beta }+f_\alpha Y^m_\beta Y_m+2\pi \alpha ^{}\sqrt{f}F_{\alpha \beta }\right)}I_i\right],$$
(10)
where $`T_3=((2\pi )^3g_s\alpha ^2)^1`$ is the D3-brane tension. Next we give an expectation value to the scalars $`\stackrel{}{Y}=2\stackrel{}{\mathrm{\Delta }}I_i`$, which means that $`r=2\mathrm{\Delta }+\delta r`$ where $`\delta r`$ is the center of mass fluctuating field in the radial direction. For large $`N_i`$ we can discard the fields in the center of the gauge group and consider only $`SU(N_i)`$ fields. Thus, setting $`f=(R_j/2\mathrm{\Delta })^4`$, expanding the DBI action and rescaling the fields according to $`(A_\alpha ,\varphi _m=(2\pi \alpha ^{})^1Y_m)g_{YM}(A_\alpha ,\varphi _m)`$ we obtain the probe Lagrangian
$$_{probe}=\frac{1}{4}\mathrm{tr}_i\left[F_{AB}F^{AB}\right]\frac{\pi ^2g_{YM}^4N_j}{(2\mathrm{\Delta }/\alpha ^{})^4}\mathrm{Str}_i\left[F_{AB}F^{BC}F_{CD}F^{DA}\frac{1}{4}\left(F_{AB}F^{AB}\right)^2\right],$$
(11)
which is in agreement with the previous result as expected from the non-renormalization of the $`F^4`$ terms.
In the context of the asymptotically flat D3-brane geometry, the form of the above deformation of the SYM Lagrangian was conjectured based on the DBI corrections to the SYM theory , or alternatively on the basis of $`PSU(2|2,4)`$ representation theory . For a geometry with harmonic function
$$H=h+\left(\frac{R}{r}\right)^4,$$
(12)
the dual field theory was conjectured to be
$$=_0\frac{h}{8T_3}𝒪_8,$$
(13)
where $`𝒪_8`$ is an irrelevant, dimension 8 operator. This deformation is irrelevant in the IR which agrees with the fact that for $`r0`$ the constant $`h`$ becomes irrelevant in the harmonic function $`H`$. What remains an open question is if the Lagrangian (13) describes the dual gravity theory as we flow from the IR. If we assume that the DBI action is dual to the full D3-brane geometry $`(h=1)`$, then we can regard (13) as the first correction to the SYM theory and determine $`𝒪_8`$ to be exactly $`𝒪_8=\mathrm{Str}[F^4\frac{1}{4}(F^2)^2]`$. However, even in this case agreement between the gravity and the field theory calculations of the cross section is not found . Fortunately the case here studied is entirely under control. While in the asymptotic flat space case the DBI action arises from integrating the massive string states that are dropped out in the decoupling limit, in our case the deformation of the IR Lagrangian arises from integrating the W’s that do survive the decoupling limit. In fact, in the decoupling limit of the double-centered D3-brane geometry the harmonic function $`H`$ is given down the $`i`$-th throat by
$$H=\left(\frac{R_j}{2\mathrm{\Delta }}\right)^4+\left(\frac{R_i}{r}\right)^4.$$
(14)
Then the deformation of the SYM action in the IR $`SU(N_i)`$ conformal fixed point is indeed given by (13) with $`h=(R_j/2\mathrm{\Delta })^4`$ and $`𝒪_8=\mathrm{Str}_i[F^4\frac{1}{4}(F^2)^2]`$ . This (non-renormalized) deformation was computed exactly as a result of integrating the W’s. Hence, if we believe the AdS/CFT correspondence to hold on the Coulomb branch, the gravity and perturbative field theory calculation of protected quantities in the IR using the Lagrangian (13) must give exactly the same answer.
## 3 Field Theory Calculation of Cross-Section
We proceed by calculating the cross section for absorption of the dilaton field using the field theory approach. This calculation was done in using a $`U(1)`$ model. We refer the reader to that reference for the details and will present here only the relevant steps necessary to obtain the correct answer. For world-volume on-shell processes that involve the coupling of the dilaton to the brane only the gauge field is relevant. In the IR $`SU(N_i)`$ theory the operator $`𝒪_\varphi `$ dual to the dilaton field reads
$$\begin{array}{ccc}\hfill S_{int}& =& d^4x\varphi \frac{1}{4}\left(\mathrm{tr}_i\left[F_{\mu \nu }F^{\mu \nu }\right]+\frac{h}{T_3}\mathrm{Str}_i\left[F_{\mu \nu }F^{\nu \eta }F_{\eta \lambda }F^{\lambda \mu }\frac{1}{4}\left(F_{\mu \nu }F^{\mu \nu }\right)^2\right]\right),\hfill \\ & & d^4x\varphi 𝒪_\varphi d^4x\varphi \frac{1}{4}\left[𝒪_4+\frac{h}{T_3}𝒪_8\right].\hfill \end{array}$$
(15)
Then the two-point function for the operator $`𝒪_\varphi `$ is
$$\begin{array}{ccc}\hfill \mathrm{\Pi }(x)& =& 𝒪_\varphi ^{}(x)𝒪_\varphi (0)_{h=\left(R_j/2\mathrm{\Delta }\right)^4}=𝒟A_\mu e^{{\scriptscriptstyle d^4z\left[{\scriptscriptstyle \frac{1}{4}}𝒪_4+{\scriptscriptstyle \frac{h}{8T_3}}𝒪_8\right]}}𝒪_\varphi (x)𝒪_\varphi (0)\hfill \\ & =& 𝒪_\varphi (x)𝒪_\varphi (0)\left(1\left(\frac{R_j}{2\mathrm{\Delta }}\right)^4\frac{1}{8T_3}d^4z𝒪_8(z)\right)_{h=0}\hfill \\ & =& \frac{1}{2^4}𝒪_4^{}(x)𝒪_4(0)_{h=0}\left(\frac{R_j}{2\mathrm{\Delta }}\right)^4\frac{1}{2^7T_3}d^4z𝒪_4^{}(x)𝒪_8(z)𝒪_4(0)_{h=0}\hfill \\ & & \mathrm{\Pi }_0(x)+\mathrm{\Pi }_1(x),\hfill \end{array}$$
(16)
where we are just keeping the leading terms that will give the logarithmic corrections in the classical gravity result for the cross section. We see that the correction to the two-point function is related to a three-point function of chiral primary operators calculated at a IR conformal fixed point.
We start by writing the Euclidean space propagator for the field strength $`(F_{\mu \nu })^{ab}`$ in the conformal theory
$$\begin{array}{c}(F\text{ }\overline{{}_{\mu \nu }{}^{})^{ab}(x)(F}\text{ }_{\alpha \beta })^{cd}(0)\left(F_{\mu \nu }\right)^{ab}(x)\left(F_{\alpha \beta }\right)^{cd}(0)=\\ =\frac{\delta ^{ad}\delta ^{bc}}{\pi ^2x^4}\left[\delta _{\mu \alpha }\delta _{\nu \beta }\delta _{\nu \alpha }\delta _{\mu \beta }\frac{2}{x_{}^2}\left(\delta _{\nu \beta }^{}x_\mu x_\alpha +\delta _{\mu \alpha }x_\nu x_\beta \delta _{\nu \alpha }x_\mu x_\beta \delta _{\mu \beta }x_\nu x_\alpha \right)\right],\end{array}$$
(17)
where $`a,b,\mathrm{}`$ are $`SU(N_i)`$ colour indices. In what follows it is convenient to write first the following contraction of the field strength propagators
$$\left((F\text{ }\overline{{}_{\alpha \beta }{}^{})^{ab}(F\text{ }\overline{{}_{}{}^{ba}{}_{}{}^{\alpha \beta }))(x)((F}\text{ }_{\mu \nu })^{cd}(F}\text{ }_{\eta \lambda })^{ef}\right)(0)=\frac{2\delta ^{de}\delta ^{cf}}{\pi ^4x^8}\left(\delta _{\mu \eta }^{}\delta _{\nu \lambda }\delta _{\mu \lambda }\delta _{\nu \eta }\right).$$
(18)
It is trivial to check that $`\mathrm{\Pi }_0(x)\frac{1}{2^4}𝒪_4^{}(x)𝒪_4(0)_{h=0}=(3N_i^{2})(\pi ^4x^8)`$ . To calculate $`\mathrm{\Pi }_1(x)`$ first we expand the symmetrized trace of $`𝒪_8`$ in terms of simple traces as
$$\begin{array}{ccc}\hfill 𝒪_8(z)=\mathrm{Str}_i\left[F^4\frac{1}{4}(F^2)^2\right]& =& \frac{2}{3}\mathrm{tr}_i[F_{\mu \nu }^{}F^{\nu \eta }F^{\mu \lambda }F_{\lambda \eta }+\frac{1}{2}F_{\mu \nu }F^{\nu \eta }F_{\eta \lambda }F^{\lambda \mu }\hfill \\ & & \frac{1}{4}F_{\mu \nu }F^{\mu \nu }F_{\eta \lambda }F^{\eta \lambda }\frac{1}{8}F_{\mu \nu }F_{\eta \lambda }F^{\mu \nu }F^{\eta \lambda }].\hfill \end{array}$$
(19)
Then using standard perturbative field theory techniques we have (see also )
$$𝒪_4^{}(x)𝒪_8(z)𝒪_4(0)_{h=0}=\frac{(3\times 2^8)N_i^{3}}{\pi ^8(xz)^8z^8},$$
(20)
and therefore
$$\mathrm{\Pi }_1(x)=\left(\frac{R_j}{2\mathrm{\Delta }}\right)^4\frac{1}{2^7T_3}d^4z\frac{(3\times 2^8)N_i^{3}}{\pi ^8(xz)^8z^8}.$$
(21)
Notice that we have to be careful in applying Wick’s theorem to obtain this result. The reason is that for large $`N_i`$ only contractions between the fields in $`𝒪_4`$ and fields that are consecutive in the trace expansion (19) of $`𝒪_8`$ will contribute. The other contractions correspond to non-planar graphs and are subleading in the large $`N_i`$ limit<sup>2</sup><sup>2</sup>2I thank Igor Klebanov for bringing this point to my attention. (see figure). A consequence of this fact is that it was essential that we used the symmetrized trace in the effective action, otherwise we would obtain a different answer. Using the result
$$d^4u\frac{e^{ipu}}{u^8}=\frac{\pi ^2}{3\times 2^6}p^4\mathrm{log}\left(p^2/\mathrm{\Lambda }^2\right),$$
(22)
a simple calculation gives the Fourier transform of $`\mathrm{\Pi }_1(x)`$
$$\mathrm{\Pi }_1(p)=\left(\frac{R_j}{2\mathrm{\Delta }}\right)^4\frac{N_i^{3}}{(3\times 2^{11})\pi ^4}p^8\left(\mathrm{log}\left(p^2/\mathrm{\Lambda }^2\right)\right)^2,$$
(23)
where $`\mathrm{\Lambda }`$ is a cut-off scale. Then the absorption cross section is related by equation (6) to the momentum space two-point function $`\mathrm{\Pi }(p)`$. The result is
$$\sigma _i=\frac{\kappa ^2\omega ^3N_i^{2}}{32\pi }\left[1\left(\frac{R_j}{2\mathrm{\Delta }}\right)^4\frac{N_i\omega ^4}{(3\times 2^3)T_3\pi ^2}\mathrm{log}\left(\frac{\omega }{\mathrm{\Lambda }}\right)\right],$$
(24)
which agrees exactly with the classical gravity prediction, including the numerical factors. We remark that in the perturbative field analysis it seems natural to set the cut-off scale to $`\mathrm{\Lambda }=m_W`$, while the strong coupled gravity calculation suggests that $`\mathrm{\Lambda }=\mathrm{\Delta }/R_i^{2}m_W/\sqrt{gN}`$. This fact may be related to the existence of colour singlet condensates of W particles at strong coupling with a large binding energy . However, a better understanding of the cut-off scale would require the extension of this calculation to the next order .
In resume, we tested the AdS/CFT duality on the Coulomb branch by finding agreement between the gravity and field theory absorption cross sections for the dilaton field near the $`SU(N_i)`$ IR conformal fixed points. This was done by determining the large $`N`$ (non-renormalized) $`F^4`$ terms in the field theory IR effective action that arise from integrating the massive W particles. Then the correction to the absorption cross section is related to a three-point function of chiral primary operators. To obtain the correct result required large $`N`$ considerations as well as the use of the symmetrized trace in the effective action. This result provides further evidence for a non-renormalization theorem for three-point functions of chiral primary operators \[27-32\].
## Acknowledgments
I would like to thank Lori Paniak and Igor Klebanov for many discussions. This work was supported by FCT (Portugal) under programme PRAXIS XXI and by the NSF grant PHY-9802484.
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