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# Coordinate Independence of of Quantum-Mechanical Path Integrals
## I Introduction
In the previous papers , we have presented a diagrammatic proof of reparametrization invariance of perturbatively defined quantum-mechanical path integrals. The proper perturbative definition of path integrals was shown to require an extension to a functional integral in $`D`$ spacetime, and a subsequent analytic continuation to $`D=1`$. In Ref. the perturbative calculations were performed in momentum space, where Feynman integrals in a continuous number of dimensions $`D`$ are known from the prescriptions of ’t Hooft and M. Veltman. In Ref. we have found the same results directly from the Feynman integrals in the $`1\epsilon `$-dimensional time space with the help of the Bessel representation of Green functions. The coordinate space calculation is interesting for many applications, for instance, if one wants to obtain the effective action of a field system in curvilinear coordinates, where the kinetic term depends on the dynamic variable. Then one needs rules for performing temporal integrals over Wick contractions of local fields.
In this note we want to show that the reparametrization invariance of perturbatively defined quantum-mechanical path integrals can be obtained in the coordinate space with the help of a simple but quite general arguments based on the inhomogeneous field equation for the Green function, and rules of the partial integration. The prove does not require the calculation of the Feynman integrals separately and remains valid for the functional integrals in an arbitrary space-time dimension $`D`$.
## II Problem with Coordinate Transformations
Recall the origin of the difficulties with coordinate transformations in path integrals. Let $`x(\tau )`$ be the euclidean coordinates of a quantum-mechanical point particle of unit mass in a harmonic potential $`\omega ^2x^2/2`$ as a function of the imaginary time $`\tau =it`$. Under a coordinate transformation $`x(\tau )q(\tau )`$ defined by $`x(\tau )=f(q(\tau ))=q(\tau )+_{n=2}^{\mathrm{}}a_nq^n(\tau )`$, the kinetic term $`\dot{x}^2(\tau )/2`$ goes over into $`\dot{q}^2(\tau )f^{}{}_{}{}^{2}(q(\tau ))/2`$. If the path integral over $`q(\tau )`$ is performed perturbatively, the expansion terms contains temporal integrals over Wick contractions which, after suitable partial integrations, are products of the following basic correlation functions
$`\mathrm{\Delta }(\tau \tau ^{})`$ $``$ $`q(\tau )q(\tau ^{})=\text{ }\text{}\text{ },`$ (1)
$`_\tau \mathrm{\Delta }(\tau \tau ^{})`$ $``$ $`\dot{q}(\tau )q(\tau ^{})=\text{ }\text{}\text{ },`$ (2)
$`_\tau _\tau ^{}\mathrm{\Delta }(\tau \tau ^{})`$ $``$ $`\dot{q}(\tau )\dot{q}(\tau ^{})=\text{ }\text{}\text{ }.`$ (3)
The right-hand sides define the line symbols to be used in Feynman diagrams for the interaction terms.
Explicitly, the first correlation function reads
$$\mathrm{\Delta }(\tau ,\tau ^{})=\frac{1}{2\omega }e^{\omega |\tau \tau ^{}|}.$$
(4)
The second correlation function (2) has a discontinuity
$$_\tau \mathrm{\Delta }(\tau ,\tau ^{})=\frac{1}{2}ϵ(\tau \tau ^{})e^{\omega |\tau \tau ^{}|},$$
(5)
where
$$ϵ(\tau \tau ^{})2_{\mathrm{}}^\tau 𝑑\tau ^{\prime \prime }\delta (\tau ^{\prime \prime }\tau ^{})$$
(6)
is a distribution which has a jump at $`\tau =\tau ^{}`$. The third correlation function (3) contains a $`\delta `$-function:
$$_\tau _\tau ^{}\mathrm{\Delta }(\tau ,\tau ^{})=\delta (\tau \tau ^{})\frac{\omega }{2}e^{\omega |\tau \tau ^{}|},$$
(7)
The temporal integrals in $`\tau `$-space over products of such distributions are undefined.
In our previous papers we have shown that a unique perturbation expansions leading to a reparametrization invariant theory is obtained by extending the path integral to a $`D`$-dimensional functional integral, and by performing the perturbation expansion in $`D`$-dimensional space, with a limit $`D1`$ taken at the end.
In this note we shall set up simple rules for integrals over products of the correlation functions in the extended coordinate space the same results.
## III Model System
To be specific, we shall prove the coordinate independence of the exactly solvable path integral of a point particle of unit mass in a harmonic potential $`\omega ^2x^2/2`$, over a large imaginary-time interval $`\beta `$,
$$Z_\omega =𝒟x(\tau )e^{𝒜_\omega [x]}=e^{\mathrm{Tr}\mathrm{log}(^2+\omega ^2)}=e^{\beta \omega /2}.$$
(8)
The action is
$$𝒜_\omega =\frac{1}{2}𝑑\tau \left[\dot{x}^2(\tau )+\omega ^2x^2(\tau )\right].$$
(9)
A coordinate transformation turns (8) into a path integral with a singular perturbation expansion.
For simplicity we assume the coordinate transformation to preserve the symmetry $`xx`$ of the initial oscillator, such its power series expansion starts out like $`x(\tau )=f(q(\tau ))=qgq^3/3+g^2aq^5/5\mathrm{}`$, where $`g`$ is a smallness parameter, and $`a`$ an extra parameter. We shall see that the perturbation expansion is independent of $`a`$, such that $`a`$ will merely serve to check the calculations. The transformation changes the partition function (8) into
$$Z=𝒟q(\tau )e^{𝒜_J[q]}e^{𝒜[q]},$$
(10)
where the transformed action $`𝒜[q]=𝒜_\omega [q]+𝒜_{\mathrm{int}}[q]`$ is decomposed into a free part
$$𝒜_\omega [q]=\frac{1}{2}𝑑\tau [\dot{q}^2(\tau )+\omega ^2q^2(\tau )],$$
(11)
and an interacting part, which reads to second order in $`g`$:
$`𝒜_{\mathrm{int}}[q]={\displaystyle \frac{1}{2}}{\displaystyle }d\tau \{g[2\dot{q}^2(\tau )q^2(\tau )+{\displaystyle \frac{2\omega ^2}{3}}q^4(\tau )]`$ (12)
$`+g^2[(1+2a)\dot{q}^2(\tau )q^4(\tau )+\omega ^2({\displaystyle \frac{1}{9}}+{\displaystyle \frac{2a}{5}})q^6(\tau )]\}.`$ (13)
The exponent in (10) contains an additional effective action $`𝒜_J[q]`$ coming from the Jacobian of the coordinate transformation:
$$𝒜_J[q]=\delta (0)𝑑\tau \mathrm{log}\frac{\delta f(q(\tau ))}{\delta q(\tau )}.$$
(14)
This has the power series expansion
$`𝒜_J[q]=\delta (0){\displaystyle 𝑑\tau \left[gq^2(\tau )+g^2\left(a\frac{1}{2}\right)q^4(\tau )\right]}.`$ (15)
For $`g=0`$, the transformed partition function (10) coincides with (8). When expanding $`Z`$ of Eq. (10) in powers of $`g`$, we obtain a sum of Wick contractions with associated Feynman diagrams contributing to each order $`g^n`$. This sum must vanish to ensure coordinate invariance of the path integral.
By considering only connected Feynman diagrams, we shall obtain an expansion for the free energy
$$F=F_\omega +\underset{n=1}{}g^nF_n,$$
(16)
where $`F_\omega `$ is the free energy of the unperturbed harmonic oscillator (8). The coordinate invariance is ensured by the vanishing of all expansion terms $`F_n`$.
## IV Expansion Terms of Free Energy Density
The graphical expansion for the ground state energy will be carried here only up to three loops. The diagrams are composed of the three types of lines in (1)–(3), and new interaction vertices for each power of $`g`$. The diagrams coming from the Jacobian action (15) are easily recognized by an accompanying power of $`\delta (0)`$.
To lowest order in $`g`$, there exists only three diagrams, two originated from the interaction (13), one from the Jacobian action (15):
$$F_1=g\left[\text{ }\text{}\text{ }+\omega ^2\text{ }\text{}\text{ }\delta (0)\text{ }\text{}\text{ }\right].$$
(17)
To order $`g^2`$, we distinguish several contributions. First there are two three-loop local diagrams coming from the interaction (13), and one two-loop local diagram from the Jacobian action (15):
$`F_2^{(1)}=g^2[\mathrm{\hspace{0.17em}3}({\displaystyle \frac{1}{2}}+a)\text{ }\text{}\text{ }+15\omega ^2({\displaystyle \frac{1}{18}}+{\displaystyle \frac{a}{5}})\text{ }\text{}\text{ }`$ (18)
(19)
$`3(a{\displaystyle \frac{1}{2}})\delta (0)\text{ }\text{}\text{ }].`$ (20)
(21)
We call a diagram local if it involves no temporal time integral. The Jacobian action (15) contributes further the nonlocal diagrams:
$`F_2^{(2)}={\displaystyle \frac{g^2}{2!}}\{2\delta ^2(0)\text{ }\text{}`$ (22)
$`4\delta (0)[\text{ }\text{}\text{ }+\text{ }\text{}\text{ }+2\omega ^2\text{ }\text{}\text{ }]\}.`$ (23)
In the perturbative calculations to follow, we shall use dimensional regularization where $`\delta (0)=0`$, according to a basic rule of t’Hooft and Veltman . As a consequence, the last terms in $`F_1`$, $`F_2^{(1)}`$, and the entire $`F_2^{(2)}`$ are zero. In fact, the term $`𝒜_J[q]`$ may be omitted completely from the path integral (10).
The remaining diagrams are either of the three-bubble type, or of the watermelon type, each with all possible combinations of the three line types (1)–(3). The former are
$`F_2^{(3)}={\displaystyle \frac{g^2}{2!}}[4\text{ }\text{}\text{ }+\mathrm{\hspace{0.17em}\hspace{0.17em}2}\text{ }\text{}\text{ }+\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}2}\text{ }\text{}`$ (24)
$`+8\omega ^2\text{ }\text{}\text{ }+8\omega ^2\text{ }\text{}\text{ }+8\omega ^4\text{ }\text{}\text{ }],`$ (25)
and the latter:
$`F_2^{(4)}={\displaystyle \frac{g^2}{2!}}\mathrm{\hspace{0.17em}4}[\text{ }\text{}\text{ }+4\text{ }\text{}\text{ }+\text{ }\text{}\text{ }+4\omega ^2\text{ }\text{}\text{ }+{\displaystyle \frac{2}{3}}\omega ^4\text{ }\text{}\text{ }].`$ (26)
(27)
Since the equal-time expectation value $`\dot{q}(\tau )q(\tau )`$ vanishes by Eq. (5), diagrams with a local contraction of a mixed line (2) are trivially zero, and have been omitted.
In our previous papers, all integrals were calculated individually in $`D=1\epsilon `$ dimensions, taking the limit $`\epsilon 0`$ at the end. Here we set up simple rules for finding the same results, which make the sum of all Feynman diagrams contributing to each order $`g^n`$ vanish.
## V Basic Properties of Dimensionally Regularized Distributions
The path integral (10) is extended to an associated functional integral in a $`D`$-dimensional coordinate space $`x`$, with coordinates $`x_\mu (\tau ,x_2,x_3,\mathrm{})`$, by replacing $`\dot{q}^2(\tau )`$ in the kinetic term by $`(_\mu q(x))^2`$, where $`_\mu =/x_\mu `$. The Jacobian action term (14) is omitted in dimensional regularization because of Veltman’s rule :
$$\delta ^{(D)}(0)=\frac{d^Dk}{(2\pi )^D}=0.$$
(28)
In our calculations, we shall encounter generalized $`\delta `$-functions, which are multiple derivatives of the ordinary $`\delta `$-function:
$`\delta _{\mu _1\mathrm{}\mu _n}^{(D)}(x)`$ $``$ $`_{\mu _1\mathrm{}\mu _n}\delta ^{(D)}(x)`$ (29)
$`=`$ $`{\displaystyle \overline{}d^Dk(ik)_{\mu _1}\mathrm{}(ik)_{\mu _n}e^{ikx}},`$ (30)
with $`_{\mu _1\mathrm{}\mu _n}_{\mu _1}\mathrm{}_{\mu _n}`$, and with $`\overline{}d^Dkd^Dk/(2\pi )^D`$. In dimensional regularization, alle these vanish at the origin as well:
$$\delta _{\mu _1\mathrm{}\mu _n}^{(D)}(0)=\overline{}d^Dk(ik)_{\mu _1}\mathrm{}(ik)_{\mu _n}=0,$$
(31)
which is a more general way of expressing Veltman’s rule. In the extended coordinate space, the correlation function (1) becomes
$$\mathrm{\Delta }(x)=\frac{d^Dk}{(2\pi )^D}\frac{e^{ikx}}{k^2+\omega ^2},$$
(32)
At the origin, it has the value
$$\mathrm{\Delta }(0)=\frac{\overline{}d^Dk}{k^2+\omega ^2}=\frac{\omega ^{D2}}{(4\pi )^{D/2}}\mathrm{\Gamma }\left(1\frac{D}{2}\right)\underset{D=1}{=}\frac{1}{2\omega }.$$
(33)
The extension of the time derivative (2),
$$\mathrm{\Delta }_\mu (x)=\overline{}d^Dk\frac{ik_\mu }{k^2+\omega ^2}e^{ikx}$$
(34)
vanishes at the origin, $`\mathrm{\Delta }_\mu (0)=0`$. This follows directly from a Taylor series expansion of $`1/(k^2+\omega ^2)`$ in powers of $`k^2`$, together with Eq. (31).
The second derivative of $`\mathrm{\Delta }(x)`$ has the Fourier representation
$$\mathrm{\Delta }_{\mu \nu }(x)=\overline{}d^Dk\frac{k_\mu k_\nu }{k^2+\omega ^2}e^{ikx}.$$
(35)
Contracting the indices yields
$$\mathrm{\Delta }_{\mu \mu }(x)=\overline{}d^Dk\frac{k^2}{k^2+\omega ^2}e^{ikx}=\delta ^{(D)}(x)+\omega ^2\mathrm{\Delta }(x),$$
(36)
which follows from the definition of the correlation function by the inhomogeneous field equation
$$(_\mu ^2+\omega ^2)q(x)=\delta ^{(D)}(x).$$
(37)
From (36) we have the relation between integrals
$$d^Dx\mathrm{\Delta }_{\mu \mu }(x)=1+\omega ^2d^Dx\mathrm{\Delta }(x),$$
(38)
Inserting Veltman’s rule (28) into (36), we obtain
$$\mathrm{\Delta }_{\mu \mu }(0)=\omega ^2\mathrm{\Delta }(0)\underset{D=1}{=}\frac{\omega }{2}.$$
(39)
This ensures the vanishing of the first-order contribution (17) to the free energy
$$F_1=g\left[\mathrm{\Delta }_{\mu \mu }(0)+\omega ^2\mathrm{\Delta }(0)\right]\mathrm{\Delta }(0)=0.$$
(40)
The same equation (36) allows us to calculate immediately the second-order contribution (21) from the local diagrams
$`F_2^{(1)}=3g^2\left[\left({\displaystyle \frac{1}{2}}+a\right)\mathrm{\Delta }_{\mu \mu }(0)5\left({\displaystyle \frac{1}{18}}+{\displaystyle \frac{a}{5}}\right)\omega ^2\mathrm{\Delta }(0)\right]`$ (41)
$`\times \mathrm{\Delta }^2(0)={\displaystyle \frac{2}{3}}\omega ^2\mathrm{\Delta }^3(0)\underset{D1}{=}{\displaystyle \frac{1}{12\omega }}.`$ (42)
The other contributions to the free energy in the expansion (16) require rules for calculating products of two and four distributions, which we are now going to develop.
## VI Integrals over Products of Two Distributions
The simplest integrals of this type are
$`{\displaystyle d^Dx\mathrm{\Delta }^2(x)}={\displaystyle \overline{}d^Dp\overline{}d^Dk\frac{\delta ^{(D)}(k+p)}{(p^2+\omega ^2)(k^2+\omega ^2)}}`$ (43)
$`=`$ $`{\displaystyle \frac{\overline{}d^Dk}{(k^2+\omega ^2)^2}}={\displaystyle \frac{\omega ^{D4}}{(4\pi )^{D/2}}}\mathrm{\Gamma }\left(2{\displaystyle \frac{D}{2}}\right)={\displaystyle \frac{(2D)}{2\omega ^2}}\mathrm{\Delta }(0),`$ (44)
and
$`{\displaystyle d^Dx\mathrm{\Delta }_\mu ^2(x)}={\displaystyle d^Dx\mathrm{\Delta }(x)\left[\delta ^{(D)}(x)+\omega ^2\mathrm{\Delta }(x)\right]}`$ (45)
$`=\mathrm{\Delta }(0)\omega ^2{\displaystyle d^Dx\mathrm{\Delta }^2(x)}={\displaystyle \frac{D}{2}}\mathrm{\Delta }(0).`$ (46)
To obtain the second result we have perfomed a partial integration and used (36).
In contrast to the integrals (43) and (46), the integral
$`{\displaystyle d^Dx\mathrm{\Delta }_{\mu \nu }^2(x)}={\displaystyle \overline{}d^Dp\overline{}d^Dk\frac{(kp)^2\delta ^{(D)}(k+p)}{(k^2+\omega ^2)(p^2+\omega ^2)}}`$ (47)
$`=`$ $`{\displaystyle \overline{}d^Dk\frac{(k^2)^2}{(k^2+\omega ^2)^2}}={\displaystyle d^Dx\mathrm{\Delta }_{\mu \mu }^2(x)}`$ (48)
diverges formally in $`D=1`$ dimension. In dimensional regularization, however, we may decompose $`(k^2)^2=(k^2+\omega ^2)^22\omega ^2(k^2+\omega ^2)+\omega ^4`$, and use (31) to evaluate further
$`{\displaystyle d^Dx\mathrm{\Delta }_{\mu \mu }^2(x)}={\displaystyle \overline{}d^Dk\frac{(k^2)^2}{(k^2+\omega ^2)^2}}=2\omega ^2{\displaystyle \frac{\overline{}d^Dk}{(k^2+\omega ^2)}}`$ (49)
$`+\omega ^4{\displaystyle \frac{\overline{}d^Dk}{(k^2+\omega ^2)^2}}=2\omega ^2\mathrm{\Delta }(0)+\omega ^4{\displaystyle d^Dx\mathrm{\Delta }^2(x)}.`$ (50)
Together with (43), we obtain the finite integrals
$`{\displaystyle d^Dx\mathrm{\Delta }_{\mu \nu }^2(x)}={\displaystyle d^Dx\mathrm{\Delta }_{\mu \mu }^2(x)}=2\omega ^2\mathrm{\Delta }(0)`$ (51)
$`+\omega ^4{\displaystyle d^Dx\mathrm{\Delta }^2(x)}=\left(1+D/2\right)\omega ^2\mathrm{\Delta }(0).`$ (52)
An alternative way of deriving the equality (48) is to use partial integrations and the identity
$$_\mu \mathrm{\Delta }_{\mu \nu }(x)=_\nu \mathrm{\Delta }_{\mu \mu }(x),$$
(53)
which follows directly from the Fourier representation (34).
Finally, from Eqs. (43), (46), and (52), we observe the useful identity
$$d^Dx\left[\mathrm{\Delta }_{\mu \nu }^2(x)+2\omega ^2\mathrm{\Delta }_\mu ^2(x)+\omega ^4\mathrm{\Delta }^2(x)\right]=\mathrm{\hspace{0.17em}0},$$
(54)
which together with the inhomogeneous field equation (36) reduces the calculation of the second-order contribution of all three-bubble diagrams (25) to zero:
$`F_2^{(3)}=g^2\mathrm{\Delta }^2(0)`$ (55)
$`\times {\displaystyle }d^Dx[\mathrm{\Delta }_{\mu \nu }^2(x)+2\omega ^2\mathrm{\Delta }_\mu ^2(x)+\omega ^4\mathrm{\Delta }^2(x)]=\mathrm{\hspace{0.17em}0}.`$ (56)
## VII Integrals Products of Four Distributions
More delicate integrals arise from the watermelon diagrams in (27) which contain products of four distributions, a nontrivial tensorial structure, and overlapping divergences . Consider the first three diagrams:
$`=`$ $`{\displaystyle d^Dx\mathrm{\Delta }^2(x)\mathrm{\Delta }_{\mu \nu }^2(x)}.`$ (57)
$`4\text{ }\text{}\text{ }`$ $`=`$ $`4{\displaystyle d^Dx\mathrm{\Delta }(x)\mathrm{\Delta }_\mu (x)\mathrm{\Delta }_\nu (x)\mathrm{\Delta }_{\mu \nu }(x)},`$ (58)
$`=`$ $`{\displaystyle d^Dx\mathrm{\Delta }_\mu (x)\mathrm{\Delta }_\mu (x)\mathrm{\Delta }_\nu (x)\mathrm{\Delta }_\nu (x)},`$ (59)
To exhibit the subtleties with the tensorial structure, we introduce the integral
$`I_D={\displaystyle d^Dx\mathrm{\Delta }^2(x)\left[\mathrm{\Delta }_{\mu \nu }^2(x)\mathrm{\Delta }_{\mu \mu }^2(x)\right]}.`$ (60)
In $`D=1`$ dimension, the bracket vanishes formally, but the limit $`D1`$ of the integral is nevertheless finite. We now decompose the Feynman diagram (57), into the sum
$`{\displaystyle d^Dx\mathrm{\Delta }^2(x)\mathrm{\Delta }_{\mu \nu }^2(x)}={\displaystyle d^Dx\mathrm{\Delta }^2(x)\mathrm{\Delta }_{\mu \mu }^2(x)}+I_D.`$ (61)
To obtain an analogous decompositions for the other two diagrams (58) and (59) we derive a few useful relations using the inhomogeneous field equation (36), partial integrations, and Veltman’s rule (31). First there is the relation
$$d^Dx\mathrm{\Delta }_{\mu \mu }(x)\mathrm{\Delta }^3(x)=\mathrm{\Delta }^3(0)\omega ^2d^Dx\mathrm{\Delta }^4(x).$$
(62)
By a partial integration, the left-hand side becomes
$$d^Dx\mathrm{\Delta }_{\mu \mu }(x)\mathrm{\Delta }^3(x)=3d^Dx\mathrm{\Delta }_\mu ^2(x)\mathrm{\Delta }^2(x),$$
(63)
leading to
$$d^Dx\mathrm{\Delta }_\mu ^2(x)\mathrm{\Delta }^2(x)=\frac{1}{3}\mathrm{\Delta }^3(0)\frac{1}{3}\omega ^2d^Dx\mathrm{\Delta }^4(x).$$
(64)
Invoking once more the inhomogeneous field equation (36) and Veltman’s rule (28), we obtain the integrals
$$d^Dx\mathrm{\Delta }_{\mu \mu }^2(x)\mathrm{\Delta }^2(x)=2\omega ^2\mathrm{\Delta }^3(0)+\omega ^4d^Dx\mathrm{\Delta }^4(x),$$
(65)
and
$$d^Dx\mathrm{\Delta }_{\mu \mu }(x)\mathrm{\Delta }_\mu ^2(x)\mathrm{\Delta }(x)=\omega ^2d^Dx\mathrm{\Delta }_\mu ^2(x)\mathrm{\Delta }^2(x).$$
(66)
Due to Eq. (64), the integral (66) takes the form
$`{\displaystyle d^Dx\mathrm{\Delta }_{\mu \mu }(x)\mathrm{\Delta }_\mu ^2(x)\mathrm{\Delta }(x)}`$ $`=`$ $`{\displaystyle \frac{1}{3}}\omega ^2\mathrm{\Delta }^3(0)`$ (67)
$``$ $`{\displaystyle \frac{1}{3}}\omega ^4{\displaystyle d^Dx\mathrm{\Delta }^4(x)}.`$ (68)
Partial integration, together with Eqs. (65) and (68), leads to
$`{\displaystyle d^Dx_\mu \mathrm{\Delta }_{\lambda \lambda }(x)\mathrm{\Delta }_\mu (x)\mathrm{\Delta }^2(x)}=`$ (69)
$`{\displaystyle d^Dx\mathrm{\Delta }_{\lambda \lambda }^2(x)\mathrm{\Delta }^2(x)}2{\displaystyle d^Dx\mathrm{\Delta }_{\lambda \lambda }(x)\mathrm{\Delta }_\mu ^2(x)\mathrm{\Delta }(x)}`$ (70)
$`={\displaystyle \frac{4}{3}}\omega ^2\mathrm{\Delta }^3(0){\displaystyle \frac{1}{3}}\omega ^4{\displaystyle d^Dx\mathrm{\Delta }^4(x)},`$ (71)
A further partial integration, and use of Eqs. (53), (66), and (71), produces the decompositions of the second and third Feynman diagrams (58) and (59):
$`4{\displaystyle d^Dx\mathrm{\Delta }(x)\mathrm{\Delta }_\mu (x)\mathrm{\Delta }_\nu (x)\mathrm{\Delta }_{\mu \nu }(x)}=`$ (72)
$`=`$ $`2I_D+4\omega ^2{\displaystyle d^Dx\mathrm{\Delta }^2(x)\mathrm{\Delta }_\mu ^2(x)},`$ (73)
and
$`{\displaystyle d^Dx\mathrm{\Delta }_\mu ^2(x)\mathrm{\Delta }_\nu ^2(x)}=`$ (74)
$`=`$ $`I_D3\omega ^2{\displaystyle d^Dx\mathrm{\Delta }^2(x)\mathrm{\Delta }_\mu ^2(x)}.`$ (75)
We now make the important observation that the subtle integral $`I_D`$ of Eq. (60) appears in Eqs. (61), (73) and (75) in such a way that it drops out from the sum of the watermelon diagrams in (27):
$`\text{ }\text{}\text{ }+4\text{ }\text{}\text{ }+\text{ }\text{}`$ (76)
$`={\displaystyle d^Dx\mathrm{\Delta }^2(x)\mathrm{\Delta }_{\mu \mu }^2(x)}+\omega ^2{\displaystyle d^Dx\mathrm{\Delta }^2(x)\mathrm{\Delta }_\mu ^2(x)}.`$ (77)
Using (64) and (65), the right-hand side becomes a sum of completely regular expressions. Moreover, adding to this sum the last two watermelon-like diagrams in Eq. (27):
$`4\omega ^2\text{ }\text{}\text{ }=4\omega ^2{\displaystyle }d^Dx\mathrm{\Delta }^2(x)\mathrm{\Delta }_\mu ^2(x),`$ (78)
and
$`{\displaystyle \frac{2}{3}}\omega ^4\text{ }\text{}\text{ }={\displaystyle \frac{2}{3}}\omega ^4{\displaystyle }d^Dx\mathrm{\Delta }^4(x),`$ (79)
we obtain for the contribution of all watermelon-like diagrams (27) the simple expression
$`F_2^{(4)}=2g^2{\displaystyle d^Dx\mathrm{\Delta }^2(x)}`$ (80)
$`\times \left[\mathrm{\Delta }_{\mu \mu }^2(x)+5\omega ^2\mathrm{\Delta }_\mu ^2(x)+{\displaystyle \frac{2}{3}}\omega ^4\mathrm{\Delta }^2(x)\right]`$ (81)
$`={\displaystyle \frac{2}{3}}\omega ^2\mathrm{\Delta }^3(0)\underset{D1}{=}{\displaystyle \frac{1}{12\omega }}.`$ (82)
This cancels the finite contribution (42), thus making also the second-order free energy in (16) vanish, and confirming the invariance of the perturbatively defined path integral under coordinate transformations up to this order.
## VIII Summary
In this note we have set up simple rules for calculating integrals over products of distributions in configuration space which produce the same results as dimensional regularization in momentum space. For a path integral of a quantum-mechanical point particle in a harmonic potential, we have shown that these rules lead to a reparametrization-invariant perturbation expansions of path integral.
Let us end with the remark that in the time-sliced definition of path integrals, reparametrization invariance has been established as long time ago in the textbook .
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# 1 Introduction
## 1 Introduction
It is well known that consistent string theory can only be formulated in space time higher than four dimensions, specifically 10 or 11 dimensions. Until recently the predominant view was that since the extra dimensions have not been detected, they must be compactified in small volumes with radii of Planck size. This was supported by the study of the heterotic string where the string scale $`M_S`$ is related to the Planck scale $`M_P`$ via
$$M_S=M_P\sqrt{k\alpha _G}$$
(1)
where $`\alpha _G`$ denotes the unified gauge coupling constant and $`k`$ is an integer of order unity which labels the level of the Kac-Moody algebra. Here the string scale is close to the Planck scale and this renders string effects unobservable. Recently, Witten observed that in nonperturbative string theory $`M_S`$ can be much lower. Lykken and Antoniadis independently suggested the possibility that $`M_S`$ can be as low as a few TeV. In this case, the fundamental string scale is not constrained theoretically as above but instead is constrained phenomenologically. This gains support from breakthroughs in understanding D-brane constructions in string theory . In the D-brane scenario the number of extra dimensions scanned by the Standard Model (SM) particles and the graviton can be very different.
Buoyed by these developments intense efforts are made to understand the main features of low scale string theories. A popular construction is to assume a factorizable geometry. The gravitational field is taken to propagate in the full 10 or 11- dimensional bulk volume whereas the SM particles are localized in the 3-brane and hence are not sensitive to the extra dimensions . The crucial observation is that the weakness of gravity in four dimension is due to its spreading in the $`n`$ extra dimensions and Eq.(1) is replaced by
$$M_P^2=M_{}^{n+2}V_n$$
(2)
where $`V_n`$ is the volume of the compactified extra space and $`M_{}`$ is the higher dimensional Planck scale. The two scales $`M_S`$ and $`M_{}`$ are related but the exact relation will not be needed in this paper. It is simplest to consider toroidal compactification, for which the volume is
$$V_n=(2\pi )^nR_1R_2R_3\mathrm{}R_n$$
(3)
and $`R_i(i=1,2,\mathrm{}n)`$ are the radii of the extra dimensions. For symmetric compactification one takes all the radii to be equal and denoted it by $`R`$.
This formulation has dramatic effects that lead to a new perspective on the mass hierarchy problem and the gauge coupling constants unification . Furthermore, these theories predict a deviation from the Newtonian $`1/r`$ law of gravitational potential at submillimeter range. The fact that this law is well tested above a distance of 1mm leads to the limit that $`n2`$ for $`M_{}=1\mathrm{T}eV`$, for the symmetric case. New experiments currently underway will probe smaller distance scales . Besides the direct search for a deviation of the Newtonian gravitational law the above scenario can also be indirectly constrained by astrophysical considerations. Although somewhat model dependent these considerations suggest that $`n2`$ and $`M_{}30`$ TeV .
A different scenario to explain the weak gravitational force is proposed by Randall and Sundrum . This mechanism does not require large extra dimensions but instead invokes a non-factorizable geometry. Here the extra dimension has the effect of modifying the space-time metric. Experiments that probe gravity at small distances are then not interpreted in terms of the scaling law of Eq. 2. The phenomenology of the graviton and its Kaluza-Klein (KK) excitations is different from the case of factorizable geometry and has been studied in Ref.. Given the differences of the these two scenarios and their extensions, and the richness of the phenomenology for each case it is clear that probes other than the graviton and its KK spectrum will be important additions to the study of higher dimensional physics.
It has been suggested recently that an extra right-handed neutrino that is a standard model gauge singlet may also be a probe of the extra dimensions. When a field is allowed to extend into the compactified extra dimensions it has associated with it Kaluza-Klein excitations which can have detectable effects in low energy experiments. Neutrinos propagating in the extra dimension will be referred to as the bulk neutrinos. The couplings of the bulk neutrino states with the active neutrino will give rise to a small neutrino mass for the active neutrino due to the spreading in the bulk volume. Some effects of these neutrinos in various low energy observables have been suggested . A solution to the solar neutrino problem in terms of matter enhanced flavor transformation of $`\nu _{eL}`$ into bulk neutrino states was studied in . Bulk neutrinos can induce energy loss in stars and the possible connection to the supernova collapse phase is studied in .
In this paper we concentrate on the constraints on the extra dimensional scenarios from universality tests of nuclear $`\beta `$-decays. Such studies have been used to provide a strong basis for the SM and more recently their usage is mainly in the test of unitarity of the quark mixings. We point out that the difference in Q-values for the nuclei which are very well measured can be used to constrain the properties of the bulk neutrinos such as the Yukawa couplings, the higher dimensional scale, $`M_{}`$, and the compactification radius R. These are very general properties of neutrinos in extra dimensions and one does not need to employ the scaling relation between the Planck scale and $`M_{}`$ as given in Eq.(2). To illustrate the effects of the bulk neutrinos we give the Kurie plots of two nuclei and the nuclear recoil spectrum of a Fermi transition. The simplest model of bulk neutrinos given below may allow a signal to be detected in the $`{}_{}{}^{3}\mathrm{H}`$ beta decay spectra. Other models which modify the tritium beta decay endpoint are discussed in . However, in searches for a small admixture of a more massive neutrino, the features of the extra dimensional model are not within currently detectable limits. However more realistic models with larger mixings with the active neutrinos may be constructed which can be probed by these experiments as well. We concentrate only on one family of bulk neutrino and the active $`\nu _{eL}`$. However, it is straightforward to extend the model to include three bulk neutrinos; one for each family. This implies further model assumptions and will take us into the realm of models for fermion families in extra dimensions which is beyond the scope of this study.
This paper is organized as follow. In section 2 we briefly outline the model of bulk neutrinos that we use and the mixings with $`\nu _{eL}`$. This serves the purpose of fixing of our notation. Sec. 3 details the study of universality using the data from Q-values of several $`\beta `$-decays. Next we present the impact of bulk neutrinos on the beta spectrum, and look at the recoil momentum for Fermi transitions. We give our conclusion in Sec 5.
## 2 A Simple Model of Bulk Neutrinos
The basic idea of producing neutrino masses using extra dimensions is given in Ref. . In the interest of being self contained and to establish our notation we give some of the construction of a simple model of bulk neutrinos. The simplest model assumes that the fermions charged under the standard model gauge group as well as the gauge bosons and the Higgs boson are localized on a 3-brane embedded in the bulk of larger dimensions. If we assume only one extra dimension then space-time is labelled by $`(x^\mu ,z)`$ where $`\mu =0,1,2,3`$. The extra coordinate $`z`$ is assumed to compactify into a circle of radius R. The brane scenario stipulates that the left-handed SM lepton doublet $`L=(\nu _{eL},e_L)`$ is given by $`L(x^\mu ,z=0)`$. The fields $`\nu _{eL}`$ and $`e_L`$ are Weyl spinors. Next we assume that a SM singlet fermion, denoted by $`\nu (x,z)`$ exists and it propagates in the full five dimensional bulk. Its right- and left-handed projections are labelled by $`\nu _R`$ and $`\nu _L`$ respectively. It is distinguished from the SM neutrino by not carrying a flavor index. In this simple model gravity is weak because there are additional dimensions in which it propagates.
In five dimensions there are five gamma matrices $`\mathrm{\Gamma }^\mu `$ and $`\mathrm{\Gamma }^5`$ where $`\mu =0,1,2,3`$. A convenient representation of the Clifford algebra in five dimensions is to choose $`\mathrm{\Gamma }^\mu =\gamma ^\mu `$ and $`\mathrm{\Gamma }^5=i\gamma ^5`$ with the usual $`\gamma ^\mu `$ of Minkowski space-time. The effective Lagrangian for generating a neutrino mass is given by
$$=_0^{2\pi R}𝑑z\overline{\nu }\left(i\gamma ^\mu _\mu +i\mathrm{\Gamma }_5_z\right)\nu +y_{}_0^{2\pi R}𝑑z\delta (z)\overline{L}H\nu _R+h.c.$$
(4)
where $`y_{}`$ is the dimensionful Yukawa coupling and $`H`$ is the Higgs doublet and it is related to the dimensionless coupling $`y`$ via
$$y_{}=\frac{y}{M_{}^{n/2}}$$
(5)
and $`n=1`$ for one extra dimension.
For simplicity we have neglected a higher dimensional bare Dirac mass term. This can be naturally implemented under $`𝒵_2`$ orbifold compactification . Following the Kaluza-Klein ansatz we Fourier expand $`\nu _R`$ and $`\nu _L`$ as follows:
$`\nu _R(x,z)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2\pi R}}}{\displaystyle \underset{k=\mathrm{}}{\overset{\mathrm{}}{}}}\nu _{kR}\mathrm{exp}\left({\displaystyle \frac{ikz}{R}}\right)`$ (6)
$`\nu _L(x,z)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2\pi R}}}{\displaystyle \underset{k=\mathrm{}}{\overset{\mathrm{}}{}}}\nu _{kL}\mathrm{exp}\left({\displaystyle \frac{ikz}{R}}\right)`$ (7)
Substituting this into Eq.(4) and integrating over $`z`$ yields the effective Lagrangian in four dimensions:
$$=\underset{k=\mathrm{}}{\overset{\mathrm{}}{}}\left[\overline{\nu }_{kL}i\gamma ^\mu _\mu \nu _{kL}+\overline{\nu }_{kR}i\gamma ^\mu _\mu \nu _{kR}\right]+\underset{k=\mathrm{}}{\overset{\mathrm{}}{}}m_k\overline{\nu }_{kL}\nu _{kR}+\frac{y_{}}{\sqrt{2\pi R}}\underset{k=\mathrm{}}{\overset{\mathrm{}}{}}\overline{L}H\nu _{kR}+h.c.$$
(8)
where
$$m_k=\frac{k}{R}$$
(9)
is the mass of the $`k\mathrm{t}h`$ KK tower state. The mass splitting between each adjacent tower state is $`1/R`$. As seen in Eq.(8) the coupling between the KK-tower neutrino states and the active $`\nu _{eL}`$ given by the Yukawa term. After spontaneous electroweak symmetry breaking a Dirac mass term is generated for the active $`\nu _{eL}`$ and is given by
$$m_D=\frac{yv}{\sqrt{4\pi RM_{}}}$$
(10)
where $`v=247`$ GeV. This mass is suppressed by a bulk volume factor as first noticed in Ref. and .
Rewriting all the mass terms in Eq.( 8) we have
$$m_D\overline{\nu }_{eL}\nu _{0R}+m_D\underset{k=1}{\overset{\mathrm{}}{}}\overline{\nu }_{eL}(\nu _{kR}+\nu _{kR})+\underset{k=1}{\overset{\mathrm{}}{}}m_k(\overline{\nu }_{kL}\nu _{kR}\overline{\nu }_{kL}\nu _{kR})+h.c.$$
(11)
We find it useful to define the following orthogonal states
$$\nu _{kR}^{^{}}=\frac{1}{\sqrt{2}}(\nu _{kR}+\nu _{kR})\nu _{kR}^{^{\prime \prime }}=\frac{1}{\sqrt{2}}(\nu _{kR}\nu _{kR})$$
(12)
and
$$\nu _{kL}^{^{}}=\frac{1}{\sqrt{2}}(\nu _{kL}\nu _{kL})\nu _{kL}^{^{\prime \prime }}=\frac{1}{\sqrt{2}}(\nu _{kL}+\nu _{kL})$$
(13)
which can be used to cast Eq.11 into the form
$$m_D\overline{\nu }_{eL}\nu _{0R}+\sqrt{2}m_D\underset{k=1}{\overset{\mathrm{}}{}}\overline{\nu }_{eL}\nu _{kR}^{^{}}+\underset{k=1}{\overset{\mathrm{}}{}}m_k(\overline{\nu }_{kL}^{^{}}\nu _{kR}^{^{}}+\overline{\nu }_{kL}^{^{\prime \prime }}\nu _{kR}^{^{\prime \prime }})+h.c.$$
(14)
The states with double prime superscripts have no low energy interactions and will be ignored. The mass terms can now be written in the familiar form of $`\overline{\nu }_LM\nu _R`$ in the bases of $`\nu _L=(\nu _{eL},\nu _{1L}^{},\nu _{2L}^{},\mathrm{})`$ and $`\nu _R=(\nu _{0R}^{},\nu _{1R}^{},\nu _{2R}^{},\mathrm{})`$. The mass matrix $`M`$ for $`k+1`$ states looks as follows
$`M=\left(\begin{array}{ccccc}m_D& \sqrt{2}m_D& \sqrt{2}m_D& \mathrm{}& \sqrt{2}m_D\\ 0& \frac{1}{R}& 0& \mathrm{}& 0\\ 0& 0& \frac{2}{R}& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& 0& 0& \frac{k}{R}\end{array}\right)`$ (20)
To find the left mass eigenstates, we consider the matrix $`MM^{}`$. Explicitly,
$`MM^{}={\displaystyle \frac{1}{R^2}}\left(\begin{array}{ccccc}(k+\frac{1}{2})\zeta ^2& \zeta & 2\zeta & \mathrm{}& k\zeta \\ \zeta & 1& 0& \mathrm{}& 0\\ 2\zeta & 0& 4& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ k\zeta & 0& 0& 0& k^2\end{array}\right)`$ (26)
where $`\zeta =\sqrt{2}m_DR`$. The mass eigenvalues $`\lambda `$ are given by the characteristic equation of Eq.(26) which is det$`[MM^{}\lambda ^2]`$. After some algebra it reduces to the transcendental equation :
$$\pi \zeta ^2\mathrm{cot}(\pi \lambda R)=2\lambda R$$
(27)
The solution of this equation yields the lowest mixed mass eigenstate $`m_D`$. For small $`\zeta `$ this is also the maximally mixed state. It is reasonable to assume that experiments that directly measures the neutrino mass will be probing this state. The current experimental limit on the mass of $`\nu _e`$ is $`m_{\nu _e}<2.5\mathrm{eV}`$ . When combined with this, Eq. (10) gives the following constraint
$$y<3.2\times 10^3\left(\frac{M_{}}{10\mathrm{T}\mathrm{e}\mathrm{V}}\right)^{\frac{1}{2}}\left(\frac{R}{1\mathrm{m}\mathrm{m}}\right)^{\frac{1}{2}}\left(\frac{m_D}{1\mathrm{e}\mathrm{V}}\right)$$
(28)
The diagonalization of Eq. (26) leads to the mixing of the $`k`$-th state to the lowest mass eigenstate $`\nu _{0L}^m`$ given by
$$\mathrm{tan}2\theta _k=\frac{2k\zeta }{k^2(k+\frac{1}{2})\zeta ^2}$$
(29)
This equation diagonalizes the submatrix between the lowest state $`\nu _{0L}`$ and the state $`\nu _{kL}`$. Where $`\zeta `$ is small, this is a good approximation to the exact mixing angle, obtainable by complete diagonalization of the matrix MM. In this limit, $`\zeta `$ can be used as an expansion parameter in the diagonalization process. For example, for $`M_{}=10\mathrm{TeV}`$, $`y=10^3`$ and $`R10^7\mathrm{m}m`$ we find $`\zeta 10^1`$. Hence, for small $`\zeta `$ we can see from above that the mixing of the KK states with the lowest mass neutrino eigenstate becomes progressively smaller as the value of $`k`$ increases. Furthermore, these KK states have no direct gauge interactions with other SM particles and their presence can only be probed through their mixing with the state $`\nu _{0L}^m`$.
In the small $`\zeta `$ limit, the vacuum oscillation survival probability for electron neutrinos is a particularly simple expression,
$$P(\nu _e\nu _e)1\frac{\pi ^2\zeta ^2}{3}+2\underset{k}{}\frac{\zeta ^2}{k^2}\left[12\mathrm{sin}^2\left(\frac{\delta m_{0k}^2t}{2E_\nu }\right)\right]+𝒪(\zeta ^4)$$
(30)
where $`\delta m_{0k}^2m_D^2m_k^2`$. The data from a reactor neutrino experiment can be used to find a limit on the Yukawa coupling. CHOOZ quotes a limit on two neutrino mixing of $`\mathrm{sin}^22\theta <10\%`$ for $`\delta m^2>10^3\mathrm{eV}^2`$. This translates into a limit on the extra dimensional parameters of
$$y<4.4\times 10^8\left(\frac{1\mathrm{m}\mathrm{m}}{R}\right)^{1/2}\left(\frac{M_{}}{10\mathrm{TeV}}\right)^{\frac{1}{2}}.$$
(31)
If combined with the limit from the neutrino mass, Eq 28, we find that the least stringent limit on $`y`$ occurs at $`R=7\times 10^6\mathrm{mm}`$.
However, Eq.(29) reveals a different phenomenon if $`\zeta `$ is large; i.e. $`\zeta 1`$. Now it is the higher KK states with $`k\zeta `$ that will have the largest mixing with the the lowest eigenstate; whereas states with lower $`k`$ values will have small mixing of order $`\frac{1}{\zeta }`$. In order for this to happen in $`n=1`$ case we require $`R10^3`$ mm. Obviously this cannot happen with $`\zeta 1`$ and a small compactification radius. When relatively large values of $`R`$ are considered, masses of these KK states can be in the $`1`$ eV range.
The above considerations can be generalized to higher extra dimensions. For $`n=2`$ and different compactification radii $`R_1`$ and $`R_2`$ the scaling equation for $`m_D`$ becomes
$$m_D=\frac{yv}{2\pi M_{}\sqrt{2R_1R_2}}$$
(32)
and yields the following constraint the on the Yukawa coupling
$$y<1.8\times 10^6\left(\frac{M_{}}{10\mathrm{T}\mathrm{e}\mathrm{V}}\right)\left(\frac{R_1R_2}{1\mathrm{m}\mathrm{m}^2}\right)^{\frac{1}{2}}\left(\frac{m_D}{1\mathrm{e}\mathrm{V}}\right)$$
(33)
For $`R_1=R_2`$ and a comparison with Eq.(28) shows that this is a much less stringent constraint on the parameters of the theory compared to the case of only one extra dimension. This is expected since the bulk neutrino state has a larger volume in which to spread and the Yukawa coupling is inversely proportional to the square root of this volume. An interesting case occurs for asymmetric compactification where one radius is much smaller than the other. Take for example $`R_2=10^{11}\mathrm{mm}`$ and $`R_1=0.1\mathrm{mm}`$ then Eq.(33) approaches the case of $`n=1`$ and gives a limit on $`y`$ similar to that of Eq.(28).
The masses of the KK tower states are generalized to
$$m_{k,l}=\sqrt{\left(\frac{k^2}{R_1^2}+\frac{l^2}{R_2^2}\right)}$$
(34)
where $`k`$ and $`l`$ denote the KK level corresponding to radii $`R_1`$ and $`R_2`$ respectively. It is seen that the smaller compactification radius $`R_2`$ will give rise to higher mass KK state in the asymmetric compactification scenario. For example if $`R_2=10^7\mathrm{mm}`$ one can have keV bulk neutrinos mixing into the active $`\nu _e`$ state that come from the small radius (i.e. for $`k=0`$ and small $`l`$ values). When both radii are large such neutrinos come only from the high KK values. Next we examine the mixings of these bulk states.
For small mixing the electron neutrino state can written as
$$\nu _e=\frac{1}{N}\left(\nu _0+\zeta _1\underset{k=1}{}\frac{1}{k}\nu _{k,0}+\zeta _1\underset{k,l1}{}\frac{1}{\sqrt{k^2+(\frac{R_1}{R_2})^2l^2}}\nu _{k,l}\right)$$
(35)
where $`\zeta _1=\sqrt{2}m_DR_1<<1`$ and $`N`$ is the normalization constant. It is seen that for the asymmetric case the KK states which come solely from the small radius may have larger masses than those from the larger radius; moreover, they are accompanied by correspondingly smaller mixings. We conclude that high mass states will have small mixing, regardless of the size of the compactification radii mainly due to the constraint on $`m_D`$ from the direct limit on the electron neutrino mass. Since the direct limits on the mu and tau neutrino masses are much less strict, large mixing may occur in these families for considerably smaller KK masses.
For two extra dimensions, the constant, $`N`$, depends on the cutoff scale $`M_{}`$, as $`N1+\zeta ^2\mathrm{log}(M_{}R)`$, if $`R_1=R_2`$ = R. Limits derived from universality of Fermi transitions which are presented in the next section do not depend on this normalization and therefore on the cutoff procedure. However, other calculations such as vacuum oscillation survival probabilities will depend on the normalization. In particular, this would become important if one were to compare the theory with the reactor neutrino data. In this case, the constraint on $`y`$ and $`M_{}`$ from universality could be used together with survival probability data to limit the size, $`R`$.
In the next section we study the constraint $`\beta `$ decay universality places on the parameters $`y,M_{}`$ and $`R`$.
## 3 Universality
In this section we study the consequences for beta decay of the model for bulk neutrinos outlined in section 2. Although bulk neutrinos do not have gauged interactions with SM particles, the allowed Yukawa coupling enables them to generate mass for the active neutrino. As seen in the last section the smallness of the neutrino mass is due to the spreading of the state in the higher dimensions. The KK excitations are also seen to have mixing effects with the active neutrino. In a given weak decay, the number of KK states that can be excited will depend on the energy released in the reaction. If $`1/R<<Q`$, where $`Q`$ is the nuclear Q-value, then many states contribute. For example, for n = 1, if $`Q=\mathrm{\hspace{0.17em}1}\mathrm{MeV}`$ and $`R=10^7\mathrm{mm}`$ then around 500 states contribute. In previous studies of such classical weak interactions one concentrates on one or two massive neutrino states and their kinematic effect. In the present scenario, towers of KK neutrinos ranging in mass from eV to a few MeV are involved in a nuclear $`\beta `$ decay. They have to be taken into account in the study of universality, electron spectra and nuclear recoil spectra. We shall examine these issues separately to see if they provide useful constraints on the Yukawa coupling, $`y`$ and the size R, of the extra dimension. We begin with universality.
The rate of beta decay from a single state in the parent nucleus to another state in the daughter nucleus is given by,
$$\lambda =\frac{\mathrm{ln}2}{ft}P(m_\nu )$$
(36)
The $`ft`$\- value is proportional to the inverse of the matrix element. $`P(m_\nu )`$ is the phase space factor, which depends on the mass of the emitted neutrino, $`m_\nu `$. We define the phase space factor as
$$P(m_\nu )=\frac{1}{m_e^5}_{m_e}^{Q+m_em_\nu }F(Z,E_e)E_ep_e(Q+m_eE_e)[(Q+m_eE_e)^2m_\nu ^2]^{1/2}𝑑E_e$$
(37)
The Q-value is the energy difference between the initial and final nuclear states, while $`E_e`$, and $`p_e`$ are the electron energy and momentum respectively. The Coulomb wave correction factor is denoted by $`F(Z,E)`$. We use the approximation of this factor found in .
For pure Fermi, $`0^+0^+`$ transitions and zero mass electron neutrinos, the corrected $`ft`$-values, $`Ft=ft(1+\delta _R)(1\delta _C)`$ for all nuclei should be the same under the assumption of universality in the standard model. Here, $`\delta _R`$ is the nucleus dependent part of the radiative correction and $`\delta _C`$ is the isospin symmetry breaking correction. These corrections contribute at the 1% level. Results of Ft-values for 10 nuclei are listed in .
The bulk neutrino scenario outlined in the previous section would cause an apparent deviation from the constant $`Ft`$ values predicted by universality in the standard model. This is because the actual phase space factor in an extra dimensional scenario is more complicated than Eq.(37). In general $`ft`$-values are found by using Eq. 36, with a measured rate and a calculated phase space factor. The transition rate for the case of one dimension is given in the small mixing limit, $`\zeta <<1`$, by
$$\lambda =\frac{\mathrm{ln}2}{ft_{xd}}\left[\left(1\frac{\pi ^2\zeta ^2}{6}\right)P(m_\nu 0)+\underset{k=1}{\overset{k_{max}}{}}P(m_{\nu _k}\frac{k}{R})\frac{\zeta ^2}{k^2}\right]$$
(38)
The maximum number of neutrinos to be summed over is determined by the maximum mass of a neutrino that can be released in the beta decay, $`k_{max}=QR`$. Transitions with a higher Q value will have more neutrino states that can contribute to the decay. Note that the $`ft`$ value in an extra dimensional scenario will not necessarily take on the same value as given in the standard model. (Although, the ‘true’ $`Ft`$\- values will be the same for all $`0^+0^+`$ transitions.) This will change the calculated value of the quark mixing angle $`V_{ud}`$ as discussed below. Because of this uncertainty, it is best to employ a normalization when looking for the effects of bulk neutrinos. In this example we will compare the relative difference in beta decays from two different nuclei.
Assuming that bulk neutrinos exist in one extra dimension, we can compare the $`ft`$-value for the extra dimensional scenario with the apparent $`ft`$-values for the standard model for the nucleus, $`A_1`$, by taking a ratio,
$$\frac{ft(A_1)_{xd}}{ft(A_1)_{SM,apparent}}\left(1\frac{\pi ^2\zeta ^2}{6}\right)+\underset{k=1}{\overset{k_{max}}{}}\frac{P(A_1,m_{\nu _k})}{P(A_1,0)}\frac{\zeta ^2}{k^2},$$
(39)
where we have used $`\lambda _{SM}/\lambda _{xd}=1`$. Since in the extra dimensional scenario, the $`ft`$-values for two nuclei are the same up to the isospin and radiative corrections, the apparent Standard Model values will not be the same,
$`{\displaystyle \frac{ft(A_2)_{SM,apparent}}{ft(A_1)_{SM,apparent}}}1`$ $`+{\displaystyle \underset{k=1}{\overset{k_{max}=Q_1R}{}}}{\displaystyle \frac{P(A_1,m_{\nu _k})}{P(A_1,0)}}{\displaystyle \frac{\zeta ^2}{k^2}}{\displaystyle \underset{k=1}{\overset{k_{max}=Q_2R}{}}}{\displaystyle \frac{P(A_2,m_{\nu _k})}{P(A_2,0)}}{\displaystyle \frac{\zeta ^2}{k^2}}`$ (40)
$`[\delta _R(A_1)\delta _R(A_2)]+[\delta _C(A_1)\delta _C(A_2)].`$
The maximum value of the second two terms may be obtained from the experimental uncertainty in the $`Ft`$-values from . The sums in the above expression are very insensitive to lower bound on the sum, so it is only necessary to consider the highest modes in the sum, where there are relatively large differences in the phase space factor ratio. The mixing of the low energy KK states is not relevant to universality considerations.
The above formula is valid for large $`\zeta `$, provided that $`\zeta ^2/k^2<<1`$, if the term $`(\zeta ^2/k^2)`$ is replaced by $`f(\zeta )^2(\zeta ^2/k^2)`$. Here $`f(\zeta )`$ is a function which must be calculated for each $`\zeta `$ by explicit diagonalization of the mass matrix given by Eq.(26). This explicit diagonalization shows that the mixing of state with $`k>>m_DR`$ is approximately $`\zeta ^{}/k=f(\zeta )\zeta /k`$. For $`\zeta =100`$, $`f(\zeta )=0.007`$. For larger $`\zeta `$, $`f(\zeta )`$ decreases, while for $`\zeta <<1`$, $`f(\zeta )=1`$.
We use $`\beta ^+`$ decay data from $`{}_{}{}^{14}\mathrm{O}`$ and $`{}_{}{}^{54}\mathrm{Co}`$ with a combined error of $`0.22\%`$. These nuclei have Q values of 1.81 MeV and 7.22 MeV respectively. This can be translated into a limit on the Yukawa coupling,$`y`$ and the size of the extra dimension, R,
$`0.0022>1.9\times 10^7\mathrm{MeV}^2y^2f(\zeta )^2\left({\displaystyle \frac{10\mathrm{TeV}}{M_{}}}\right)\left({\displaystyle \frac{1\mathrm{mm}}{R}}\right)`$
$`\left[{\displaystyle \underset{k=1}{\overset{Q_1R}{}}}\left({\displaystyle \frac{P(A_1,m_{\nu _k})}{P(A_1,0)}}{\displaystyle \frac{P(A_2,m_{\nu _k})}{P(A_2,0)}}\right){\displaystyle \frac{1}{m_{\nu _k}^2}}{\displaystyle \underset{k=Q_1R}{\overset{Q_2R}{}}}{\displaystyle \frac{P(A_2,m_{\nu _k})}{P(A_2,0)}}{\displaystyle \frac{1}{m_{\nu _k}^2}}\right]`$ (41)
Converting the sums in the above equation into integrals and integrating over all the available neutrino states produces a limit on the Yukawa coupling. As long as the small mixing limit applies, the limit on $`y`$ has little R dependence and remains fairly constant for example, at $`y<10^3`$ for a scale of $`M_{}=10\mathrm{TeV}`$. The limit becomes less strict for a larger high dimensional scale. The lack of R dependence can be seen when converting the sum in Eq. (3) to an integral, and making the change of variables $`dkRdm_\nu `$. In the case of large mixing, the limit is less strict, due to the factor $`f(\zeta )`$.
When comparing with Eqs.(28, 31), one can see that the limit on the Yukawa coupling from universality is always less stringent than the limit on the Yukawa coupling derived from the maximum neutrino mass and mixing determined by tritium beta decay and reactor neutrino experiments, respectively. This latter limit is shown as the bottom curve in Figure 1. Therefore, one does not obtain a useful constraint from universality when considering only one extra dimension.
One can perform the same exercise for two extra dimensions which both have the same size R. In the perturbative limit, the mass eigenstates are $`m_{\nu _{k,n}}^2(k^2+n^2)/R^2`$. In this case the sum ranges over two indices, $`k`$ and $`n`$. An equation similar to Eq.(3) may be derived for this case,
$`0.0022>6\times 10^{25}\mathrm{MeV}^2y^2\left({\displaystyle \frac{10\mathrm{TeV}}{M_{}}}\right)^2\left({\displaystyle \frac{\mathrm{mm}}{R}}\right)^2[`$
$`{\displaystyle \underset{k,n=1}{\overset{k^2+n^2=Q_1R}{}}}({\displaystyle \frac{P(A_1,m_{\nu _{k,n}})}{P(A_1,0)}}{\displaystyle \frac{P(A_2,m_{\nu _{k,n}})}{P(A_2,0)}}){\displaystyle \frac{1}{m_{\nu _{k,n}}^2}}{\displaystyle \underset{k,n=Q_1R}{\overset{k^2+n^2=Q_2R}{}}}{\displaystyle \frac{P(A_2,m_{\nu _{k,n}})}{P(A_2,0)}}{\displaystyle \frac{1}{m_{\nu _{k,n}}^2}}]`$ (42)
All two dimensional cases considered here, $`M_{}>1\mathrm{T}\mathrm{e}\mathrm{V}`$, fulfill the condition $`\zeta ^2<<1`$. We solve the above equation by converting the sums to integrals, and find that the limit is again insensitive to R. Figure 1 also shows the constraint on $`M_{}`$ and $`y`$ for the two dimensional case. It can be seen that greater precision is needed to probe scales greater than 10 TeV, since above this scale the Yukawa coupling is not constrained to be less than of order 1.
The Standard Model value of $`V_{ud}`$ can be obtained from measured $`Ft`$ values as in with the additional input of $`G_F`$, the weak coupling constant from muon decay. Using these quantities, the unitarity sum $`V_{ud}^2+V_{us}^2+V_{ub}^2`$ is calculated to be smaller than one by two standard deviations. The extra dimensional scenario would effect both the muon decay measurement as well as the $`Ft`$ values from beta decay. This introduces additional parameters and model uncertainties, such as the coupling of the muon neutrino to the same or another Kaluza Klein tower of neutrinos. If, for example, the muon neutrinos did not couple to any bulk neutrinos, then the calculated value of $`V_{ud}`$ should appear too large and an extra dimensional scenario could not account for the shortfall in the unitarity sum through the determination of $`V_{ud}`$ from beta decay.
Next we make contact with another much discussed test of universality using charged pion decays . The pion decays, $`\pi e\nu _e`$ and $`\pi \mu \nu _\mu `$ can be probed kinematically. Since the electron in the decay is relativistic, while the muon is not, the contributions of a massive neutrino differ considerably in the two decays; a massive electon neutrino will make a larger difference than a massive muon neutrino. In the case of massless neutrinos the standard model ratio is $`\mathrm{\Gamma }(\pi e\nu )/\mathrm{\Gamma }(\pi \mu \nu )=1.233\times 10^4`$. The experimentally measured quantity is $`R_{\pi \mu }^{\pi e}=(1.230\pm 0.004)\times 10^4`$ . Therefore, the extradimensional contribution to these decays can not exceed 0.3%. The contribution has been calculated explicitly by and for the $`n=2`$ case with two families of bulk neutrinos. The contribution is approximately given by
$$7.52\pi \times 10^2y_e^2\left(\frac{1\mathrm{T}\mathrm{e}\mathrm{V}}{M_{}}\right)^26.23\pi \times 10^1\left(\frac{1\mathrm{T}\mathrm{e}\mathrm{V}}{M_{}}\right)^2y_\mu ^2<0.003.$$
(43)
With this reaction alone, the only way to obtain a constraint on either $`y_e`$ or $`y_\mu `$ as a function of $`M_{}`$ is to make an assumption about one of the couplings such as, $`y_e1`$ or $`y_\mu 1`$. Ideally, one would obtain a separate constraint on one of the couplings and use the pion decay data to limit the other. Taking the constraint on $`y_e`$ from beta decay universality discussed previously we find that another two orders of magnitude in precision in the beta decay data would be required in order to probe $`y_\mu `$. However, a complete discussion must include an analysis of neutrino oscillations which is beyond the scope of this paper; however, see for a discussion.
## 4 Kinematic searches
In addition to changing the apparent $`ft`$ values, bulk neutrino scenarios can in principle also be probed kinematically. Here we demonstrate the consequences of the extra dimensional scenario on nuclear recoil momenta, and electron spectra. Massive neutrino searches in beta decay have focused on examining beta spectra and on examining the nuclear recoil spectra . A tower of Kaluza-Klein bulk states will produce a somewhat different signal than just one massive neutrino. We examine the signature in both of these types of spectra from the KK tower of states. Larger mixing scenarios may be probed by tritium beta decay searches for eV range neutrinos. However, as we shall see below,the model of extra dimensions cannot currently be constrained by searches for heavy (keV-MeV) neutrinos.
With a massive electron neutrino, the number of electron neutrinos as a function of momentum would be,
$$\frac{N(p_e)}{p_e^2F(Z,p_e)}=\frac{1}{m_e^5}\frac{1}{ft}\left[(QE_e+m_e)^2m_\nu ^2\right]^{1/2}(QE_e+m_e)$$
(44)
We show the normalized Kurie plot, $`(N(p_e)/(p_e^2F(Z,p_e)))^{1/2}`$ vs. $`E_e`$ for $`{}_{}{}^{38m}\mathrm{K}`$ at the top of Figure 3, for the case of $`m_\nu =0`$. The Q- value for this decay is 5.02 MeV.
In the case of one extra dimension, the function is modified to
$$\frac{N(p_e)}{p_e^2F(Z,p_e)}=\frac{1}{m_e^5}\frac{1}{ft_{xd}}\underset{k=0}{\overset{k_{max}}{}}|U_{ke}|^2\left[(QE_e+m_e)^2m_{\nu _k}^2\right]^{1/2}(QE_e+m_e).$$
(45)
Here $`|U_{ke}|^2`$ is the mixing of the kth mass eigenstate with the electron neutrino. For heavy massive neutrino searches (keV-MeV) very small mixing is relevant as discussed below. However, for massive neutrino searches in the eV range, such as tritium beta decay, a slightly larger mixing is applicable. Several states will give a different signature than just one or two massive states in a kinematic study. We illustrate this in the lower panel of Figure 2. This panel shows the beta endpoint spectrum for tritium for three cases. The dashed line shows the spectrum for a single massive neutrino of 2.3 eV. The solid line shows the result for a scenario with one extra dimensions, and $`\zeta =.13`$, $`1/R=25`$ eV giving $`m_D=2.3\mathrm{eV}`$. The first ten mass eigenstates were calculated by explicit diagonalization of the matrix, Eq. (26). The first three occur at 2.3 eV, 25 eV and 50 eV, with mixings $`|U_{ke}|^2`$ of 0.974, 0.017 and 0.004 respectively. As can be seen from the figure, the Kaluza-Klein tower produces several bumps as well as less counts slightly away from the endpoint than a single massive neutrino. For comparison a two neutrino mixing scenario is shown as the dot-dashed line. The mass eigenstates are at 2.3 eV and 25 eV with a mixing of $`\mathrm{sin}^2\theta 2.5\%`$.
We now turn to searches for massive neutrinos in the keV and MeV range. In the case of one extra dimension and small mixing, the function is modified to
$`{\displaystyle \frac{N(p_e)}{p_e^2F(Z,p_e)}}={\displaystyle \frac{1}{m_e^5}}{\displaystyle \frac{1}{ft_{xd}}}`$ $`[(1{\displaystyle \frac{\pi ^2\zeta ^2}{6}})(QE+m_e)+`$ (46)
$`{\displaystyle \underset{k=1}{\overset{k_{max}}{}}}{\displaystyle \frac{\zeta ^2}{k^2}}[(QE_e+m_e)^2m_{\nu _k}^2]^{1/2}(QE_e+m_e)]`$
We have normalized to the number of counts in order to eliminate the uncertainty in $`ft_{xd}`$. Therefore, we replace $`N(p_e)N(p_e)/N`$, where for one extra dimension, $`N=N(p_e)𝑑p_e`$. We plot the results for the extra dimensional scenario as the ratio,
$$=\left[\frac{N(p_e)}{Np^2F(Z,p)}\right]_{xd}^{1/2}/\left[\frac{N(p_e)}{Np^2F(Z,p)}\right]_{SM}^{1/2}.$$
(47)
Two cases are shown in Figure 3. The panel in the bottom left hand corner shows a situation where many modes can contribute to the decay; $`R=10^8\mathrm{mm}`$, $`m_D=1\mathrm{eV}`$. The mass of the lightest mode is at about 0.02 MeV and has a mixing of around $`\zeta ^2=5\times 10^9`$. Nuclei with smaller Q values will have a slightly larger signature for the continuum case.
The experimental limits on the mixing angle for a single 17 keV neutrino, for example are of order $`10^3`$. Since the mixing of the $`20\mathrm{keV}`$ modes from this extra dimensional scenario are considerably weaker, these modes are not likely to be detectable with present experimental data. The mixing of a $``$1 MeV neutrino is smaller still, since the mixing goes as $`\zeta /k=\sqrt{2}m_DR/k`$. If the lowest KK state occurs around 1 MeV, this implies a size $`R2\times 10^{10}\mathrm{mm}`$ and a mixing of $`\zeta ^23\times 10^{13}`$. The point here is that the mixing of the a high mass KK mode is proportional to $`m_D`$ which is constrained by tritium beta decay experiments. Therefore, in this simple model the mixing of keV and MeV neutrinos will always be small. The bottom right hand corner of Figure 3 shows contributions from three separate neutrino masses, 1, 2 and 3 MeV. The bumps shown in the lower right panel are unobservable, since current constraints on the mixing of an MeV neutrino are many orders of magnitude less ($`\mathrm{sin}^22\theta 10^3`$).
We also consider the method of detecting massive neutrinos by measuring nuclear recoil spectra. In the presence of a massive neutrino, the nuclear recoil spectrum for a pure Fermi transition has the shape,
$$P_r(Z,r)=\frac{1}{2}_{E_{min}}^{E_{max}}F(Z,E_e)\left[rE_eE_\nu +r\frac{a}{2}(r^2p_e^2p_\nu ^2)\right]𝑑E_e$$
(48)
Here, $`r`$ is the recoil momentum of the nucleus and $`a=1`$. For a pure Gamow-Teller transition, $`a=1/3`$. The limits of integration depend on the mass of the neutrino and are given by
$$E_{max,min}=\frac{1}{2}\left[\frac{E_0(E_0^2r^2+me^2m_\nu ^2)\pm r\sqrt{(E_0^2r^2m_e^2m_\nu ^2)^24m_\nu ^2m_e^2}}{E_0^2r^2}\right]$$
(49)
Here $`E_0=Q+m_e`$ is the total energy available for the decay. We again normalize the distribution by integrating over all possible recoil energies and finding $`P_r(m_\nu )`$, for each given neutrino mass. The distribution for $`m_\nu =0`$ is shown on the top panel of Figure 4. The normalized probability distribution for the scenario with one extra dimension looks like,
$$P_r(r)=\frac{(1\frac{\pi ^2\zeta ^2}{6})P_r(r,m_\nu 0)+\underset{k}{}\frac{\zeta ^2}{k^2}P_r(r,m_\nu k/R)}{P_r}$$
(50)
where the normalization is
$$P_r=(1\frac{\pi ^2\zeta ^2}{6})P_r(r,m_\nu 0)𝑑r+\underset{k}{}\frac{\zeta ^2}{k^2}P_r(r,m_\nu k/R)𝑑r$$
(51)
To see the effect of the bulk neutrino states on recoil spectra, we again plot normalized ratios, $`(P_r(r)/P_r)_{xd}/(P_r(r)/P_r)_{SM}`$ The results are shown in Figures 4 for the Fermi transition in $`{}_{}{}^{38m}\mathrm{K}`$. The top panels show the normalized spectra for a single zero mass neutrino. The bottom left panels show the effect when many neutrino modes can contribute to the decay, while the bottom right panels show the case where only a few neutrinos can contribute. However, as in the case with the Kurie plots, the mixing of the KK modes is several orders of magnitude too small to be detected with current experimental measurements.
## 5 Conclusions
We have studied the effects of bulk neutrinos in nuclear $`\beta `$ decays in a simple model which only allows the SM singlet state to propagate in one or more compactified extra dimensions. All other SM particles are confined to a 3-brane and have no KK excitations. There are many well studied nuclei with different Q values that allow us to use universality to place constraints on the Yukawa coupling and the fundamental scale $`M_{}`$. This information is complementary to the various studies of physics of higher dimensions using the gravitational force or KK graviton probes.
Our study shows that for $`n=1`$ the most stringent constraint on the parameters of the bulk neutrinos arise from the direct measurement of the mass of $`\nu _e`$ and from reactor neutrino oscillation experiments. Universality tests do not give additional information. On the other hand for $`n=2`$ the universality test restricts the Yukawa coupling to be less than unity for $`M_{}10\mathrm{TeV}`$. Future experiments with radioactive beams measuring decays with Q values of order 10 MeV and improved accuracy with current measurements that can push this by a factor of ten will be most welcome.
The Kurie plot and the recoil nuclear spectrum in some selected decays are also given. For probes of massive Kaluza-Klein states we find the nuclei with smaller Q values are more sensitive. Currently, for medium and heavy nuclei, the sensitivity of experiments designed for massive neutrino searches are several orders of magnitude below what is needed to probe mixing of the massive KK neutrinos in these reactions. This may be due to the simplicity of the model we have considered and hopefully such studies will prove useful for more realistic models with enhanced mixings that also address the family problem. We also point out that our study may be used in connection with other data to constrain parameters of second and third generation mixing with bulk neutrinos. These scenarios are more model dependent and contain more parameters. We give an illustration of this point for charged pion decays.
On the other hand we find that the tritium $`\beta `$ decay spectrum can be significantly altered by the presence of KK bulk neutrino states in the eV range. They can reveal themselves in the shape of endpoint region of the Kurie plot if the mixings are favorable.
It is useful to compare the results from universality limits to the ones derived from astrophysical constraints which are indirect. The latter relies on the energy loss carried away by bulk neutrinos. For example with limits derived from the neutrino magnetic moment for $`n=2`$ there is essentially no constraint on the neutrino magnetic moment induced from the sort of extra dimensional scenario presented here, while the limit on the Yukawa coupling is of order 1 for $`M_{}10\mathrm{TeV}`$ from universality. For $`n=1`$ one obtains a constraint from the mu and tau neutrino magnetic moment whereas universality gives no useful constraints. As mentioned before in this case for electron neutrinos the strongest constraint comes from the electron neutrino mass limit and oscillation experiments.
## 6 Acknowledgements
We wish to thank J. Behr and M.Trinczek for helpful discussions. One of us (J.N.N.) would like to thank Prof. T.K. Lee of the National Center for theoretical Science for his hospitality.
This research is partly supported by the Natural Science and Engineering Council of Canada.
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# Nonlinear Decoherence in Quantum State Preparation of a Trapped Ion
## Abstract
We present a nonlinear decoherence model which models decoherence effect caused by various decohereing sources in a quantum system through a nonlinear coupling between the system and its environment, and apply it to investigating decoherence in nonclassical motional states of a single trapped ion. We obtain an exactly analytic solution of the model and find very good agreement with experimental results for the population decay rate of a single trapped ion observed in the NIST experiments by Meekhof and coworkers (D. M. Meekhof, et al., Phys. Rev. Lett. 76, 1796 (1996)).
PACS numbers: 32.80.Pj, 42.50.Lc, 03.65.Bz, 05.45.+b
In recent years, much progress has been made in preparation, manipulation, and measurement of quantum states of the center-of-mass vibrational motion of a single trapped ion experimentally \[1-8\] and theoretically \[9-16\], which are not only of fundamental physical interest but also of practical use for sensitive detection of weak signals and quantum computation in an ion trap . In particular, the NIST group has experimentally created and observed nonclassical motional states of a single trapped ion. In the NIST experiments , an anti Jaynes-Cummings model (JCM) interaction between the internal and motional states of a trapped ion is realized through stimulated Raman transitions, which couple internal states of the trapped ion to its motional states, when the Lamb-Dicke limit is satisfied and the driving laser fields are tuned to the first blue sideband. Detection of motional states is carried out by observing the evolution characteristics of quantum dynamics of internal levels of the trapped ion under influence of the anti JCM-typed interaction. The NIST experiments revealed the fact that the population of the low atomic state ($`P_{}`$) evolves according to the following phenomenological expression
$$P_{}(t)=\frac{1}{2}\{1+\underset{n}{}p_n\mathrm{cos}(2gt\sqrt{n+1})e^{\gamma _nt}\}$$
(1)
where $`p_n`$ is the initial probability distribution of motional states of the trapped ion in the Fock representation, $`g`$ is a coupling constant between the atomic internal and motional states, $`\gamma _n`$ is a decay rate. The experimentally observed decay rate is of the following form
$$\gamma _n=\gamma _0(n+1)^\nu $$
(2)
where the observed value of $`\nu `$ is $`\nu 0.7`$.
A question that naturally arises is: how to explain the above experimentally observed decay rate? It is generally accepted that the appearance of the decay factor $`\gamma _n`$ in the evolution of internal states is a consequence of decoherence. It is of practical significance to well understand decoherence for preparation of nonclassical states and quantum computation in ion traps. There are various sources of decoherence , such as ion vibrational decoherence, ion internal-state decoherence, decoherence caused by non-ideal external fields, and so on. Recently, Schneider and Milburn have investigated decoherence due to laser intensity and phase fluctuations and obtained the power $`\nu `$ in Eq.(2) being $`\nu 0.5`$ instead of the experimentally observed value $`0.7`$. More recently, Murao and Knight , using master equation method, have studied decoherence due to the imperfect dipole transitions and fluctuation of vibrational potential in the NIST experiments. In spite of these efforts, the problem of decoherence in quantum state preparation of a trapped ion has been not satisfactorily solved, and its character and microscopic origin still call for further attention. In particular, it should be pointed out that the experimentally observed decay rate indicated in (2) is a collective effect caused by various decohering sources, not by a specific decohering source. Nevertheless, authors in refs investigated the decay rate caused only by a specific source of decoherence, not by various sources of decoherence. So how to model the experimentally observed decay rate caused by various decohering sources is an interesting subject in quantum state preparation and manipulation of a trapped ion. In this paper, we present a nonlinear decoherence model to model decoherence effects caused by various decohering sources in a quantum system. We shall show that our theoretical model can well describe the experimentally observed decay rate in the NIST experiments .
We consider a single trapped ion with mass $`m`$ and laser cooled to the Lamb-Dicke limit. Following symbols Ref., we denote three related internal states and motional states of the ion by $`|i`$ ($`i=\underset{¯}{0},,`$) and $`|n`$ ($`n=0,1,2,\mathrm{}`$), respectively. The free Hamiltonian of the trapped ion is given by $`\widehat{H_0}=\mathrm{}\omega _x\widehat{a}^+\widehat{a}\mathrm{}\omega _{01}||\mathrm{}\omega _{02}||`$, where $`\omega _{01}(\omega _{02})`$ is the transition frequency between states $`|(|)`$ and $`|\underset{¯}{0}`$, $`\widehat{a}^+(\widehat{a})`$ is the creation (annihilation) operator of the motional states with the corresponding frequency $`\omega _x`$. Two driving laser beams with detuning $`\mathrm{\Delta }`$, wave vector $`\stackrel{}{k}_1(\stackrel{}{k}_2)`$ and frequency $`\omega _1(\omega _2)`$ are used to cause dipole transitions between the level $`|(|)`$ and $`|\underset{¯}{0}`$.
With the dipole and rotating wave approximations, under large detuning condition the intermediate level $`|\underset{¯}{0}`$ can be adiabatically eliminated when the Lamb-Dicke limit is met and the driving laser beam is tuned to the first blue sideband. Then, in the interaction picture of $`\widehat{H_0}`$, the effective Hamiltonian of the system has the anti JCM-typed form
$$\widehat{H}_S=\mathrm{}g(\widehat{a}^+\sigma _++\widehat{a}\sigma _{})$$
(3)
where $`g`$ is a coupling constant, which depends on the coupling strength between internal and motional states of the trapped ion and the Lamb-Dicke parameter defined by $`\eta =\delta kx_0`$, where $`\delta k`$ is the wave-vector difference of the two Raman beams along $`x`$, and $`x_0=\sqrt{\mathrm{}/2m\omega _x}`$. For simplicity, we set $`\mathrm{}=1`$ throughout this paper.
The Hamiltonian (3) is diagonal in the dressed-state representation with the following basis
$$|\phi (n,i)=\frac{1}{\sqrt{2}}(|,n(1)^i|,n+1),i=1,2$$
(4)
$$|\phi (0,3)=|,0$$
(5)
And we have $`\widehat{H}_S|\phi (n,i)=E_{ni}|\phi (n,i)`$ with eigenvalues $`E_{ni}=(1)^{i+1}g\sqrt{n+1}`$ for $`i=1,2`$, and $`E_{03}=0`$.
Before going to our model, let us briefly recall a few basic facts about the interaction between a quantum system and its environment. The interaction between the system and its environment may create two types of effects \[20-34\]: decoherence and dissipation, which can be mathematically described by decaying of the off-diagonal and diagonal elements of the reduced density operator of the system, respectively. These two effects have been paid much attention in various areas, for instance, quantum measurement \[20,25-28\], condensed matter physics \[21-23\], quantum computation \[29-31\], and so on. The decoherence effect causes the states of the system continuously decohere to approach classical states . The dissipation effect dissipates energy of the system to environment \[21-23\]. The two effects can be understood in terms of Hamiltonian formalism \[32-34\]. If we assume that the total Hamiltonian of the system plus environment to be $`\widehat{H}_T=\widehat{H}_S+\widehat{H}_R+\widehat{H}_I`$, where $`\widehat{H}_S`$ and $`\widehat{H}_R`$ are Hamiltonians of the system and environment, respectively, and $`\widehat{H}_I`$ is the interaction Hamiltonian between them, when the Hamiltonian of the system commutes with that of the interaction between the system and environment, i.e., $`[\widehat{H}_S,\widehat{H}_I]=0`$, which means that there is no energy transfer between the system and the environment, energy of the system is conservative, so that what interaction between the system and environment describes is decoherence effect. When $`[\widehat{H}_S,\widehat{H}_I]0`$, there is energy transfer between the system and environment, so that what interaction between the system and environment describes is the dissipation effect. It should be pointed out that the decoherence and dissipation happen at different time scales . The dissipation effect occurs at the relaxation time $`\tau _{rel}`$, while the decoherence time scale $`\tau _d`$ is much shorter than $`\tau _{rel}`$ with the time evolution of a quantum system. Hence, we here restrict our attention on decoherence effect.
We now present our model. We use a reservoir consisting of an infinite set of harmonic oscillators to model the environment of the single trapped ion in the NIST experiments, and assume that in the interaction picture of $`\widehat{H}_0`$ the total Hamiltonian is of the following phenomenological form
$`\widehat{H}_T`$ $`=`$ $`\widehat{H}_S+{\displaystyle \underset{k}{}}\omega _k\widehat{b}_k^{}\widehat{b}_k+F(\{\widehat{O}_S\}){\displaystyle \underset{k}{}}c_k(\widehat{b}_k^{}+\widehat{b}_k)`$ (7)
$`+F^2(\{\widehat{O}_S\}){\displaystyle \underset{k}{}}{\displaystyle \frac{c_k^2}{\omega _k^2}}.`$
Here the first term is the Hamiltonian of the system in the interaction picture given by Eq.(3); the second term is the Hamiltonian of the reservoir; the third one represents the interaction between the system and the reservoir with a coupling constant $`c_k`$, where $`\{\widehat{O}_S\}`$ is a set of linear operators of the system or their linear combinations in the same picture as that of $`\widehat{H}_S`$, $`F(\{\widehat{O}_S\})`$ is an operator function of $`\{\widehat{O}_S\}`$. In order to enable what the interaction between the system and the reservoir describes in Eq.(6) is decoherence not dissipation, we require that the linear operator $`\widehat{O}_S`$ commutes with the Hamiltonian of the system, i.e., $`[\widehat{O}_S,\widehat{H}_S]=0`$. It is well known that the decohering process can indeed be considered as a quantum measurement process. The conventional definition of a quantum measurement involves any form of interaction between a quantum object and a classical system. Therefore, the interaction function $`F(\{\widehat{O}_S\})`$ in the model (6) can involve any form of interaction between the system and environment. This enables it to model the collective decohering behavior caused by various decohering sources. The concrete form of the function $`F(\{\widehat{O}_S\})`$ may be regarded as an experimentally determined quantity. The last term in Eq.(6) is a renormalization term, which is discussed in Ref.. When $`F(\{\widehat{O}_S\})`$ is a linear and nonlinear function of the linear operator $`\widehat{O}_S`$, We call decoherence described by the interaction between the system and the reservoir linear and nonlinear decoherence, respectively, in the similar sense of the linear and nonlinear dissipation implied in Ref.. In this sense, decoherence investigated in Ref. is a kind of linear decoherence. In what follow we shall show that nonlinear decoherence can better describe the decay rate in the the NIST experiments.
The Hamiltonian (6) can be exactly solved by making use of the following unitary transformation
$$\widehat{U}=\mathrm{exp}[F(\{\widehat{O}_S\})\underset{k}{}\frac{c_k}{\omega _k}(\widehat{b}_k^{}\widehat{b}_k)].$$
(8)
After applying the unitary transformation (7) to the total Hamiltonian (6), we get a decoupled Hamiltonian $`\widehat{H}_T^{}=\widehat{H}_S+_k\omega _k\widehat{n}_k`$, where $`\widehat{n}_k=\widehat{b}_k^{}\widehat{b}_k`$. The density operator associated with the decoplued Hamiltonian is given by
$$\widehat{\rho }_T^{}(t)=e^{i\widehat{H}_T^{}t}\widehat{\rho }_T^{}(0)e^{i\widehat{H}_T^{}t}$$
(9)
where $`\widehat{\rho }_T^{}(0)=\widehat{U}\widehat{\rho }_T(0)\widehat{U}^1`$, with $`\widehat{\rho }_T(0)`$ being the initial total density operator. Through a converse transformation of (7), it is straightforward to obtain the total density operator associated with the original Hamiltonian (6) with the following expression
$`\widehat{\rho }_T(t)`$ $`=`$ $`e^{i\widehat{H}_St}\widehat{U}^1e^{it_k\omega _k\widehat{n}_k}\widehat{U}\widehat{\rho }_T(0)`$ (11)
$`\times \widehat{U}^1e^{it_k\omega _k\widehat{n}_k}\widehat{U}e^{i\widehat{H}_St}.`$
We assume that the system and reservoir are initially in thermal equilibrium and uncorrelated, so that $`\widehat{\rho }_T(0)=\widehat{\rho }_S(0)\widehat{\rho }_R(0)`$, where $`\widehat{\rho }_S(0)`$ and $`\widehat{\rho }_R(0)`$ are the initial density operator of the system and the reservoir, respectively. $`\widehat{\rho }_R(0)`$ can be expressed as $`\widehat{\rho }_R=_k\widehat{\rho }_k(0)`$ where $`\widehat{\rho }_k(0)=(1e^{\beta \omega _k})e^{\beta \omega _k\widehat{n}_k}`$ is the density operator of the $`k`$-th harmonic oscillator in thermal equilibrium, where $`\beta =1/k_BT`$, $`k_B`$ and $`T`$ being the Boltzmann constant and temperature, respectively. After taking the trace over the reservoir, from Eq.(9) we can get the reduced density operator of the system, denoted by $`\widehat{\rho }(t)=tr_R\widehat{\rho }_T(t)`$, its matrix elements in the dressed state representation are explicitly written as
$`\rho _{(m^{},i^{})(m,i)}(t)`$ $`=`$ $`\rho _{(m^{},i^{})(m,i)}(0)R_{m^{}i^{}mi}(t)`$ (13)
$`\times e^{i\varphi _{m^{}i^{}mi}(t)},`$
Here the phase is defined by
$$\varphi _{m^{}i^{}mi}(t)=[E_{m^{}i^{}}E_{mi}],$$
(14)
and $`R_{m^{}i^{}mi}(t)`$ is a reservoir-dependent quantity given by
$`R_{m^{}i^{}mi}(t)`$ $`=`$ $`{\displaystyle \underset{k}{}}Tr_R\{D(\alpha _{mik})e^{it\omega _k\widehat{n}_k}D(\alpha _{mik})`$ (16)
$`\times D(\alpha _{m^{}i^{}k})e^{it\omega _k\widehat{n}_k}D(\alpha _{m^{}i^{}k})\widehat{\rho }_k(0)\},`$
where $`\alpha _{mik}=f(\{O_{mi}\})c_k/\omega _k`$ with $`O_{mi}`$ being an eigenvalue of the linear operator $`\widehat{O}_S`$ in a dressed state, i.e., $`\widehat{O}_S|\phi (m,i)=O_{mi}|\phi (m,i)`$, and $`D(\alpha )=\mathrm{exp}(\alpha \widehat{b}_k^+\alpha ^{}\widehat{b}_k)`$ is a displacement operator.
Making use of properties of the displacement operator:
$`D(\alpha )D(\beta )=D(\alpha +\beta )\mathrm{exp}[iIm(\alpha \beta ^{})],`$ (17)
$`\mathrm{exp}(x\widehat{n}_k)D(\alpha )\mathrm{exp}(x\widehat{n}_k)=\mathrm{exp}(\alpha e^x\widehat{b}_k^+\alpha ^{}e^x\widehat{b}_k),`$ (18)
and the following formula
$$Tr_R[D(\alpha )\widehat{\rho }_k(0)]=\mathrm{exp}[\frac{1}{2}|\alpha |^2\mathrm{coth}(\frac{\beta \omega _k}{2})],$$
(19)
we find that the reservoir-dependent quantity $`R_{m^{}i^{}mi}(t)`$ can be written as the following factorized form
$$R_{m^{}i^{}mi}(t)=e^{i\delta \varphi _{m^{}i^{}mi}(t)}e^{\mathrm{\Gamma }_{m^{}i^{}mi}(t)},$$
(20)
with the following phase shift and damping factor
$$\delta \varphi _{m^{}i^{}mi}(t)=[F^2(\{O_{m^{}i^{}}\})F^2(\{O_{mi}\})]Q_1(t),$$
(21)
$$\mathrm{\Gamma }_{m^{}i^{}mi}(t)=[F(\{O_{m^{}i^{}}\})F(\{O_{mi}\})]^2Q_2(t).$$
(22)
Here the two reservoir-dependent functions are given by
$$Q_1(t)=_0^{\mathrm{}}𝑑\omega J(\omega )\frac{c^2(\omega )}{\omega ^2}\mathrm{sin}(\omega t),$$
(23)
$$Q_2(t)=2_0^{\mathrm{}}𝑑\omega J(\omega )\frac{c^2(\omega )}{\omega ^2}\mathrm{sin}^2(\frac{\omega t}{2})\mathrm{coth}(\frac{\beta \omega }{2}),$$
(24)
where we have taken the continuum limit of the reservoir modes: $`_k_0^{\mathrm{}}𝑑\omega J(\omega )`$, where $`J(\omega )`$ is the spectral density of the reservoir, $`c(\omega )`$ is the corresponding continuum expression for $`c_k`$.
We assume that the system is initially in a state $`\widehat{\rho }(0)=||_np_n|nn|`$. Then, from Eqs.(10)-(16) we find that at time $`t`$ the population of the lower atomic state is given by
$$P_{}(t)=\frac{1}{2}\{1+\underset{n}{}p_n\mathrm{cos}[\varphi _{n1n2}(t)+\delta \varphi _{n1n2}(t)]e^{\mathrm{\Gamma }_{n1n2}(t)},$$
(25)
which indicates that the interaction between the system and reservoir induces a phase shift $`\delta \varphi _{n1n2}(t)`$ and a damping factor $`\mathrm{\Gamma }_{n1n2}(t)`$ in the time evolution of the atomic population.
Taking into account the experimental expression (1), we choose the following linear operator and interaction function:
$`\widehat{O}_S=\widehat{a}^+\sigma _++\widehat{a}\sigma _{},`$ (26)
$`F(\{\widehat{O}_S\})=\widehat{O}_S^{2d+1},`$ (27)
where $`d`$ is an adjustable parameter to describe the nonlinearity in the interaction, which reflects the deviation degree of the nonlinearity of $`F(\{\widehat{O}_S\})`$ with respect to the linear operator $`\widehat{O}_S`$. The value of the parameter $`d`$ is determined by the experimental results. With these choices, it is easy to find that
$`F^2(O_{n1})F^2(O_{n2})=0,`$ (28)
$`F(O_{n1})F(O_{n2})=2(\sqrt{n+1})^{2d+1}.`$ (29)
Then the phase shift in Eq.(21) naturally vanishes, and the damping factor becomes
$$\mathrm{\Gamma }_{n1n2}=4(n+1)^{2d+1}Q_2(t).$$
(30)
So that we can find from Eq.(21) that
$$P_{}(t)=\frac{1}{2}\{1+\underset{n}{}p_n\mathrm{cos}(2gt\sqrt{n+1})e^{4(n+1)^\nu Q_2(t)}\}$$
(31)
where $`\nu =2d+1`$ and $`Q_2(t)`$ is given by Eq.(20). From Eq.(27) we see that the argument of the cosine function on the RHS of Eq.(27) does have the same form as that in the experimental expression (1). Comparing the theoretical expression (27) with the experimental result (1), we find that when the nonlinear deviation $`d0.15`$, the $`n`$-dependence of the damping factor in Eq.(27) is completely in agreement with that seen in the experimental expression (1). The final step is to determine the time dependence of the damping factor in Eq.(27). From Eqs.(19), (20), and (27), we see that all necessary information about the effects of the environment is contained in the spectral density of the reservoir. Eq.(27) indicates that the time dependence of the damping factor is completely determined by the spectral density of the reservoir. The experimental expression (1) requires that the time dependence of the damping factor must be linear, so that if we choose the spectral density such that
$$Q_2(t)=\frac{1}{4}\gamma _0t$$
(32)
where $`\gamma _0`$ is a characteristic parameter, then we can get an expression of $`P_{}(t)`$, which has exactly the same form as the experimental result (1). It is possible to find a spectral density of the reservoir to satisfy the condition (28). For instance, for the case of zero temperature, if we take the spectral density $`J(\omega )=\gamma _0/(2\pi c^2(\omega ))`$, substituting it to Eq.(20) we can realize Eq.(28).
In conclusion, we have present a nonlinear decoherence model, and obtained its exactly analytic solution. It has been shown that our model can give precisely the same expression of the population decay rate of the single trapped ion as that observed in the NIST experiments . The nonlinear decoherence model can describe the NIST experiments so well. This indicates that the reservior and the nonlinear coupling between the system and the reservoir, which we design, properly model the real environment of the single trapped ion in the experiments. It is worthwhile to emphasize that the nonlinearity in the coupling describes a collective contribution of various decohering sources to the decay rate. Hence, what the nonlinear decoherence describes is a collective decoherence effect caused by various decohering sources not a specific decoherence source. We have noted that authors in ref. obtained the decay rate in Eq.(1), but decoherence which they considered is a specific decoherence caused by the imperfect dipole transitions and fluctuation of vibrational potential, so their results can not cover the contribution of other decohering sources to the decay rate in the NIST experiment. It can be expected that the nonlinear decoherence model proposed in the present paper can describe decoherence behaviors of a wide variety of quantum systems.
L. M. Kuang thanks Profs. Changpu Sun , Hong Chen , and Dr. Shao-Ming Fei a for enlightening discussions. This work was supported in part by 95-climbing project of China, NSF of China, Educational Committee Foundation and NSF of Hunan Province, and special project of NSF of China via Institute of Theoreical Physics, Academia Sinica.
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# A Giant Glitch in the Energetic 69 ms X-ray Pulsar AXS J161730-505505
## 1. Introduction
Radio pulsars are thought to be highly magnetized ($`10^{12}`$ G), rapidly spinning neutron stars whose luminosity is powered by rotational energy loss. The study of young ($`<10^5`$ yrs) rotation-powered pulsars provides an important laboratory for understanding the early evolution (thermal, spin, and magnetic) of these embers of stellar collapse. The most energetic of these pulsars are observable at X-ray wavelengths, which allow us to probe these extreme, but rare, examples.
In addition to uniform spin-down corresponding to the rotational energy loss, rotation powered pulsars show sudden discontinuities in their rotation periods (see Lyne & Graham-Smith 1998, chapter 6). These rare phenomena, known as “glitches” are considered to arise from sudden changes in the configuration of super-dense material in the neutron star interior. To date, a total of 71 glitches with $`|\mathrm{\Delta }P/P|>10^9`$ have been reported in 30 pulsars (Urama & Okeke 1999). Observation of pulsar glitches gives us insights into the structure and physical processes inside the neutron stars, such as the interactions of neutron superfluid and crust components (e.g., Anderson & Itoh 1975).
In this article we report the detection of a giant glitch from AXS J161730$``$505505 by using the newly acquired multi-mission X-ray data. This source is an unusual case of a young pulsar discovered first by its X-ray emission, revealed during the course of an archival X-ray study of the SNR RCW 103 (Gotthel, Petre, & Hwang 1997). Further analysis detected highly significant pulsations from photons attributed to this source (Torii et al. 1998). The 69 ms pulse period was consistent with that reported from a GINGA observation of the region (Aoki, Dotani, & Mitsuda 1992). Recent radio observation has confirmed AXS J161730$``$505505 as a young energetic rotation-powered pulsar (Kaspi et al. 1998). The lack of evidence for this object in archival soft X-ray images of this field suggested that the source might be highly absorbed, making it difficult to determine whether it is associated with a SNR.
## 2. Observations
A set of day-long X-ray observations of the field containing RCW 103, 1E 161348$``$5055, and AXS J161730$``$505505 were performed with the ASCA (Tanaka, Inoue, & Holt 1994), SAX (Boella et al. 1997), and XTE (Bradt et al. 1993) observatories. We summarize in table 1 the set of observations presented in this work.
Imaging data were acquired with the Gas Imaging Spectrometers (GIS) on-board ASCA and with the Medium Energy Concentrators (MECS) instruments on SAX. These instruments have moderate imaging ($`2^{}`$) and spectral resolution ($`8\%`$ at 6 keV) over an energy band pass of $`0.710`$ keV (GIS) and $`1.512`$ keV (MECS), with a field-of-view (FOV) large enough to cover the SNR and pulsar simultaneously. Non-imaging data were obtained with the Proportional Counter Array (PCA) on-board XTE which provides broader energy band-pass ($`260`$ keV) at lower spectral resolution ($`16\%`$ at 6 keV). The PCA FOV is roughly circular with a $`1^{}`$ FWHM response. The GIS data were collected in the highest time resolution mode ($`0.5\mathrm{ms}`$ or $`61\mu \mathrm{s}`$, depending on data acquisition mode) whose measured absolute accuracy is $`200\mu \mathrm{s}`$ in this mode (Saito et al. 1997). The PCA data were collected using the Good Xenon mode with $`0.9\mu `$s timing resolution. For the current analysis, the absolute timing uncertainty is $`100\mu `$s (Rots et al. 1998). Photons collected by the MECS are time tagged with 15 $`\mu `$s resolution. We do not include data from ASCA’s Solid-state Imaging Spectrometers (SIS), as the pulsar fell just off the edge of its FOV. Nor do we include data from SAX’s other instruments as the observing time is insufficient for these instruments to measure the pulsar periodicity reliably, as a part of the observation was interrupted prematurely.
Each data set was processed through its standard pipeline reduction for that mission and edited to exclude times of high background contamination using the standard screening criteria. This rejects time intervals of South Atlantic Anomaly passages, Earth occultations, bright Earth limb in the FOV (ASCA and SAX only), and other periods of high particle activity. The resulting effective observation times are summarized in Table 1. For each observation, event data from all detectors were co-added and the arrival times of each event were corrected to the solar system barycenter using the software TIMECONV (ASCA), BARYCONV (SAX), or FXBARY (XTE).
## 3. Results
### 3.1. Timing
The X-ray images obtained with both ASCA and SAX above 3 keV reveal AXS J161730–505505 $`4^{}`$ outside the SNR shell (see figure 1a and 1b of Gotthelf et al. 1997; Gotthelf, Petre, & Vasisht 1999). To increase the signal-to-noise ratio for detecting pulsations from the pulsar, we extracted photons from an 8 diameter aperture centered on the pulsar, restricting the energy range of extracted photons to $`310`$ keV for GIS and $`312`$ keV for the other instruments. For the PCA data, in this energy band, we further restrict our search to Layer 1 data only, which provides the best sensitivity for a Crab-like spectrum; For the higher energy analysis afforded by the PCA, above $`12`$ keV, we used data from all three PCA layers.
We searched each data set for the expected 69 ms period predicted from the initial period and period derivative measurement (Torii et al. 1998). A periodgram was constructed using the $`\chi ^2`$ statistic to test against a null hypothesis. For each trial period, we folded the data into 10 bins and computed the $`\chi ^2`$ of the resultant profile. We search a narrow range of periods centered on the expected period $`\pm 0.1`$ ms, sampled in increments of $`0.1\times P^2/T`$, where $`T`$ is the observation duration, and $`P`$ is the test period. A highly significant signal was detected from each of our data sets.
As well as the newly obtained data, we have re-analyzed the previous GINGA and ASCA data (Aoki et al. 1992; Torii et al. 1998) in a uniform way and revised the period and its error by using the method of Leahy (1987). Our X-ray timing results derived from these 13 measurements are listed in table 1, along with an updated radio ephemeris (Kaspi 1999, Private communication).
### 3.2. Spectrum
We search for spectral dependence of the pulse profile by comparing the folded light curves in several energy bands. No strong energy dependence is evident in the energy resolved light curves. Furthermore, the pulse amplitude and pulse profile remained unchanged between observational epochs.
We examined the ASCA and SAX data on AXS J161730$``$505505 for any long term changes in its energy spectrum or flux. As for the timing analysis, we restrict our comparison to the energy range above $`3`$ keV and extract photons from an $`8^{}`$ diameter aperture centered on the source. We fitted the spectrum with a power law function modified by interstellar absorption. The absorption was fixed at $`6.8\times 10^{22}\mathrm{cm}^2`$ (Torii et al. 1998). Spectra from each observation were found to be consistent with each other. Combining the 7 ASCA observations, we obtain the pulse phase averaged photon index, $`\mathrm{\Gamma }=1.4\pm 0.2`$, and the observed flux of $`(3.6\pm 0.2)\times 10^{12}\mathrm{ergs}\mathrm{s}^1\mathrm{cm}^2`$ (90% confidence errors) in the 3-10 keV range, which is consistent with the previous measurement (Torii et al. 1998).
## 4. Discussion
A $`\chi ^2`$ fit to the all 14 data points as summarized in table 1 gives the mean spin-down rate of $`\dot{P}=1.3611(1)\times 10^{13}`$ s/s and $`P=0.069347150(1)`$ at MJD 50,000.0 (Figure 1). However, the quality of the fit is bad, with $`\chi ^2/d.o.f.=4808/12`$. The residual of the fit shows a jump of $`\mathrm{\Delta }P1.2\times 10^7`$ s between the observations of 1993 August (MJD 49,217.6) and 1997 September (MJD 50,696.0) (Figure 2, top panel). Within the observation span of 10 years, the residual is neither periodic nor smooth. A sudden change in the period between MJD 49,217.6 and MJD 50,696.0 is suggested. We consider the most likely explanation for these residuals is due to a glitch, similar to those observed in several young rotation-powered pulsars.
Given the clear evidence of glitch activity, we next attempted to model the spin-down data with a single glitch followed by an exponential recovery. The data coverage is limited and we simply assumed the following relation for the spin-down.
$$P(t)=P_0+\dot{P}(tt_0)+\mathrm{\Delta }P\mathrm{exp}(\frac{tt_0}{\tau })$$
(1)
where $`\mathrm{\Delta }P=0`$ for $`t<t_0`$ and $`\mathrm{\Delta }P`$ is a negative constant for $`tt_0`$. This model contains five parameters, which, except for the depth of the glitch $`\mathrm{\Delta }P`$, are found to be independent of the time of the glitch, $`t_0`$. The derived parameters, $`P_0`$, $`\dot{P}`$, $`\mathrm{\Delta }P`$, and $`\tau `$ are summarized in table 2 for assumed values of $`t_0=49,300.0`$ MJD, $`t_0=50,000.0`$ MJD, and $`t_0=50,600.0`$ MJD. The residual for $`t_0=50,000`$ is shown in the middle panel of figure 2. The quality of the fit is now characterized by $`\chi ^2/d.o.f.=22.6/10`$.
The size of the glitch depends strongly upon the unknown glitch epoch $`t_0`$. The fractional increase in rotation was found to be $`\mathrm{\Delta }P/P=4.2\times 10^6`$ for $`t_0=50,000.0`$ MJD (fixed) while it changes between $`\mathrm{\Delta }P/P=11\times 10^6`$ for $`t_0=49,300.0`$ MJD (fixed) and $`\mathrm{\Delta }P/P=1.8\times 10^6`$ for $`t_0=50,600.0`$ MJD (fixed). The minimum fractional increase in rotation rate is therefore comparable with those of the largest known pulsar glitches (Lyne, et al. 1996a; Shemar & Lyne 1996).
Using the above model, the recovery time following the glitch episode is found to be $`\tau =700`$ days. This duration is somewhat unusual compared to the radio pulsars whose recovery time is seen to bifurcate between $`\tau 100`$ days and $`\tau 1,000`$ days (Shemar & Lyne 1996). Our derived value lies squarely between these two time scales, perhaps due to the simple model we invoked which allows for only a single glitch. Because of the sparse data coverage between August 1993 and September 1997, however, we cannot determine if the recovery time could be expressed as the sum of the two timescales.
Glitches found in radio pulsars may be classified into three groups (Lyne & Graham-Smith 1998). The first is a Crab-like glitch which is characterized by steps mainly in $`\dot{P}`$. The second is a Vela-like glitch which is characterized by large changes in $`P`$ ($`\mathrm{\Delta }P/P10^6`$) and the exponential recoveries. For the Vela pulsar, linear sawtooth changes in $`\dot{P}`$ have been observed between glitches (Lyne et al. 1996b). The third kind is often found in old pulsars which is characterized by a change in $`P`$.
Since the large glitch found for AXS J161730$``$505505 is similar in its magnitude to those in the Vela pulsar, a phenomenological model taking into account the sawtooth behavior may be a good description. Apart from short term effects, the spin-down for the Vela pulsar is expressed by a linear change in spin-down rate. Therefore, the following function is appropriate if the transient effects have already ceased by MJD 50,696.
$$P(t)=P_0+\mathrm{\Delta }P_0+(\dot{P}+\mathrm{\Delta }\dot{P})(tt_0)+\frac{1}{2}\ddot{P}(tt_0)^2$$
(2)
Here, $`\mathrm{\Delta }P_0=0`$ and $`\mathrm{\Delta }\dot{P}=0`$ for $`t<t_0`$ and they are constant values ($`\mathrm{\Delta }P_0<0`$ and $`\mathrm{\Delta }\dot{P}>0`$) for $`tt_0`$. This model contains six parameters. Again, the time of the glitch, $`t_0`$, had to be assumed. The derived parameters are summarized in table 3. For the condition that $`\mathrm{\Delta }P_0<0`$, $`t_0`$ was restricted to $`t_0\stackrel{<}{_{}}50,205`$. The residual for $`t_0=50,000`$ is shown in the bottom panel of figure 2. The quality of the fit is significantly improved to $`\chi ^2/d.o.f.=17.1/9`$. Compared to the fit to Equation (1), the f-test gives a chance probability of $`0.1`$. This result suggests that transient effects had indeed ceased by MJD 50,696. The reduced $`\chi ^2`$ is still larger than unity, suggesting the presence of timing noise and smaller glitches.
We can estimate the expectancy of large ($`|\mathrm{\Delta }P/P|>10^7`$) glitches by using the semi-empirical relation based upon the superfluid vortex unpinning model (Alpar & Baykal 1994). Using this relation, the expected number of large glitches between the first GINGA observation and the last ASCA observations is 3.7. Therefore, there should have been about 4 glitches of $`|\mathrm{\Delta }P/P|>10^7`$. Indeed the residuals to the fit of Equation (1) or (2) still hint a small jump of $`\mathrm{\Delta }P=5\times 10^9`$ s between MJD 51,263.8 (1999 Mar.) and MJD 51,394.2 (1999 Aug). This may be another glitch of $`\mathrm{\Delta }P/P=7\times 10^8`$, much smaller in magnitude than the one near MJD 50,000, but still relatively large compared to those seen for most radio pulsars (Shemar & Lyne 1996). We have thus found a giant glitch of $`|\mathrm{\Delta }P/P|\stackrel{>}{_{}}10^6`$ at $`t_050,000`$ MJD ($`49,218t_050,696`$) and possibly a glitch of $`|\mathrm{\Delta }P/P|7\times 10^8`$ at $`t_051,300`$ MJD ($`51,264t_051,394`$).
The detection of a giant glitch from AXS J161730$``$505505 gives a rare sample for studying the origin of pulsar glitches and the interior structure of neutron stars. In this context, regular timing observation of this pulsar in the radio band is quite important to monitor the onset of a glitch and following transient effects. Long term timing observations for measuring the braking index are desired. Searches for the highly absorbed X-ray emission from the supernova remnant associated with this pulsar should be undertaken with the next generation X-ray observatories.
Acknowledgments — We thank Jules Halpern for a critical reading of the manuscript, the referee Andrew Lyne for invaluable comments and suggestions on the original draft, and Victoria Kaspi for communicating the revised radio ephemeris. Special thanks to Alan Smale for expediting the delivery of the XTE data and to the XTEhelp@athena.gsfc.nasa.gov crew for their patience and assistance. Part of this research has made use of data obtained through the HEASARC data center, provided by NASA/GSFC. E. V. G. and G. V. acknowledge the supported of NASA LTSA grant NAG 5-7935.
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# Random Matrix Theory and Chiral Symmetry in QCD
## 1 INTRODUCTION
Well before the advent of quantum chromodynamics (QCD), the theory of the strong force, it was realized that the essential ingredients of the hadronic world at low energies are chiral symmetry and its spontaneous breaking. (see e.g. Refs. ). Mainly through lattice QCD simulations, it has become well established that chiral symmetry breaking by the vacuum state of QCD is a nonperturbative phenomenon that results from the interaction of many microscopic degrees of freedom. We argue in this review that the complexity of the QCD vacuum leads to a low-energy description that is completely dictated by the global symmetries of QCD. This interpretation of Goldstone’s theorem provides a natural duality between a strongly interacting fundamental theory and a weakly interacting low-energy effective theory.
### 1.1 QCD and Chiral Symmetry
We illustrate the concept of spontaneous symmetry breaking using the simpler example of a classical spin system with two rotational degrees of freedom. The Hamiltonian of this system has a certain symmetry: It is invariant under rotations. In mathematical language, the symmetry group is $`G=\mathrm{O}(3)`$. However, at low temperatures, the ground state of the system does not exhibit this symmetry. In a small external magnetic field, which breaks the rotational invariance explicitly, the spins will polarize in the direction of the magnetic field. In the thermodynamic limit, the spins remain polarized even if the magnetic field is switched off completely. This phenomenon is known as spontaneous magnetization. The ground state is no longer invariant under O(3) rotations but only under O(2) rotations in the plane perpendicular to the spontaneous magnetization. In mathematical language, the full symmetry group $`G=\mathrm{O}(3)`$ is spontaneously broken to a smaller symmetry group $`H=\mathrm{O}(2)`$. The spontaneously broken phase is characterized by low-energy excitations in the form of spin waves in the plane perpendicular to the spontaneous magnetization. This is a consequence of a general theorem known as Goldstone’s theorem , which tells us that spontaneous breaking of a continuous symmetry leads to low-lying excitations, the Goldstone modes, with a mass that vanishes in the absence of a symmetry-breaking field.
The Goldstone modes are given by the fluctuations in the plane perpendicular to the direction of the spontaneous magnetization. Thus, the spin system has two Goldstone modes. In general, spontaneous symmetry breaking in a system of spins with $`n`$ components is associated with $`n1`$ Goldstone modes. This number is also equal to the number of generators of the coset $`G/H`$ \[the number of generators of O($`n`$) is $`n(n1)/2`$\].
A spontaneously broken symmetry is characterized by an order parameter, which in this case is the spontaneous magnetization. At nonzero temperature, the alignment of the spins is counteracted by their thermal motion, and above a critical temperature (the Curie temperature) the spontaneous magnetization vanishes.
Let us now consider the hadronic world and interpret the particle spectrum in terms of the concepts discussed above. We will look for the simplest theory consistent with the following two empirical facts: $`(a)`$ there are three bosonic particles, the pions, that are much lighter than all other hadrons and $`(b)`$ the proton and the neutron have almost the same mass. Fact $`(a)`$ implies that there are three Goldstone bosons associated with a spontaneously broken symmetry. Assume that the underlying theory is based on an $`n`$-component field without preferred directions, i.e. the theory is invariant under O($`n`$) transformations of the fields. Spontaneous symmetry breaking means that the ground state has a preferred direction, leaving $`n1`$ directions for the Goldstone modes. Since there are three pions, this suggests that $`n=4`$ and that the ground state of the theory should be symmetric under O(3). The O(3) symmetry is the familiar isospin symmetry, which results in the near equality of the masses of the pions and of the mass of the proton and the neutron, fact $`(b)`$. As was first conjectured by Gell-Mann and Lévy , the theory of the strong force is based on an O(4) invariance spontaneously broken to O(3) with nucleons transforming according to SU(2). This is the familiar linear $`\sigma `$-model. Since the pions are not completely massless, a small O(4) symmetry-breaking mass term should also be present in the theory.
The $`\sigma `$-model is a phenomenological model for the interactions of pions and nucleons. It has been very successful in explaining many previously known empirical relations. Our aim, however, is to understand chiral symmetry in terms of QCD, the fundamental theory of the strong interactions. Quantum chromodynamics is a gauge theory of quarks that come in six different flavors and three different colors. The gauge field interaction is according to the non-Abelian SU(3) color group. Two of the six quarks are nearly massless, and at low energies QCD is well approximated by a theory with only the two lightest quarks. They are mixed by the SU(2) isospin symmetry group. This symmetry is exact for degenerate quark masses. For massless quarks, there is an additional symmetry. The helicity of a particle is a good quantum number, and the right-handed and left-handed quarks can be rotated independently. The isospin or vector symmetry rotates both chiralities in the same way whereas the axial SU(2) group rotates them in the opposite direction. This explains that the chiral symmetry group in the massless case is $`G=\mathrm{SU}(2)\times \mathrm{SU}(2)`$, which is isomorphic to O(4). The mass term breaks this symmetry explicitly and thus plays the role of the magnetic field in the spin system. Even in the massless case, however, the vacuum state of QCD is characterized by a nonzero expectation value of the chiral condensate, which, like the mass term, mixes right-handed and left-handed quarks. Thus, in the vacuum, the axial part of the symmetry group $`G`$ is broken spontaneously. The ground state is not unique, and the degenerate states are connected by the broken group $`G`$. The degeneracy can be lifted by means of a small mass term, and in the thermodynamic limit the system will be frozen in this direction, i.e. the direction of the QCD vacuum state is determined by the mass term. This is exactly the same situation as in the spin system discussed above. The Goldstone modes analogous to spin waves are the pions, with a mass that is well below the typical hadronic mass scale of about 1 GeV. As in the spin system, we expect that the expectation value of the chiral condensate will become zero above a critical temperature. This phenomenon is known as the restoration of chiral symmetry. The chirally symmetric phase probably existed in the early universe and may be produced in relativistic heavy-ion collisions.
Although QCD is the consistent theory of the strong interactions, many questions remain unanswered. For example, why have its constituents never been observed in nature? This phenomenon is known as confinement and means that all physical states are color singlets. QCD is best understood at high energies, where, because of asymptotic freedom, quarks and gluons become weakly interacting and perturbative calculations are possible. At low energies, on the other hand, it is necessary to rely on nonperturbative approaches. One approach is to study the QCD partition function by means of Monte Carlo simulations of a discretized version of the QCD action. This approach has been very fruitful, and a great deal of our understanding of low-energy QCD is based on such calculations (for a recent review, see Ref. ). The drawback is that large-scale simulations do not necessarily provide a simple picture of the relevant degrees of freedom. Therefore, it is often advantageous to study QCD by means of effective models or theories. One example is the instanton liquid vacuum , which is a model for an ensemble of relevant gauge field configurations. Effective low-energy theories are a second example. We already mentioned the $`\sigma `$-model for pions and nucleons. However, because of the spontaneous breaking of chiral symmetry and the existence of a mass gap in QCD, one can do better. Based on chiral symmetry one can formulate an exact low-energy theory for the Goldstone modes. This nonlinear $`\sigma `$-model is the basis for a systematic low-energy expansion that is discussed in detail in Secs. 1.3 and 4.
The QCD partition function can be written as a Euclidean path integral that can be expressed as the expectation value of the fermion determinant,
$$Z=\underset{f}{}det(𝒟+m_f).$$
(1)
Here the average is over all gauge fields weighted by the Euclidean Yang-Mills action, the product is over quark flavors of mass $`m_f`$, and $`𝒟`$ is the Dirac operator, which we introduce in great detail in Sec. 2. The fermion determinant can be expressed as a product over the eigenvalues $`i\lambda _n`$ of the Dirac operator. Therefore, we may also interpret the average as an average over the eigenvalues with probability distribution determined by the gauge field dynamics. We have argued that spontaneous breaking of chiral symmetry means that a small quark mass leads to a macroscopic realignment of the QCD vacuum. It is clear from the QCD partition function that this is only possible if there is an accumulation of Dirac eigenvalues near zero. Otherwise, a small mass term would be completely dominated by the much larger eigenvalues in the factors of $`(i\lambda _n+m_f)`$. For a free Dirac operator in a four-dimensional box, the eigenvalue density is proportional to $`\lambda ^3`$ and thus vanishes near zero. Therefore, the small eigenvalues must be due to interactions mediated by the gauge fields.
There are two possibilities. First, the eigenvalues may originate from the exactly zero Dirac eigenvalues in the field of an instanton. At a nonzero density of the liquid of instantons and anti-instantons, the zero eigenvalues are distributed over a band because of interactions that lift the degeneracy of the eigenvalues . The eigenvalue repulsion results in a nonzero eigenvalue density near zero. Second, the eigenvalues may originate from the bulk of the Dirac spectrum. As is the case for any interaction in quantum mechanics, the interactions mediated by the gauge fields lead to a repulsion of the eigenvalues. For increasing interaction strength, it is then advantageous for the eigenvalues to move to a region with a low eigenvalue density. Both mechanisms of spontaneous chiral symmetry breaking rely on the repulsion of eigenvalues in interacting systems. This topic has been investigated in detail by means of random matrix theory, discussed in the next subsection.
As mentioned above, we expect chiral symmetry to be restored at high temperatures . This is confirmed by lattice QCD simulations, which show a critical temperature of about 150 MeV . However, there is another direction in which the QCD partition function can be explored, namely at nonzero baryon number density. In this domain, rigorous results are available only at infinite baryon number density . This raises the question to what extent these results can be extrapolated to physically interesting densities. Because of the complex phase of the fermion determinant, Monte Carlo simulations are not possible in this case. Therefore, it is not surprising that many of the recent developments in QCD at nonzero density are based on the analysis of effective models . The picture that has emerged for two massless flavors is that a first-order chiral phase transition occurs at zero temperature along the chemical potential axis. Renormalization-group arguments and lattice QCD simulations indicate that for two massless flavors, a second-order phase transition occurs at zero chemical potential along the temperature axis. The expectation is that a line of first-order transitions and a line of second-order transitions will extend in the $`\mu T`$ plane and will join at the tricritical point, as indicated in Fig. 1. A color-superconducting phase is conjectured to exist at higher densities but the discussion of this phase is beyond the scope of this review. One of the questions we address is the robustness of this picture based on the dynamics of the eigenvalues of the QCD Dirac operator.
### 1.2 Random Matrix Theory
Random matrix theory (RMT) first appeared in the mathematical literature in 1928 and was first applied to physics in the context of nuclear resonances by Wigner almost 50 years ago . At that time, theoretical approaches such as the shell model had been very successful in describing the low-lying excitations of complex nuclei. However, highly excited resonances, which can be observed experimentally by neutron scattering, could not be described by the microscopic theory. The problem is generic: For any complex quantum system containing many degrees of freedom with complicated dynamics, it is very hard, if not impossible, to obtain exact results for the energy levels far above the ground state of the system.
Having acknowledged that the highly excited states cannot be predicted individually, one can ask whether the experimental data have some generic statistical features that can be described theoretically. This is where RMT comes in. Every quantum system is described by a Hamilton operator that can be expressed in matrix form. For a complex quantum system such as a large nucleus, this matrix is very complicated, and we may not even know its details. In this case, one approach is to assume that all interactions that are consistent with the symmetries of the system are equally likely. This means replacing the elements of the Hamilton matrix by random numbers that are uncorrelated and distributed according to the same probability distribution. To obtain definite results, observables such as the level density must then be averaged over the random matrix elements. This defines a statistical theory of energy levels, which is mathematically known as RMT. An important point is that the random matrix must have the same symmetries as the original Hamilton matrix. A collection of early papers on RMT can be found in the book by Porter , and the standard reference on RMT is the book by Mehta .
Can such an enormous simplification of the real problem describe empirical data? A similar puzzle occurred in the theory of critical phenomena, in which the critical behavior does not depend on the detailed dynamics of the theory. The reason for this simplification is the appearance of a length scale, the correlation length, that diverges at the critical point. Because of the corresponding separation of scales, it is possible to integrate out the short-wavelength fluctuations and renormalize the theory to a fixed-point theory that does not depend on the details of the initial theory. What is the separation of scales that takes place in quantum spectra? The two basic scales are the average level spacing and the scale of the variation of the average level spacing. The equivalent of the correlation length is the inverse average level spacing, which diverges in the thermodynamic limit or the semiclassical limit. We thus expect that spectral properties on the scale of the average level spacing do not depend on the details of the underlying dynamics. They are universal.
Universal properties can be studied in the simplest model of a given universality class and, in this way, exact analytical results for the correlation functions can be obtained. In the case of spectral correlations the simplest models of the universality classes are the Gaussian RMTs, and in the past three decades many exact results have been derived for these models. Unfortunately, in most cases it is not possible to prove whether a given theory belongs to one of the RMT universality classes. Therefore, random matrix predictions are usually verified by comparisons with empirical data.
Since Wigner’s original proposal, universal quantities have been identified and computed in a variety of fields including nuclear physics, atomic and molecular physics, disordered mesoscopic systems, quantum systems with classically chaotic analogs, two-dimensional quantum gravity, conformal field theory, and QCD. A recent comprehensive review of the applications of RMT can be found in Ref. . RMT is now an independent subfield of mathematical physics. It provides a unifying description of universal statistical features of many different quantum systems and is applicable whenever a system is sufficiently complex.
Let us raise an interesting point here. The eigenvalues of the Hamilton matrix for a given quantum system are the quantities of primary interest. Instead of going through RMT to compute the eigenvalues, an alternative — and much simpler — approach might be to postulate random eigenvalues. It turns out that this does not describe empirical data, at least not if the random eigenvalues are uncorrelated. Therefore, in addition, one would have to postulate how the eigenvalues are correlated, and it is not at all clear how to do this. In RMT, on the other hand, one starts with uncorrelated random matrix elements. At the end of the calculation, one finds that the resulting eigenvalues are strongly correlated in precisely the right way to describe the data.
In most applications, RMT is used to describe the statistics of energy levels, i.e. of the eigenvalues of the Hamilton operator. For the Euclidean QCD partition function, it is more natural to construct a random matrix model for the Dirac operator. As shown in the previous subsection, the spectrum of this operator is intimately related to the phenomenon of chiral symmetry breaking. This establishes the connection between RMT and chiral symmetry announced in the title of this review. We use random matrix methods to study the Dirac spectrum and its implications for chiral symmetry breaking.
It is important to note that RMT can only provide an exact description of universal quantities. It cannot be used for the calculation of nonuniversal observables such as the average level spacing. Such behavior is well known in statistical mechanics where, for example, the Ising model describes the critical exponents of the liquid-gas phase transition but does not give the critical temperature. In the case of the distribution of the eigenvalues of a system, the global spectral density is not described by RMT. For example, for the Gaussian random matrix ensembles the average spectral density is a semicircle, whereas in real systems it typically increases strongly with excitation energy. In contrast, universal quantities do not depend on the details of the dynamics of the system. In RMT, they do not depend on the distribution of the random matrix elements. It is crucial to distinguish the universal quantities from the model-dependent ones. In the first five sections of this review, we mainly concentrate on the universal, model-independent features that yield exact quantitative results. However, it is sometimes useful to construct random matrix models to obtain a qualitative description of the physics in a disordered system. As we show in Sec. 6, schematic random matrix models for QCD at nonzero temperature and density can yield important qualitative insights into problems that are otherwise difficult to tackle.
### 1.3 Effective Low-Energy Theories and Chiral Random Matrix Theory
Let us now make the connection between QCD, effective low-energy theories for QCD, and chiral RMT. We concentrate here on the big picture, postponing a more detailed discussion to Secs. 3 and 4.
At low energies, quarks and gluons are confined in hadrons, i.e. the particles we observe in nature are composite objects. Instead of attempting to describe the results of low-energy scattering experiments in terms of quarks and gluons, it is often simpler to use an effective theory whose elementary degrees of freedom are the lightest particles of the theory. As mentioned earlier, in QCD the low-lying degrees of freedom are the pions resulting from the spontaneous breaking of chiral symmetry. Therefore, an essential ingredient in the construction of the effective low-energy theory is the requirement that it correctly incorporates the chiral symmetries of the original theory. Since the up- and down-quark are not completely massless in nature, the pions are not massless but have a small mass of about 140 MeV. They are much lighter than the lightest non-Goldstone particles, such as the $`\rho `$-meson or the nucleons, which have a mass of about $`\mathrm{\Lambda }`$ 1 GeV. This means that at sufficiently low energies, the QCD partition function is well approximated by the partition function of an effective low-energy theory involving only pions.
We now consider QCD in a finite Euclidean volume $`V_4=L^4`$. The partition function is then dominated by the pions if
$$\frac{1}{\mathrm{\Lambda }}L.$$
(2)
This statement follows from Eq. (4) below by comparing the contribution of the pion, $`\mathrm{exp}(m_\pi L)`$, to that of a heavier particle, $`\mathrm{exp}(\mathrm{\Lambda }L)`$. The interactions of the pions are described by an effective chiral Lagrangian that will be given in Sec. 4. The fields in this Lagrangian can be separated into zero-momentum modes (constant fields) and nonzero-momentum modes. It was realized by Gasser and Leutwyler that there exists a kinematic regime where the fluctuations of the zero-momentum modes dominate the fluctuations of the nonzero-momentum modes. This regime is given by the condition
$$L\frac{1}{m_\pi },$$
(3)
where $`m_\pi `$ is the pion mass. Intuitively, this means that the wavelength of the pion is much larger than the linear extent of the box. Thus, the pion field does not vary appreciably over the size of the box, so the derivative terms are small. (These statements are quantified in Sec. 4.) We need consider only the zero-momentum modes in the regime of Eq. (3). For constant fields, the spacetime integral in the action can be replaced by the four-volume, and, effectively, we only have to deal with a much simpler zero-dimensional theory. However, the global symmetries remain important.
We are now ready to make the connection to chiral RMT. A random matrix, as introduced in the previous subsection, contains independently distributed random variables. Each element of the matrix has the same average size. On the other hand, the matrix elements of the Dirac operator contain specific correlations due to the gauge fields and the kinetic terms, which are not included in the RMT. Only the global symmetries of QCD are included in the RMT. Since the random matrix elements do not have any spacetime dependence, we expect that in the domain of Eq. (3), which is dominated by constant fields, the RMT reproduces the mass dependence of the finite-volume partition function.
To summarize, the three ingredients in the construction of the finite-volume effective theory from QCD in the regime of Eq. (3) are the following: $`(a)`$ the global symmetries of QCD, $`(b)`$ the spontaneous breaking of chiral symmetry, and $`(c)`$ the fact that the partition function is dominated by the constant Goldstone fields. There is an almost one-to-one correspondence to the properties of the chiral random matrix model : $`(a)`$ the Dirac matrix is replaced by a random matrix with the same global symmetries, $`(b)`$ chiral symmetry is broken spontaneously in the limit of infinitely large random matrices, and $`(c)`$ the random matrix elements do not have any spacetime dependence.
## 2 QCD PARTITION FUNCTION AND DIRAC SPECTRUM
After the introductory remarks on our basic philosophy in Sec. 1, we now present a detailed discussion of the properties of the QCD partition function and the Dirac operator.
### 2.1 Basic Definitions
The QCD partition function is defined as
$$Z^{\mathrm{QCD}}=\mathrm{Tr}e^{\beta H},$$
(4)
where $`\beta `$ is the inverse temperature and $`H`$ is the QCD Hamiltonian in a box of volume $`V_3=L^3`$. It can be rewritten as a (suitably regularized) functional integral in Euclidean space,
$$Z^{\mathrm{QCD}}=DA_\mu \underset{f=1}{\overset{N_f}{}}det(𝒟+m_f)e^{S_{\mathrm{YM}}},$$
(5)
where $`N_f`$ is the number of quark flavors and $`S_{\mathrm{YM}}`$ is the Euclidean Yang-Mills action. The Euclidean four-volume is $`V_4=V_3\beta `$. The $`A_\mu `$ are non-Abelian gauge fields, which can be represented as
$$A_\mu =A_\mu ^a\frac{T_a}{2},$$
(6)
where the $`T_a`$ are the generators of the gauge group SU($`N_c`$). The number of colors is denoted by $`N_c`$. The QCD Dirac operator is given by
$$𝒟=\gamma _\mu (_\mu +iA_\mu ).$$
(7)
This operator is anti-Hermitian, $`𝒟^{}=𝒟`$. The $`\gamma _\mu `$ are Euclidean gamma matrices with $`\{\gamma _\mu ,\gamma _\nu \}=2\delta _{\mu \nu }`$. We use the chiral representation in which $`\gamma _5\gamma _1\gamma _2\gamma _3\gamma _4=\mathrm{diag}(1,1,1,1)`$.
### 2.2 Global Symmetries
The structure of the QCD Lagrangian is to a large extent determined by symmetries and renormalizability. As noted above, it is important to analyze these symmetries to construct the correct effective low-energy theory and the correct random matrix model. We now discuss three important global symmetries of the partition function and the Dirac operator.
#### 2.2.1 CHIRAL SYMMETRY AND TOPOLOGY
The Dirac operator satisfies
$$\{\gamma _5,𝒟\}=0.$$
(8)
This relation is a compact expression of chiral symmetry, i.e. of the fact that right-handed and left-handed quarks can be rotated independently. One can write down an eigenvalue equation for $`𝒟`$,
$$𝒟\psi _n=i\lambda _n\psi _n,$$
(9)
where the eigenvalues and eigenfunctions depend on the gauge field in Eq. (7). Using Eq. (8) one can show that the nonzero eigenvalues of $`𝒟`$ occur in pairs $`\pm i\lambda _n`$ with eigenfunctions $`\psi _n`$ and $`\gamma _5\psi _n`$. There can also be eigenvalues equal to zero, $`\lambda _n=0`$. The corresponding eigenfunctions can be arranged to be simultaneous eigenfunctions of $`\gamma _5`$ with eigenvalue $`\pm 1`$, i.e. these states have definite chirality. Denoting the number of zero eigenvalues with positive and negative chirality by $`N_+`$ and $`N_{}`$, respectively, the Atiyah-Singer index theorem states that $`\nu N_+N_{}`$ is a topological invariant that does not change under continuous changes of the gauge field. However, the individual numbers $`N_+`$ and $`N_{}`$ are not protected by topology, i.e. very small deformations of the gauge field will lift accidental zero modes. Thus, unless we impose very special constraints on the gauge fields, we generically have either $`N_+=0`$ or $`N_{}=0`$.
In a chiral basis with $`\gamma _5\psi ^{R/L}=\pm \psi ^{R/L}`$, one can use Eq. (8) to show that $`\overline{\psi }_m^R|𝒟|\psi _n^L=0=\overline{\psi }_m^L|𝒟|\psi _n^R`$ for all $`m`$ and $`n`$, where $`\overline{\psi }=\psi ^{}\gamma _0`$. From this property and the fact that $`𝒟`$ is anti-Hermitian, it follows that the Dirac operator has the matrix structure
$$𝒟=\left(\begin{array}{cc}0& iW\\ iW^{}& 0\end{array}\right).$$
(10)
This off-diagonal block structure is characteristic for systems with chiral symmetry. If there are $`n+\nu `$ right-handed and $`n`$ left-handed modes, the matrix $`W`$ has dimension $`(n+\nu )\times n`$, and the matrix $`𝒟`$ in Eq. (10) has $`|\nu |`$ eigenvalues equal to zero.
The QCD partition function can be decomposed into sectors of definite topological charge $`\nu `$,
$$Z^{\mathrm{QCD}}(\theta )=\underset{\nu =\mathrm{}}{\overset{\mathrm{}}{}}e^{i\nu \theta }Z_\nu ^{\mathrm{QCD}}.$$
(11)
Here, the $`\theta `$-angle is the coefficient of the topological $`F\stackrel{~}{F}`$ term (which violates $`P`$ and $`CP`$ conservation) in the Lagrangian. We now show that the $`\theta `$-dependence of the partition function can be absorbed into the phase of the quark masses. To this end, we introduce complex masses $`m_f`$ for the right-handed quarks and the complex conjugate masses $`m_f^{}`$ for the left-handed quarks. The partition function then reads
$$Z_\nu ^{\mathrm{QCD}}=\underset{f}{}\widehat{m}_f^{|\nu |}\underset{\lambda _n>0}{}(\lambda _n^2+m_fm_f^{})_\nu ,$$
(12)
where $`\widehat{m}_f=m_f`$ ($`m_f^{}`$) for $`\nu 0`$ ($`\nu <0`$). In Eq. (12), the average is only over the gauge field configurations with topological charge $`\nu `$, weighted by the Yang-Mills action. Since
$$e^{i\nu \theta }\underset{f}{}\widehat{m}_f^{|\nu |}=\{\begin{array}{cc}\left(e^{i\theta }_fm_f\right)^\nu \hfill & \mathrm{for}\nu 0\hfill \\ \left[\left(e^{i\theta }_fm_f\right)^{}\right]^\nu \hfill & \mathrm{for}\nu <0,\hfill \end{array}$$
(13)
it follows that the $`\theta `$-dependence of $`Z^{\mathrm{QCD}}`$ is entirely determined by the combination $`e^{i\theta }_fm_f`$.
#### 2.2.2 FLAVOR SYMMETRIES
The second global symmetry is flavor symmetry, whose spontaneous breaking has profound implications for the hadron spectrum. To make this symmetry explicit, we rewrite the fermion determinant in $`Z^{\mathrm{QCD}}`$ as a Grassmann integral,
$$\underset{f=1}{\overset{N_f}{}}det(𝒟+m_f)=𝑑\psi 𝑑\overline{\psi }\mathrm{exp}\left[d^4x\underset{f=1}{\overset{N_f}{}}\overline{\psi }_f(𝒟+m_f)\psi _f\right],$$
(14)
where $`\overline{\psi }=\psi ^{}\gamma _0`$. Again going to a chiral basis with $`\gamma _5\psi ^{R/L}=\pm \psi ^{R/L}`$, the exponent on the right-hand side of Eq. (14) can be rewritten as
$$d^4x\underset{f=1}{\overset{N_f}{}}(\overline{\psi }_f^RiW^{}\psi _f^R+\overline{\psi }_f^LiW\psi _f^L+\overline{\psi }_f^Rm_f\psi _f^L+\overline{\psi }_f^Lm_f\psi _f^R).$$
(15)
In the chiral limit where all $`m_f=0`$, the fermion determinant is invariant under the transformations
$`\begin{array}{cc}\psi ^LU\psi ^L\hfill & \overline{\psi }^L\overline{\psi }^LU^1\hfill \\ \overline{\psi }^R\overline{\psi }^RV^1\hfill & \psi ^RV\psi ^R.\hfill \end{array}`$ (18)
The only condition on $`U`$ and $`V`$ is that their inverses must exist. If the number of right-handed states $`N_R`$ is equal to the number of left-handed states $`N_L`$, the symmetry is thus<sup>1</sup><sup>1</sup>1Actually, the symmetry group is $`\mathrm{Gl}(N_f)\times \mathrm{Gl}(N_f)`$. The restriction to $`\mathrm{U}(N_f)`$ is a consequence of the Riemannian nature of the integration manifold (see Sec. 4.3). $`\mathrm{U}(N_f)\times \mathrm{U}(N_f)`$. However, if $`N_RN_L`$, i.e. if $`\nu 0`$, the axial symmetry group $`\mathrm{U}_\mathrm{A}(1)`$ is broken explicitly by instantons or the anomaly. A second U(1) group, $`\mathrm{U}_\mathrm{V}(1)`$, corresponds to the conservation of baryon number. The full flavor symmetry group in the chiral limit is thus given by $`G=\mathrm{SU}(N_f)\times \mathrm{SU}(N_f)`$.
If the quark masses are nonzero, the axial SU($`N_f`$) subgroup with $`U=V^1`$ is broken explicitly by the mass term. The SU($`N_f`$) vector symmetry (with $`U=V`$) is good for degenerate quark masses ($`m_f=m`$ for all $`f`$) but is broken explicitly for different quark masses ($`m_fm_f^{}`$ for all $`ff^{}`$).
What is much more important than the explicit breaking, however, is the spontaneous breaking of the axial-flavor symmetry. For an axial-flavor–symmetric ground state, the vacuum expectation value $`\overline{\psi }\psi =\overline{\psi }^R\psi ^L+\overline{\psi }^L\psi ^R`$ would be zero. However, phenomenological arguments and lattice QCD simulations indicate that $`\overline{\psi }\psi (240\mathrm{MeV})^3`$. The spontaneous breaking of the axial symmetry also follows from the absence of parity doublets and the presence of Goldstone bosons, the pions. The quantity $`\overline{\psi }\psi `$ is only invariant if $`U=V`$, i.e. the vacuum state is symmetric under the flavor group $`H=`$ SU<sub>V</sub>($`N_f`$). Thus, the vector symmetries are unbroken whereas the axial symmetries are maximally broken. The first of these statements is in agreement with the Vafa-Witten theorem , which states that in vector-like theories such as QCD, vector symmetries cannot be spontaneously broken. The reasons why the axial symmetries are maximally broken are less well understood. The Goldstone manifold is $`G/H=\mathrm{SU}(N_f)`$, and so there are $`N_f^21`$ Goldstone bosons.
#### 2.2.3 ANTI-UNITARY SYMMETRIES
Third, we consider the anti-unitary symmetries of the Dirac operator. According to a fundamental theorem by Wigner, a symmetry in quantum mechanics is either unitary or anti-unitary. An anti-unitary symmetry operator $`A`$ can always be written as $`A=UK`$, where $`U`$ is a unitary operator and $`K`$ is the complex conjugation operator. Below we always consider spectra of an irreducible subspace of the unitary symmetries. If $`A=UK`$ is an anti-unitary symmetry of the Dirac operator, then the symmetry operator $`A^2=(UK)^2=UU^{}`$ is unitary, and in an irreducible subspace it is necessarily a multiple of the identity, $`UU^{}=\lambda \mathrm{𝟏}`$. Because of this relation, $`U`$ and $`U^{}`$ commute so that $`\lambda `$ is real. By unitarity we have $`|\lambda |=1`$, which yields $`\lambda =\pm 1`$. Therefore, the anti-unitary symmetries can be classified according to the sign of $`A^2`$. There are three possibilities for the classification of Hermitian (or anti-Hermitian) operators: $`(a)`$ there are no anti-unitary symmetries (denoted by the Dyson index $`\beta =2`$); $`(b)`$ if $`A^2=1`$, it is possible to construct a basis in which the operator is real (denoted by the Dyson index $`\beta =1`$); $`(c)`$ if $`A^2=1`$, it is possible to construct a basis in which the matrix elements of the operator can be organized into real (or self-dual) quaternions (denoted by the Dyson index $`\beta =4`$).
For QCD with $`N_c3`$ and fermions in the fundamental representation of the gauge group, there are no anti-unitary operators that commute with the Dirac operator. This means that the matrix $`W`$ in Eq. (10) is a general complex matrix with no further symmetries. There are two cases with nontrivial anti-unitary symmetries: QCD with two colors and fermions in the fundamental representation, and QCD with any number of colors and adjoint fermions. We will discuss them next.
For QCD with two colors and fermions in the fundamental representation, the Dirac operator is given by Eq. (7) with $`A_\mu =A_\mu ^a\tau _a/2`$, where the $`\tau _a`$ are the SU(2) Pauli matrices in color space. The anti-unitary symmetry in this case is
$$[C\tau _2K,i𝒟]=0,$$
(19)
where $`C=\gamma _2\gamma _4`$ is the charge conjugation matrix. The square of the anti-unitary operator is $`(C\tau _2K)^2=1`$. In this case, it is possible to find a basis in which the matrix $`W`$ in Eq. (10) is real for all gauge field configurations .
As a consequence of the pseudoreality of SU(2), the symmetry of the QCD partition function in the chiral limit is enlarged to SU(2$`N_f`$) . An axial U(1) is broken explicitly by instantons or the axial anomaly, as for $`N_c3`$. The chiral condensate is only invariant under an $`\mathrm{Sp}(2N_f)`$ subgroup of $`\mathrm{SU}(2N_f)`$. The Vafa-Witten theorem, prohibiting the spontaneous breaking of global vector symmetries and assuming maximum breaking of the axial symmetries, thus predicts a pattern of spontaneous chiral symmetry breaking given by SU($`2N_f`$) $``$ Sp($`2N_f`$). The Goldstone manifold is the coset SU($`2N_f`$)/Sp($`2N_f`$), which is the set of antisymmetric unitary matrices. We thus have $`2N_f^2N_f1`$ Goldstone bosons .
For fermions in the adjoint representation of the gauge group and any number of colors, the Dirac operator is $`𝒟_{ab}=\gamma _\mu (_\mu \delta _{ab}+f_{abc}A_\mu ^c)`$, where $`a`$ and $`b`$ are color indices and the $`f_{abc}`$ are the structure constants of SU($`N_c`$). In this case, the anti-unitary symmetry is
$$[CK,i𝒟]=0.$$
(20)
Because of $`(CK)^2=1`$, it follows that the eigenvalues of the Dirac operator are twofold degenerate with linearly independent eigenfunctions. In this case, it is possible to choose a basis in which the matrix elements of $`W`$ are real (or self-dual) quaternions for all gauge fields. The eigenvalues of such a matrix are unit quaternions and therefore doubly degenerate in the representation of the Dirac matrix as complex numbers.
Restricting ourselves to the case of even $`N_f`$, we can show that for $`N_f`$ Majorana fermions, the flavor symmetry group is now SU($`N_f`$). In this case, the chiral condensate is only invariant under an O($`N_f`$) subgroup of SU($`N_f`$). Applying again the Vafa-Witten theorem with maximum breaking of axial symmetry, we expect the pattern of spontaneous chiral symmetry breaking according to SU($`N_f`$) $``$ O($`N_f`$), with the Goldstone manifold given by the coset SU($`N_f`$)/O($`N_f`$). This is the set of symmetric unitary matrices. We thus have $`(N_f+2)(N_f1)/2`$ Goldstone bosons .
On the lattice, the symmetries of the Dirac operator may be different from the continuum symmetries . We return to this point in Sec. 5.2.
### 2.3 Dirac Spectrum
Based on the eigenvalue equation (9), we define the spectral density of the Dirac operator by
$$\rho (\lambda )=\underset{n}{}\delta (\lambda \lambda _n),$$
(21)
where the average is over gauge fields weighted by the full QCD action. The spectral density is important because of its relation with the order parameter for spontaneous chiral symmetry breaking, the chiral condensate $`\overline{\psi }\psi `$. It was shown by Banks and Casher that
$$\mathrm{\Sigma }|\overline{\psi }\psi |=\frac{\pi \rho (0)}{V}.$$
(22)
To be precise, we should have written $`\mathrm{\Sigma }=lim_{\epsilon 0}lim_{m0}lim_V\mathrm{}\pi \rho (\epsilon )/V`$, where it is important that the limits are taken in the order indicated. (In the normalization of Eq. (21) the spectral density is proportional to the volume, so the explicit factor of $`1/V`$ in (22) is canceled to yield a finite result.)
The relation (22) can readily be derived. The chiral condensate is given by
$$\overline{\psi }\psi =\underset{m0}{lim}\underset{V\mathrm{}}{lim}\frac{1}{VN_f}\frac{}{m}\mathrm{log}Z^{\mathrm{QCD}}(m).$$
(23)
From Eq. (5), this yields
$$\overline{\psi }\psi =\underset{m0}{lim}\underset{V\mathrm{}}{lim}\frac{1}{V}\underset{n}{}\frac{1}{i\lambda _n+m}.$$
(24)
Since the nonzero eigenvalues occur in pairs $`\pm i\lambda _n`$, their contribution to the sum can be written as $`2m/(\lambda _n^2+m^2)`$ (with $`\lambda _n>0`$). For gauge fields with topological charge $`\nu `$, the zero modes contribute a term $`|\nu |/(mV)`$. Assuming $`\nu ^2V`$, we can drop these contributions in the limit $`V\mathrm{}`$. In the same limit, the sum in Eq. (24) can be converted to an integral. In the limit $`m0`$, we have $`2m/(\lambda ^2+m^2)\pi \delta (\lambda )`$, which yields Eq. (22). As discussed above in Sec. 1.1, spontaneous chiral symmetry breaking is encoded in an accumulation of the small Dirac eigenvalues; for the order parameter to be nonzero, we need $`\rho (0)/V>0`$.
An immediate consequence of the Banks-Casher relation is that the small eigenvalues are spaced as
$$\mathrm{\Delta }\lambda =\frac{1}{\rho (0)}=\frac{\pi }{V\mathrm{\Sigma }},$$
(25)
provided that $`\rho (0)/V>0`$. This naturally defines a scale
$$z=\lambda V\mathrm{\Sigma }$$
(26)
for the study of the distribution of individual eigenvalues. For this purpose, it is convenient to define the so-called microscopic spectral density
$$\rho _s(z)=\underset{V\mathrm{}}{lim}\frac{1}{V\mathrm{\Sigma }}\rho \left(\frac{z}{V\mathrm{\Sigma }}\right).$$
(27)
This function describes the extreme infrared properties of the Dirac spectrum. Based on the arguments in the introduction, we expect it to be completely determined by the global symmetries of the Dirac operator. As we demonstrate below, $`\rho _s(z)`$ can be computed both from the low-energy effective theory and from chiral RMT; the results coincide. Further confirmation of these ideas will come from results of lattice QCD simulations.
At this point we would like to add some remarks on possible ultraviolet divergences. As an example, we consider the chiral condensate in the $`V\mathrm{}`$ limit, but before the limit $`m0`$ is taken. Converting the sum in Eq. (24) to an integral and dropping the contribution from the zero eigenvalues, we obtain
$$\overline{\psi }\psi =\underset{\mathrm{\Lambda }\mathrm{}}{lim}\underset{m0}{lim}_0^\mathrm{\Lambda }𝑑\lambda \frac{2m\rho (\lambda )/V}{\lambda ^2+m^2},$$
(28)
where we have introduced an ultraviolet cutoff $`\mathrm{\Lambda }`$ that must be removed at the end of the calculation. Asymptotically, the spectral density behaves as $`\rho (\lambda )V\lambda ^3`$, which means that the integral is ultraviolet-divergent. However, this divergence does not contribute to $`\overline{\psi }\psi `$ if the $`m0`$ limit is taken first to yield the usual Banks-Casher relation. Alternatively, we can first subtract the divergent contributions to the integral (28) and remove the cutoff before taking the chiral limit.
A slightly different situation arises if we consider the dependence of the chiral condensate on a valence (or spectral) mass $`m_v`$ that does not appear in the average, so that $`\rho (\lambda )`$ is independent of $`m_v`$. The quantity
$$\mathrm{\Sigma }(m_v)=\underset{\mathrm{\Lambda }\mathrm{}}{lim}_0^\mathrm{\Lambda }𝑑\lambda \frac{2m_v\rho (\lambda )/V}{\lambda ^2+m_v^2}$$
(29)
is subject to ultraviolet divergences. Using again the fact that $`\rho (\lambda )V\lambda ^3`$ for large $`\lambda `$, we find that the leading divergence is $`m_v\mathrm{\Lambda }^2`$. If we consider valence masses on the microscopic scale (26) we have $`m_v1/(V\mathrm{\Sigma })`$, and if the $`V\mathrm{}`$ limit is taken before the $`\mathrm{\Lambda }\mathrm{}`$ limit, the ultraviolet divergences are removed.
In lattice QCD simulations with lattice spacing $`a`$, the cutoff is $`\mathrm{\Lambda }=1/a`$. In the continuum limit, $`a0`$, both the coupling constant and the $`\overline{\psi }\psi `$-operator have to be renormalized. For a discussion of the ultraviolet divergences in a recent lattice study of $`\mathrm{\Sigma }(m_v)`$ at finite volume, finite quark mass, and finite lattice spacing, see Ref. .
## 3 CHIRAL RANDOM MATRIX THEORY
### 3.1 Introduction of the Model
In this section we introduce a chiral random matrix theory (chRMT) with the global symmetries of the QCD Dirac operator. In the spirit of the invariant random matrix ensembles, we construct a model with eigenfunctions distributed uniformly over the unitary unit sphere. This is achieved by choosing Gaussian-distributed random matrix elements. We thus arrive at the following chRMT :
$$Z_{N_f,\nu }^\beta (m_1,\mathrm{},m_{N_f})=DW\underset{f=1}{\overset{N_f}{}}det(𝒟+m_f)e^{\frac{N\beta }{4}\mathrm{Tr}v(W^{}W)},$$
(30)
where $`\beta `$ is the Dyson index,
$$𝒟=\left(\begin{array}{cc}0& iW\\ iW^{}& 0\end{array}\right),$$
(31)
and $`W`$ is an $`n\times m`$ matrix with $`\nu =mn`$ and $`N=n+m`$. The interpretation of this model is that $`N`$ low-lying modes interact via a random interaction. A natural representation of this model is in the form of gauge field configurations given by a liquid of instantons. Then the low-lying modes are the zero modes of each instanton. We assume that $`\nu `$ does not exceed $`\sqrt{N}`$ so that, to a good approximation, $`n=N/2`$ for large $`N`$. The parameter $`N`$ is identified as the dimensionless volume of spacetime. For the formulation of this model with explicit factors of $`N/V`$ included, see Ref. . The potential $`v`$ is defined by
$$v(\varphi )=\underset{k1}{}a_k\varphi ^k.$$
(32)
The simplest case is the Gaussian case, where $`v(\varphi )=\mathrm{\Sigma }^2\varphi `$. It can be shown (see Sec. 3.5) that the microscopic spectral density does not depend on the higher-order terms in this potential provided that the average spectral density near zero remains nonzero. The matrix elements of $`W`$ are either real \[$`\beta =1`$, chiral Gaussian Orthogonal Ensemble (chGOE)\], complex \[$`\beta =2`$, chiral Gaussian Unitary Ensemble (chGUE)\], or quaternion real \[$`\beta =4`$, chiral Gaussian Symplectic Ensemble (chGSE)\]. In the latter case, the eigenvalues of $`𝒟`$ are doubly degenerate, and the use of Majorana fermions is implemented by replacing the determinant by its square root. For a non-Gaussian potential $`v(\varphi )`$, we will omit the G in the abbreviations and use chOE, chUE, and chSE, respectively. Two earlier attempts to describe QCD Dirac eigenvalues used the Wigner-Dyson ensembles instead of the above chiral ensembles .
This model reproduces the following symmetries of the QCD partition function:
* The $`\mathrm{U}_\mathrm{A}(1)`$ symmetry. All eigenvalues of the random matrix Dirac operator occur in pairs $`\pm i\lambda _n`$ or are zero.
* The topological structure of the QCD partition function. The Dirac matrix has exactly $`|\nu |=|nm|`$ zero eigenvalues. This identifies $`\nu `$ as the topological sector of the model.
* The flavor symmetry, which is the same as in QCD. For $`\beta =2`$ it is $`\mathrm{SU}(N_f)\times \mathrm{SU}(N_f)`$, for $`\beta =1`$ it is $`\mathrm{SU}(2N_f)`$, and for $`\beta =4`$ it is $`\mathrm{SU}(N_f)`$ (each Majorana flavor counts as 1/2 Dirac flavor).
* The chiral symmetry, which is broken spontaneously with a chiral condensate given by
$$\mathrm{\Sigma }=\underset{N\mathrm{}}{lim}\pi \rho (0)/N.$$
(33)
($`N`$ is interpreted as the dimensionless volume of spacetime.) The symmetry-breaking pattern is $`\mathrm{SU}(N_f)\times \mathrm{SU}(N_f)\mathrm{SU}(N_f)`$, $`\mathrm{SU}(2N_f)\mathrm{Sp}(2N_f)`$, and $`\mathrm{SU}(N_f)\mathrm{O}(N_f)`$ for $`\beta =2`$, 1, and 4, respectively — the same as in QCD .
* The anti-unitary symmetries. These are implemented by choosing the matrix elements of $`W`$ to be real, complex, or quaternion real for $`\beta =1`$, $`\beta =2`$, and $`\beta =4`$, respectively.
Along with the invariant random matrix ensembles, the chiral ensembles are part of a larger classification scheme that also includes ensembles for the description of disordered superconductors . In total, 10 different families of random matrix ensembles have been identified. They correspond one-to-one with the Cartan classification of symmetric spaces .
The uniform distribution of the eigenfunctions over the unitary unit sphere is expressed as the invariance
$$WU^{}WV,$$
(34)
where the $`n\times n`$ matrix $`U`$ and the $`m\times m`$ matrix $`V`$ are orthogonal matrices for $`\beta =1`$, unitary matrices for $`\beta =2`$, and symplectic matrices for $`\beta =4`$. This invariance makes it possible to express the partition function in terms of eigenvalues of $`W`$ defined by
$$W=U^{}\mathrm{\Lambda }V.$$
(35)
Here, $`\mathrm{\Lambda }`$ is a diagonal matrix with real diagonal matrix elements $`\lambda _k0`$. In terms of the eigenvalues, the partition function (30) is given by
$$Z_{N_f,\nu }^\beta (m_1,\mathrm{},m_{N_f})=𝑑\lambda |\mathrm{\Delta }(\lambda ^2)|^\beta \underset{k}{}\lambda _k^{\beta |\nu |+\beta 1}e^{\frac{N\beta }{4}v(\lambda _k^2)}\underset{f}{}m_f^{|\nu |}(\lambda _k^2+m_f^2),$$
(36)
where the Vandermonde determinant is defined by
$$\mathrm{\Delta }(\lambda ^2)=\underset{k<l}{}(\lambda _k^2\lambda _l^2).$$
(37)
In Eq. (36) and elsewhere in this review, we have omitted the normalization constant of the partition function. From the joint eigenvalue distribution (which is the integrand of Eq. (36)), we see that a nonzero topological charge $`\nu `$ can be introduced by adding $`\beta |\nu |/2`$ massless flavors to the theory with $`\nu =0`$. Therefore, this duality between flavor and topology is a general feature of all correlation functions. For a discussion of this duality in terms of finite volume partition functions, see Refs. .
### 3.2 Sigma Model Representation of Chiral Random Matrix Models
The chiral random matrix partition function can be evaluated using standard random matrix methods. For simplicity, let us consider the case $`\beta =2`$. The fermion determinant can be written as a Grassmann integral, and averaging over the Gaussian distribution function results in a four-fermion interaction. Because of the underlying unitary invariance (34) of the chRMT, the fermionic variables only appear in invariant combinations of the form $`\sigma ^{fg}\overline{\psi }_i^f\psi _i^g`$, and the partition function can be rewritten identically in terms of these variables as
$$Z_\nu (m)=D\sigma e^{n\mathrm{\Sigma }^2\mathrm{Tr}\sigma \sigma ^{}}\stackrel{\nu }{det}(\sigma +m\mathrm{𝟏})\stackrel{n}{det}[(\sigma +m\mathrm{𝟏})(\sigma ^{}+m\mathrm{𝟏})],$$
(38)
where the quark masses have been taken real and degenerate for simplicity. The integration is over the real and imaginary parts of the $`N_f\times N_f`$ arbitrary complex matrix $`\sigma `$. For $`m=0`$ this integral is invariant under $`\sigma U\sigma V^1`$ with $`U`$ and $`V\mathrm{SU}(N_f)`$. Let us calculate the integrals by a saddle-point approximation. The saddle-point equation reads (notice that $`\nu n`$)
$$\mathrm{\Sigma }^2(\sigma ^{}+m)\sigma =1.$$
(39)
The condensate at the saddle point is given by
$$|\overline{\psi }\psi |=\frac{1}{2nN_f}_m\mathrm{log}Z=\frac{\mathrm{\Sigma }^2}{2N_f}\mathrm{Tr}(\sigma +\sigma ^{}),$$
(40)
where the expectation value is with respect to the partition function (38). The solution of the saddle-point equation, $`\overline{\sigma }`$, is proportional to the identity matrix. In the chiral limit, we then find $`\overline{\sigma }=1/\mathrm{\Sigma }`$, and we can identify $`\mathrm{\Sigma }`$ as the chiral condensate.
The fluctuations about the saddle point can be separated into massive modes with a curvature of order $`n`$ and Goldstone modes with a curvature of order $`mn`$. In the limit $`n\mathrm{}`$, the partition function (38) can be simplified by keeping the integrals over the Goldstone manifold and performing the remaining integrals by a saddle-point integration. This results in the partition function
$$Z_\nu (m)=_{U\mathrm{U}(N_f)}DU\stackrel{\nu }{det}Ue^{n\mathrm{\Sigma }\mathrm{Tr}(MU+M^{}U^1)},$$
(41)
which is the familiar finite volume partition function to be discussed in Sec. 4.1.
### 3.3 $`\theta `$-Dependence of the Partition Function
The $`\theta `$-dependence of the partition function is obtained by summing over all topological sectors of the partition function according to Eq. (11), which applies to full QCD where the partition function is obtained by integrating over the gauge fields. The terms $`Z_\nu `$ in this equation reflect the probability of encountering gauge fields with topological charge $`\nu `$.
For the effective theory, the partition function in the topological sector $`\nu `$ is given in Eq. (41). However, this $`\nu `$-dependence is only due to the fermion determinant. In the sum over $`\nu `$, we should therefore add weight factors $`P(\nu )`$ that take into account the distribution of topological charge of the quenched gauge fields (i.e. gauge fields generated without the fermion determinant). This yields
$$Z^{\mathrm{eff}}(\theta ,m)=\underset{\nu }{}P(\nu )e^{i\nu \theta }Z_\nu ^{\mathrm{eff}}(m).$$
(42)
Below, we argue that the weight factors can be ignored for light quarks but that they are necessary when one considers the quenched theory (which corresponds to $`N_f=0`$ or, equivalently, to the limit of very heavy quarks).
Our starting point is that the topological susceptibility of the quenched gauge fields is nonvanishing, i.e.
$$\frac{\nu ^2_q}{V}0.$$
(43)
Invoking the central limit theorem, we assume that $`P(\nu )`$ is Gaussian,
$$P(\nu )=\frac{1}{\sqrt{2\pi \nu ^2_q}}e^{\frac{\nu ^2}{2\nu ^2_q}},$$
(44)
which has been verified by quenched lattice QCD simulations.
We analyze the partition function (42) for $`Z_\nu ^{\mathrm{eff}}(m)`$ given by the finite-volume partition function in Eq. (41). Replacing the sum over $`\nu `$ by an integral, we obtain after integration
$$Z^{\mathrm{eff}}(\theta ,m)=_{U\mathrm{U}(N_f)}DUe^{\frac{\nu ^2_q}{2}(\theta i\mathrm{Tr}\mathrm{log}U)^2+\frac{1}{2}\mathrm{\Sigma }V\mathrm{Tr}(MU+M^{}U^1)}.$$
(45)
The exponent in the integrand corresponds to the effective static potential of QCD with a large number of colors , and therefore the chRMT partition function reproduces all identities that have been derived from this effective potential . For example, the topological charge is screened by light quarks. For a careful analysis of the periodicity requirements of the $`\theta `$-dependence of the effective partition function (45), see Ref. .
Let us illustrate the screening of topological charge for $`N_f=1`$ and real $`m`$, for which
$$Z^{\mathrm{eff}}(\theta ,m)=𝑑\varphi e^{\frac{\nu ^2_q}{2}(\theta \varphi )^2+\mathrm{\Sigma }Vm\mathrm{cos}\varphi }.$$
(46)
The integral over $`\varphi `$ can be performed by a saddle-point approximation resulting in
$$Z^{\mathrm{eff}}(\theta ,m)=e^{mV\mathrm{\Sigma }\mathrm{cos}\theta }.$$
(47)
The topological susceptibility is given by
$$\nu ^2=\frac{^2}{\theta ^2}\mathrm{log}Z=mV\mathrm{\Sigma }(\mathrm{for}\theta =0)$$
(48)
so that the topological charge is completely screened in the chiral limit. This also shows that the results are insensitive to our choice of the distribution function $`P(\nu )`$. For example, in an instanton liquid interpretation of the chiral random matrix model, the natural choice for $`P(\nu )`$ is the binomial distribution $`B(N\nu ,\nu )`$, resulting in the same effective potential .
In fact, it can be concluded from Eq. (48) that for light quarks the distribution $`P(\nu )`$ can be ignored altogether. Let us again illustrate this for $`N_f=1`$, for which the integral in Eq. (41) results in the finite volume partition function
$$Z_\nu ^{\mathrm{eff}}(m)=I_\nu (mV\mathrm{\Sigma }),$$
(49)
where $`I_\nu `$ is a modified Bessel function. Summing over $`\nu `$ according to Eq. (11), i.e. without the weight factors $`P(\nu )`$, we obtain
$$Z^{\mathrm{eff}}(\theta ,m)=\underset{\nu }{}e^{i\nu \theta }I_\nu (mV\mathrm{\Sigma })=e^{mV\mathrm{\Sigma }\mathrm{cos}\theta },$$
(50)
where we have used a summation formula for modified Bessel functions. This is exactly the same result as in Eq. (47). For an extension of these results to more than one flavor, see Ref. .
### 3.4 Spectral Correlation Functions
In this section, we discuss the spectral correlation functions and the microscopic spectral density corresponding to the chRMT partition function. The global spectral density is a semicircle. It is not universal and we do not expect to find it in realistic physical systems. The universal quantities are local or microscopic spectral correlation functions. In general, a $`k`$-point spectral correlation function is defined by
$$R_k(\lambda _1,\mathrm{},\lambda _k)=\underset{\begin{array}{c}j_1\mathrm{}j_k=1\\ j_pj_q\end{array}}{\overset{N}{}}\delta (\lambda _1E_{j_1})\mathrm{}\delta (\lambda _kE_{j_k}),$$
(51)
where the $`E_{j_p}`$ are eigenvalues. The quantity $`R_k`$ represents the probability density to find $`k`$ eigenvalues, regardless of labeling, at $`\lambda _1,\mathrm{},\lambda _k`$. In particular, $`R_1(\lambda )=\rho (\lambda )`$. By “local” we mean that the energy differences $`|\lambda _p\lambda _q|`$ are of the order of a few mean level spacings.
Because of the $`\mathrm{U}_\mathrm{A}(1)`$ symmetry of QCD, all nonzero eigenvalues come in pairs $`\pm i\lambda _n`$ (see Sec. 2.2.1). This implies that the origin, $`\lambda =0`$, is a special point of the spectrum and that we have to consider the $`R_k`$ separately in the bulk of the spectrum and near $`\lambda =0`$. The latter is called the microscopic region or the “hard edge” of the spectrum.<sup>2</sup><sup>2</sup>2The reason for this nomenclature is that in a Fermi-gas formulation of RMT, where the eigenvalues are interpreted as the positions of particles, this symmetry corresponds to choosing a potential that is infinite for negative values of the eigenvalues. In this picture, the soft edge corresponds to a potential with a finite slope. In addition, there is the tail or “soft edge” of the spectrum, which we do not discuss. The bulk of the spectrum is the middle region, far from either edge.
Because the random matrix partition function is known in terms of an integral over the eigenvalues of the Dirac operator, see Eq. (36), it is possible to obtain the spectral density and all spectral $`k`$-point correlation functions by integration over $`Nk`$ eigenvalues. In the present case, the integrals can be performed most conveniently by the orthogonal polynomial method. We mention only the most important results in the sections below.
#### 3.4.1 BULK CORRELATIONS
Our conventions are such that if we have $`N`$ eigenvalues, the global spectral density is normalized to $`N`$ so that the mean level spacing is of order $`1/N`$. Local spectral correlation functions $`R_k(\lambda _1,\mathrm{},\lambda _k)`$ in the bulk of the spectrum are thus characterized by $`\lambda _p𝒪(1)`$ and $`|\lambda _p\lambda _q|𝒪(1/N)`$. Universal correlation functions are obtained by rescaling the eigenvalues according to the average local level spacing. This procedure, known as unfolding, is discussed in more detail in Sec. 5.1. The rescaled eigenvalues with average level spacing equal to unity are denoted by $`x_k`$.
In the bulk of the spectrum, the RMT results for the $`R_k`$ are identical for the chiral ensembles and the corresponding nonchiral ensembles . For simplicity, let us consider the simplest ensemble, the chUE. On the unfolded scale, we obtain in the limit $`N\mathrm{}`$
$$R_k(x_1,\mathrm{},x_k)=det[K(x_p,x_q)]_{p,q=1,\mathrm{},k}$$
(52)
with the sine kernel
$$K(x,y)=\frac{\mathrm{sin}\pi (xy)}{\pi (xy)}.$$
(53)
The fact that $`K(x,y)`$ depends only on $`|xy|`$ reflects the translational invariance of the spectral properties after unfolding. The results for the chOE and the chSE are somewhat more complicated than those for the chUE .
Other quantities can be derived from the $`R_k`$. As a measure of short-range correlations between the eigenvalues, one considers the nearest-neighbor spacing distribution $`P(s)`$, which is the probability density to have a spacing of $`s`$ between adjacent levels. In random matrix theories, one finds that $`P(s)s^\beta `$ for small values of $`s`$ and has a Gaussian tail for large spacings. Long-range spectral correlations are characterized by the level-number variance $`\mathrm{\Sigma }^2(L)`$ given by
$$\mathrm{\Sigma }^2(L)=(n(L)L)^2,$$
(54)
where $`n(L)`$ is the number of levels in an interval of length $`L`$, or by the spectral rigidity $`\mathrm{\Delta }_3(L)`$ , which is defined as an integral transform of $`\mathrm{\Sigma }^2(L)`$,
$$\mathrm{\Delta }_3(L)=\frac{2}{L^4}_0^L𝑑u(Lu)(L^2Luu^2)\mathrm{\Sigma }^2(u).$$
(55)
The advantage of using the $`\mathrm{\Delta }_3`$-statistic is that its statistical fluctuations are much smaller than those of the $`\mathrm{\Sigma }^2`$-statistic. The quantities $`\mathrm{\Sigma }^2`$ and $`\mathrm{\Delta }_3`$ can be derived from the two-point function $`R_2`$. To compute $`P(s)`$, one needs all $`k`$-point functions. These quantities are very different for systems with correlated eigenvalues (described by RMT) and systems with uncorrelated eigenvalues (described by the Poisson ensemble). For example, the large-$`L`$ behavior of the chUE results is given by $`\mathrm{\Sigma }^2(L)(\mathrm{ln}L)/\pi ^2`$ and $`\mathrm{\Delta }_3(L)(\mathrm{ln}L)/2\pi ^2`$, whereas for uncorrelated eigenvalues one finds $`\mathrm{\Sigma }^2(L)=L`$ and $`\mathrm{\Delta }_3(L)=L/15`$.
#### 3.4.2 MICROSCOPIC CORRELATIONS
In the context of QCD, the small eigenvalues are more interesting than those in the bulk because of their relation to spontaneous chiral symmetry breaking via the Banks-Casher relation. The local spectral correlation functions at the hard edge of the spectrum are characterized by $`\lambda _p𝒪(1/N)`$ and $`|\lambda _p\lambda _q|𝒪(1/N)`$. Unfolding in the microscopic region is easy; the eigenvalues are simply rescaled by the mean level spacing at $`\lambda =0`$, which is given by Eq. (25). By convention, we define the unfolded variables by $`x_p=\lambda _pV\mathrm{\Sigma }`$, i.e. the factor of $`\pi `$ is omitted. Note that $`VN`$. Thus, in the unfolded variables we have $`x_p𝒪(1)`$ and $`|x_px_q|𝒪(1)`$. As advertised above, the functional form of the $`R_k`$ is different in the microscopic region. The $`R_k`$ are still given by the determinant of a kernel according to Eq. (52), but the sine kernel is now replaced by the Bessel kernel ,
$$K(x,y)=\sqrt{xy}\frac{xJ_{\alpha +1}(x)J_\alpha (y)yJ_\alpha (x)J_{\alpha +1}(y)}{x^2y^2},$$
(56)
where $`J_\alpha `$ denotes Bessel functions and $`\alpha =N_f+|\nu |`$, with $`N_f`$ the number of massless flavors and $`\nu `$ the topological charge. A mathematical discussion of the Bessel kernel can be found in Ref. .
The RMT results for the chOE and the chSE are more complicated . Note that, in contrast to the situation in the bulk, the microscopic correlations depend on $`N_f`$ and $`\nu `$.
The one-point function is universal in the microscopic region. For the chUE, the microscopic spectral density, defined above in Eq. (27), is related to the kernel by $`\rho _s(z)=K(z,z)`$ and can be expressed as
$$\rho _s(z)=\frac{z}{2}\left[J_{N_f+|\nu |}^2(z)J_{N_f+|\nu |+1}(z)J_{N_f+|\nu |1}(z)\right].$$
(57)
It is also the generating function for the Leutwyler-Smilga sum rules, which are essentially inverse moments of $`\rho _s`$ (or of higher-order spectral correlation functions). They are discussed below. The results (56) and (57) were first derived using chiral RMT. In the meantime, Eq. (57) has been confirmed by an explicit calculation starting from a partially quenched chiral Lagrangian , which is discussed in the next section. The microscopic spectral density and the two-point correlation function for the case $`N_f+|\nu |=0`$ have also been derived by means of the supersymmetric method of RMT .
There is a very interesting connection between the Bessel kernel and the effective partition function. In addition to the partition function $`Z^{(N_f)}`$ for $`N_f`$ massless sea quarks, one can also compute the partition function $`Z^{(N_f+2)}`$ for $`N_f`$ massless quarks and two additional flavors with imaginary masses. One then finds the relation
$$K(x,y)=\frac{1}{2}(xy)^{N_f+1/2}\frac{Z^{(N_f+2)}(ix,iy)}{Z^{(N_f)}},$$
(58)
which can also be generalized to massive sea quarks. Thus, to obtain the $`R_k`$ for $`N_f`$ flavors, one needs to know the partition function for two additional flavors with imaginary masses. This result has a natural interpretation in terms of the effective spectral partition function (discussed below in Sec. 4). The two additional flavors are the spectral quark and its bosonic superpartner . The relation (58) once again demonstrates the equivalence of the random matrix formulation and the effective field theory in the zero mode domain. Additional consistency conditions for the finite volume partition function are discussed in Ref. .
### 3.5 Universality Proofs
In most random matrix calculations, the probability distribution of the random matrix elements is assumed to be Gaussian to simplify the calculation. Nonuniversal features such as the global spectral density depend on the choice of the probability distribution. On the other hand, universal results of RMT do not depend on this choice, nor on other deformations of the random matrix model. In addition to numerous empirical verifications of universal behavior in RMT, analytical proofs of universality have recently been constructed for two different types of deformations.
First, consider deformations that do not break the invariance of the theory under unitary transformations of the random matrix. In this case, the Gaussian distribution $`P(W)\mathrm{exp}(Na_1\mathrm{Tr}WW^{})`$ is replaced by the more general distribution $`P(W)\mathrm{exp}(N_{k=1}^{\mathrm{}}a_k\mathrm{Tr}(WW^{})^k)`$, i.e. the quadratic term in the exponent is replaced by an arbitrary polynomial. The advantage of the unitary invariance is that the probability distribution depends only on the eigenvalues of $`WW^{}`$. For the chiral ensembles, it has been shown that both the bulk and the microscopic spectral correlations remain unchanged on the unfolded scale . The weight function of QCD also contains the fermion determinant, and it has been shown that the universal results are unchanged if the random Dirac matrix in the determinant is replaced by a polynomial of that matrix . Universal microscopic correlations have also been found for the so-called Ginsparg-Wilson Dirac operator . Along with work on universality for the much more widely studied Wigner-Dyson ensembles , these findings make it clear that spectral correlations on the scale of the average level spacing are strongly universal, i.e. they do not change despite substantial variations of the average spectral density. There have also been several interesting results on the validity of universality for wide correlators, i.e. on macroscopic scales .
Second, consider deformations that do break the unitary invariance by adding an arbitrary deterministic matrix $`Y`$ to the random matrix $`W`$ in the Dirac operator. As we discuss in Sec. 6, such a model can provide a schematic description of QCD at nonzero temperature. In particular, the matrix $`Y`$ can be chosen in such a way that $`\rho (0)/V`$ vanishes so that chiral symmetry is restored. It has been proven that the bulk correlations on the unfolded scale are unaffected by the matrix $`Y`$, and that the same is true for the microscopic spectral correlations as long as $`\rho (0)/V`$ remains nonzero .
Most universality proofs have been performed for random matrix ensembles with $`\beta =2`$. However, it is possible to establish relations between the $`\beta =2`$ kernel and the kernels of the other two ensembles with $`\beta =1`$ and $`\beta =4`$ . From these relations and the universality of the $`\beta =2`$ results, one can then infer the universality of the $`\beta =1`$ and $`\beta =4`$ results. Recently, interesting connections between massive correlators and the massless correlators for $`\beta =1`$ and $`\beta =4`$ have been made , and the transition between these ensembles and the $`\beta =2`$ ensemble has been studied .
The microscopic spectral correlations are universally given by RMT only if $`\rho (0)/V>0`$. An interesting question is what happens to the $`R_k`$ if there is some phase transition so that $`\rho (0)/V`$ vanishes. This transition could be triggered by, for example, nonzero temperature or a large number of flavors. Beyond the transition point, there will be a gap in the spectrum, and the smallest eigenvalues beyond this gap are presumably described by the soft-edge results of RMT. What is more interesting, however, is the transition point at which $`\rho (0)/V`$ becomes zero so that there is no gap. This situation has been investigated in two different ways. First, keeping the unitary invariance, one can fine-tune the polynomial in the exponent of $`P(W)`$ to make $`\rho (0)/V`$ just vanish . Second, breaking the unitary invariance, one can choose a “critical” deterministic matrix $`Y`$ so that $`\rho (0)/V=0`$ without a gap . In both cases, one obtains new functional forms for the microscopic spectral correlations and a different scaling with the volume (or matrix dimension). However, the results are not unique; they depend on the way one chooses to approach the $`\rho (0)/V=0`$ limit. It is, therefore, an open question as to whether the microscopic spectral correlations of the QCD Dirac operator at the chiral phase transition agree with one of the results obtained in a random matrix model.
## 4 EFFECTIVE THEORIES AT LOW ENERGIES
### 4.1 Finite-Volume Partition Function
As we have seen before, chiral symmetry is spontaneously broken by the QCD vacuum. According to Goldstone’s theorem, this leads to the appearance of massless modes. The Lagrangian describing the dynamics of these modes can be obtained solely from the symmetries of the QCD partition function. Goldstone modes corresponding to the symmetry-breaking pattern $`\mathrm{SU}_R(N_f)\times \mathrm{SU}_L(N_f)\mathrm{SU}_V(N_f)`$ can be parameterized as
$$U=U_RU_L^1$$
(59)
with $`U_R\mathrm{SU}_R(N_f)`$ and $`U_L\mathrm{SU}_L(N_f)`$. The chiral Lagrangian is constructed by the requirement that it should have the same invariance properties as the QCD Lagrangian. In particular, for $`m=0`$ the Lagrangian should be both Lorentz invariant and invariant under $`\mathrm{SU}_R(N_f)\times \mathrm{SU}_L(N_f)`$. To lowest order in the momenta (or derivatives), the kinetic term in the effective Lagrangian is uniquely given by
$$_{\mathrm{kin}}=\frac{F^2}{4}\mathrm{Tr}_\mu U_\mu U^1.$$
(60)
The parameter $`F`$ is the pion decay constant. The mass term in the QCD Lagrangian breaks the full flavor symmetry. However, the full symmetry can be restored if the mass matrix is transformed as well,
$$MU_LMU_R^1.$$
(61)
We require that the mass term in the effective Lagrangian satisfies this extended symmetry. To lowest order in $`M`$, it is therefore uniquely given by
$$_m=\frac{1}{2}\mathrm{\Sigma }\mathrm{Tr}(MU+M^{}U^1).$$
(62)
For a diagonal mass matrix, the action is therefore minimized by $`U=\mathrm{𝟏}`$. The normalization of the mass term is such that the mass derivative of the partition function according to Eq. (23) is equal to the chiral condensate $`\mathrm{\Sigma }`$. The chiral Lagrangian is valid in the domain
$$m\mathrm{\Lambda }\mathrm{and}p\mathrm{\Lambda },$$
(63)
where $`p`$ is the momentum and $`\mathrm{\Lambda }`$ is a typical hadronic mass scale.
We will study the effective Lagrangian at finite volume in a box of volume $`L^4`$. Then the smallest nonzero-momentum modes are of the order
$$p1/L.$$
(64)
The fluctuations of the zero-momentum modes, the constant fields, are not affected by the kinetic term and are limited only by the mass term. Comparing Eqs. (60) and (62), we see that in a domain where
$$\frac{m\mathrm{\Sigma }}{F^2}\frac{1}{L^2},$$
(65)
the fluctuations of the zero-momentum modes completely dominate the fluctuations of the nonzero-momentum modes. The mass-dependence of the effective partition function is then given by
$$Z(m,\theta )=_{U\mathrm{SU}(N_f)}DUe^{V\mathrm{\Sigma }\mathrm{Re}\mathrm{Tr}Ume^{i\theta /N_f}},$$
(66)
where we have set $`M=m\mathrm{𝟏}`$ and have introduced the $`\theta `$-dependence by the substitution $`mm\mathrm{exp}(i\theta /N_f)`$ just as for the QCD partition function. Combining the conditions (63) and (65), we find that the finite volume partition function is valid in the domain
$$m\frac{1}{\mathrm{\Lambda }L^2}\mathrm{\Lambda }.$$
(67)
The partition function for a given topological charge can be extracted from its $`\theta `$-dependence. Since
$$Z(m,\theta )=\underset{\nu }{}e^{i\nu \theta }Z_\nu (m),$$
(68)
we obtain $`Z_\nu `$ by Fourier inversion. Thus, the finite volume partition function in the sector of topological charge $`\nu `$ is
$$Z_\nu (m)=\frac{1}{2\pi }𝑑\theta e^{i\nu \theta }Z(m,\theta )=_{U\mathrm{U}(N_f)}DU\stackrel{\nu }{det}Ue^{V\mathrm{\Sigma }\mathrm{Re}\mathrm{Tr}mU},$$
(69)
which coincides with Eq. (41). The partition function (69) can be evaluated for arbitrary masses. The integrals were first calculated in the context of one-plaquette lattice QCD models but later rederived by means of Itzykson-Zuber integrals .
### 4.2 Leutwyler-Smilga Sum Rules
For masses in the domain of Eq. (67), we can equate the mass dependence of the QCD partition function and the mass dependence of the effective partition function. Equating the coefficients of the expansion in powers of the mass then gives us the so-called Leutwyler-Smilga sum rules . The QCD partition function can be expanded as
$$Z_\nu ^{\mathrm{QCD}}(m)=\frac{m^{|\nu |}_k^{}(m+i\lambda _k)}{_k^{}i\lambda _k}=m^{|\nu |}\left(1+m^2\underset{\lambda _k>0}{}\frac{1}{\lambda _k^2}+\mathrm{}\right),$$
(70)
where the prime indicates that the product is over nonzero eigenvalues only. The effective partition function can be expanded in a power series in $`m`$ as well. From the U(1) part of the group integral, it is clear that $`Z_\nu ^{\mathrm{eff}}(m)m^{|\nu |}`$. Thus,
$$Z_\nu ^{\mathrm{eff}}=m^{|\nu |}(a_0+a_1m^2+\mathrm{}),$$
(71)
where the coefficients $`a_i`$ are obtained by calculating the group integrals. Let us consider two examples. The simplest example is $`N_f=1`$, for which
$$Z_\nu ^{\mathrm{eff}}=\frac{1}{2\pi }𝑑\theta e^{i\nu \theta }e^{mV\mathrm{\Sigma }\mathrm{cos}\theta }=I_\nu (mV\mathrm{\Sigma }).$$
(72)
The series expansion of the modified Bessel function is given by
$$I_\nu (x)=\frac{1}{\nu !}\left(\frac{x}{2}\right)^\nu +\frac{1}{(\nu +1)!}\left(\frac{x}{2}\right)^{\nu +2}+\mathrm{}$$
(73)
for $`\nu 0`$. We also have $`I_\nu (x)=I_\nu (x)`$. Matching the normalizations of Eqs. (70) and (71) and equating the coefficients of the $`O(m^2)`$ terms, we find
$$\frac{1}{V^2}\underset{\lambda _k>0}{}\frac{1}{\lambda _k^2}=\frac{\mathrm{\Sigma }^2}{4(|\nu |+1)}.$$
(74)
A second example is the case of $`\nu =0`$ and arbitrary $`N_f`$. In this case, we use the group integral
$$DUU_{ij}U_{kl}^{}=\frac{1}{N_f}\delta _{jk}\delta _{il}$$
(75)
to derive the result
$$\frac{1}{V^2}\underset{\lambda _k>0}{}\frac{1}{\lambda _k^2}=\frac{\mathrm{\Sigma }^2}{4N_f}.$$
(76)
Since asymptotically for large $`\lambda `$ the spectral density is proportional to $`V\lambda ^3`$, the sum over the eigenvalues has to be regularized. As discussed above in Sec. 2.3, a finite result is obtained by taking the thermodynamic limit before removing the cutoff.
What do we learn from these sum rules? Since the total number of eigenvalues is of order $`V`$, the obvious interpretation is that the smallest eigenvalues are of order $`1/V`$. Indeed, this is in agreement with the Banks-Casher formula. What is more important, however, is that in the derivation of the sum rules we have relied only on the chiral symmetry of the partition function and its spontaneous breaking by the formation of a chiral condensate. Therefore any theory with the same pattern of chiral symmetry breaking as QCD should obey the same spectral sum rules. In particular, the eigenvalues distributed according to the chRMT introduced in the previous section should obey the same sum rules. In fact, all Leutwyler-Smilga sum rules can be derived systematically from the joint eigenvalue chRMT probability distribution in Eq. (36) using the known Selberg integrals . In addition to these sum rules, it is possible to derive massive sum rules . A relation between sum rules and partially quenched effective theories was considered in Ref. . Spectral sum rules for Ginsparg-Wilson fermions on the lattice were considered in Ref. . Inverse moments of the small Dirac eigenvalues in generalized chiral perturbation theory were derived in Ref. .
We have seen that chRMT allows us to calculate the correlations of the smallest eigenvalues of the Dirac operator on the scale of the average level spacing. This raises the question of whether it is possible to derive the microscopic properties of the eigenvalues from the low-energy effective partition function. The answer is no. The spectral density cannot be derived from the mass dependence of the chiral condensate,
$$\mathrm{\Sigma }(m)=\frac{1}{V}\underset{k}{}\frac{1}{i\lambda _k+m},$$
(77)
because the average contains the fermion determinant with the same mass. As we show in the next section, this problem can be circumvented by the introduction of a spectral mass.
### 4.3 Spectral Mass and Partially Quenched Partition Function
The Dirac spectrum can be obtained from the resolvent
$$\mathrm{\Sigma }(z)=\frac{1}{V}\mathrm{Tr}\frac{1}{𝒟+z},$$
(78)
were the spectral mass $`z`$ is an independent complex variable that does not occur in the average. With purely imaginary eigenvalues the spectral density is given by the discontinuity of the resolvent across the imaginary axis,
$$\frac{\rho (\lambda )}{V}=\underset{ϵ0}{lim}\frac{1}{2\pi }[\mathrm{\Sigma }(i\lambda +ϵ)\mathrm{\Sigma }(i\lambda ϵ)].$$
(79)
This follows immediately upon writing the trace as a sum over the eigenvalues of $`𝒟`$. The generating function for the resolvent can be written down easily ,
$$Z^{\mathrm{sp}}=\frac{det^{N_f}(𝒟+m)det(𝒟+z)}{det(𝒟+z^{})}_{\mathrm{YM}},$$
(80)
where we have explicitly displayed the fermion determinant in the measure. The resolvent is then given by
$$\mathrm{\Sigma }(z)=\frac{1}{V}_zZ^{\mathrm{sp}}(z,z^{})|_{z^{}=z}.$$
(81)
In QCD, this method was first introduced in order to derive the quenched strong coupling expansion . In nuclear physics and condensed matter physics, this method is known as the supersymmetric method or Efetov method for quenched disorder .
The determinants in the numerator can be written as fermionic integrals whereas the determinant in the denominator can be written as a bosonic integral. The partition function (80) is thus invariant under flavor symmetries that mix commuting and anticommuting degrees of freedom. More precisely, for $`m=z=z^{}=0`$ the partition function is invariant under the supergroup $`\mathrm{Gl}_R(N_f+1|1)\times \mathrm{Gl}_L(N_f+1|1)`$. In order to obtain a consistent effective theory, it is essential to extend the unitary symmetry to the general linear group Gl. We expect, as for the QCD vacuum, that chiral symmetry is spontaneously broken to the diagonal subgroup $`\mathrm{Gl}_\mathrm{V}(N_f+1|1)`$. The mass term explicitly breaks the full symmetry to this subgroup as well.
The effective partition function can be derived in the same way as the regular chiral Lagrangian with one complication. To obtain finite integrals, we must make sure that the integration manifold is Riemannian. This is why we have extended the flavor symmetry to the full general linear group. The integration manifold is then given by the maximum Riemannian submanifold of $`\mathrm{Gl}_\mathrm{A}(N_f+1|1)`$. In plain language, this means that we compensate the extra minus sign in the supertrace by complexifying the group parameters. The generating function is thus given by
$$Z^{\mathrm{pq}}=\underset{U\widehat{\mathrm{Gl}}(N_f+1|1)}{}DUe^{{\scriptscriptstyle d^4x\left[{\scriptscriptstyle \frac{F^2}{4}}\mathrm{Str}_\mu U_\mu U^1{\scriptscriptstyle \frac{F^2m_0^2}{12}}\left({\scriptscriptstyle \frac{\sqrt{2}\mathrm{\Phi }_0}{F}}\theta \right)^2+{\scriptscriptstyle \frac{\mathrm{\Sigma }}{2}}\mathrm{Str}(MU+M^{}U^1)\right]}},$$
where $`i\sqrt{2}\mathrm{\Phi }_0/F=\mathrm{Str}\mathrm{log}U`$. This partition function is also known as the partially quenched effective partition function. It has been used to obtain a better understanding of quenched lattice QCD results . The hat on Gl denotes the maximum Riemannian submanifold of $`\mathrm{Gl}(N_f+1|1)`$, and Str stands for the supertrace. An example of an explicit parameterization of $`U`$ will be given in the next subsection. The parameter $`m_0^2`$ is proportional to the topological susceptibility and results in a mass for the singlet channel that does not vanish in the chiral limit.
The perturbative formulation of this partition function was first given by Bernard and Golterman , and a perturbative calculation of the dependence of the condensate on the spectral mass was first performed in Ref. . In that case, it is acceptable to ignore the convergence properties of the effective partition function and integrate over the unitary supergroup instead.
An alternative generating function would be obtained by replacing the ratio of the two determinants by $`det^n(𝒟+z)`$ and putting $`n0`$ after having calculated the resolvent. This procedure, known as the replica trick , successfully reproduces the asymptotic expansions of the spectral correlation functions . There have been recent claims that it is possible to derive truly nonperturbative results by means of the replica trick . We do not believe that this can be done in general . As mentioned, because of the compact/noncompact structure of the final answer, the supersymmetric formulation is the natural approach to this problem.
### 4.4 Domains of the Partially Quenched Effective Theory
The Goldstone fields can be written as
$$U=e^{i\sqrt{2}\mathrm{\Pi }/F}.$$
(82)
To second order in the fields, the effective Lagrangian in momentum space is given by
$$=\frac{1}{V}\underset{a}{}\underset{k}{}(k^2+M_a^2)\mathrm{\Pi }_a^2(k),$$
(83)
where the sum is over the momenta in a box of length $`L`$ (including the zero-momentum state) and the masses of the Goldstone bosons are denoted by $`M_a`$. The sum over $`a`$ is over the different Goldstone modes. In addition to the usual Goldstone modes with mass
$$M_{mm}^2=\frac{2m\mathrm{\Sigma }}{F^2},$$
(84)
there are both fermionic and bosonic Goldstone bosons with mass
$$M_{mz}^2=\frac{(m+z)\mathrm{\Sigma }}{F^2}$$
(85)
and fermionic and bosonic Goldstone bosons with mass
$$M_{zz}^2=\frac{2z\mathrm{\Sigma }}{F^2}.$$
(86)
In the Lagrangian (83), one can distinguish the zero-momentum modes from the nonzero-momentum modes. The magnitude of the fluctuations of the nonzero-momentum modes is of order $`1/(k^2+M_a^2)`$, whereas the magnitude of the fluctuations of the zero-momentum modes is of order $`1/M_a^2`$. Since the smallest nonzero momenta are of order $`1/L`$, the fluctuations of the zero-momentum modes are dominant if
$$M_a^2\frac{1}{L^2}.$$
(87)
This means that the Compton wavelength of the “pion” is much larger than the size of the box. For nonzero quark masses of order $`O(L^0)`$, the inequality (87) is never satisfied. However, the spectral mass $`z`$ is a free parameter, and the inequality can be rewritten as
$$zE_c\frac{F^2}{\mathrm{\Sigma }L^2}.$$
(88)
We sometimes refer to the quantity $`E_c`$ as the Thouless energy for reasons that will become clear in Sec. 7.3. In the domain of Eq. (88), the dominant contributions to the resolvent are from the zero-momentum modes. Thus, the partially quenched effective partition function (4.3) can be reduced to the partition function in the zero-momentum sector ,
$$Z=_{U\widehat{\mathrm{Gl}}(N_f+1|1)}DUe^{\frac{1}{2}V\mathrm{\Sigma }\mathrm{Str}(MU+M^{}U^1)}.$$
(89)
The microscopic spectral density and the spectral correlations computed from Eq. (89) are identical to those obtained in chRMT . Let us show this by an explicit calculation of the dependence of the chiral condensate on the spectral mass $`z`$. To cover the complete range (88), the integral over $`U`$ has to be done nonperturbatively. This is a straightforward superintegral that can be evaluated using standard methods. The calculation of the integration measure requires an explicit parameterization of the integration manifold. For example, in the quenched limit, $`N_f=0`$, a possible choice is
$$U=\left(\begin{array}{cc}e^{i\varphi }& \alpha \\ \beta & e^s\end{array}\right)$$
(90)
with $`\alpha `$ and $`\beta `$ Grassmann variables, $`\varphi [0,2\pi ]`$, and $`s(\mathrm{},\mathrm{})`$. A characteristic feature of the integration manifold is that it consists of a compact and a noncompact component. The valence-quark mass dependence of this partition function and its generalization to arbitrary topological charge coincide with the valence quark mass dependence of the chUE partition function. We merely quote the final result for the resolvent for $`N_f`$ massless flavors in the sector of topological charge $`\nu `$ ,
$$\frac{\mathrm{\Sigma }(u)}{\mathrm{\Sigma }}=u\left[I_a(u)K_a(u)+I_{a+1}(u)K_{a1}(u)\right]+\frac{|\nu |}{u},$$
(91)
where $`a=N_f+|\nu |`$ and $`u=zV\mathrm{\Sigma }`$. The compact/noncompact symmetries are reflected in the appearance of the $`I_a/K_a`$-Bessel functions and are thus a natural ingredient of the underlying integration manifold. This result was first obtained from chRMT by integrating the microscopic spectral density in Eq. (57) as follows,
$$\frac{\mathrm{\Sigma }(u)}{\mathrm{\Sigma }}=_0^{\mathrm{}}𝑑\zeta \frac{2u}{\zeta ^2+u^2}\rho _s(\zeta ).$$
(92)
Alternatively, the microscopic spectral density can be obtained by taking the discontinuity of $`\mathrm{\Sigma }(u)`$ according to Eq. (79). By integrating back, we find that the mass dependence of the zero-momentum partially quenched partition function coincides with that of the chRMT partition function for $`\beta =2`$.
Figure 2 presents a schematic picture of the different domains in the Dirac spectrum. There are other non-QCD theories that can be reduced to the same partially quenched partition function. Probably the best known example is the instanton liquid model of QCD . Another example, more closely related to disordered condensed matter systems, is the random-flux model , which is in essence quenched lattice QCD with Kogut-Susskind fermions but without the phase factors due to the $`\gamma `$-matrices. Related examples are so-called two-sublattice models with disorder , and disordered lattice models with the chiral and flavor symmetries of QCD .
In the range
$$E_cz\mathrm{\Lambda },$$
(93)
the effective action (4.3) is still valid but chiral random matrix theory no longer applies. In this range, the $`U`$ fields can be expanded to second order in the pion fields, and the spectral mass dependence of the chiral condensate can be obtained to one loop order. The spectral density follows by taking the discontinuity of $`\mathrm{\Sigma }(z)`$ across the imaginary axis. This calculation can be performed both for the partition function of Eq. (4.3) and the effective partition functions for $`\beta =1`$ and $`\beta =4`$. The result for the slope of the spectral density at zero, valid for $`N_f2`$, reads
$$\frac{\rho ^{}(0)}{\rho (0)}=\frac{(N_f2)(N_f+\beta )}{16\pi \beta N_f}\frac{\mathrm{\Sigma }}{F^4}.$$
(94)
The result for $`\beta =2`$ was first derived by Smilga and Stern from the scalar susceptibility in standard chiral perturbation theory. We emphasize that this is an exact result for the QCD Dirac spectrum that is valid in the thermodynamic limit. The vanishing of the slope for $`N_f=2`$ has been confirmed by instanton liquid simulations .
The two-point correlation function can also be obtained from a supersymmetric generating function with, in this case, two different spectral masses and two bosonic superpartners . Again the zero-momentum contribution coincides with the chUE result. Let us consider the perturbative expansion of the two-point correlation function. Generalizing Eq. (79), we obtain a relation between the two-point level correlation function and the pion susceptibility ,
$$\rho (\lambda )\rho (\lambda ^{})^C=\frac{1}{4\pi ^2}\frac{\mathrm{\Sigma }^2}{F^4}\mathrm{Disc}|_{z=i\lambda ,z^{}=i\lambda ^{}}\underset{q}{}\frac{1}{(q^2+M^2)^2},$$
(95)
where the meson mass is given by $`M^2=(\sqrt{z^2}+\sqrt{z^2})\mathrm{\Sigma }/F^2`$ and the superscript $`C`$ denotes the connected part of the two-point function. By taking the discontinuities of the zero-momentum contribution, we find the two-point correlation function
$$\rho (\lambda )\rho (\lambda ^{})^C\frac{1}{2\pi ^2}\left[\frac{1}{(\lambda \lambda ^{})^2}+\frac{1}{(\lambda +\lambda ^{})^2}\right],$$
(96)
which is the correct asymptotic result for the chUE.
### 4.5 QCD in Three Euclidean Dimensions
QCD in three Euclidean dimensions can be analyzed in much the same way as discussed above. The main difference is the absence of the $`\mathrm{U}_\mathrm{A}(1)`$ symmetry and the absence of instantons. In terms of the Dirac spectrum, the eigenvalues do not occur in pairs $`\pm i\lambda `$, and strictly zero eigenvalues are absent. The Dirac spectrum of this theory on the microscopic scale can be analyzed along the same lines as the four-dimensional theory. The microscopic spectral density and the Leutwyler-Smilga sum rules have been derived , the low-energy effective theory has been identified , the mass-dependent microscopic spectral density has been found , universality has been studied , and lattice QCD studies have confirmed the theoretical analysis .
## 5 UNIVERSAL PROPERTIES OF THE LATTICE QCD DIRAC SPECTRUM
### 5.1 Unfolding
In order to compare eigenvalues of any physical system with RMT results, it is necessary first to “unfold” the empirical spectrum (see e.g. Ref. ). The spacing $`\mathrm{\Delta }`$ of the eigenvalues is related to the spectral density by
$$\mathrm{\Delta }(E)=1/\rho (E).$$
(97)
Unfolding is a local rescaling of the energy scale so that the mean level spacing of the unfolded eigenvalues is equal to unity throughout the spectrum. This is achieved by splitting the spectral density in an average density, $`\overline{\rho }(E)`$, and a fluctuating piece,
$$\rho (E)=\overline{\rho }(E)+\rho _{\mathrm{fl}}(E).$$
(98)
The average spectral density can be obtained in different ways. In some cases, it can be computed analytically using semiclassical arguments. However, in most cases $`\overline{\rho }(E)`$ is obtained by averaging over many level spacings. There are two essentially different procedures to do this, spectral averaging and ensemble averaging. Spectral averaging is appropriate if only one or a few spectra are available. The average level spacing at $`E`$ is then obtained by averaging over many level spacings around $`E`$. Ensemble averaging is appropriate if there is a large ensemble of spectra all drawn from the same statistical distribution. In this case, the average spacing can be obtained by averaging the spacing at $`E`$ over each member of the ensemble. In both cases, the separation of the spectral density into an average and a fluctuating piece requires a separation of scales that is typically achieved only in the thermodynamic limit or in the semiclassical limit. The unfolded spectrum, $`\{x_n\}`$, is then obtained from the average spectral density (equal to the inverse average level spacing) by
$$x_n=_{\mathrm{}}^{E_n}\overline{\rho }(\lambda )𝑑\lambda ,$$
(99)
where $`\{E_n\}`$ is the original sequence of eigenvalues. The difference between the two procedures is that spectral averaging yields $`\overline{\rho }(E)`$ separately for each spectrum whereas ensemble averaging yields a single $`\overline{\rho }(E)`$ for all spectra. In both cases, the average spectral density of the unfolded spectrum becomes $`\overline{\rho }(x)=1`$, so that the sequence $`\{x_n\}`$ has an average level spacing equal to unity. The correlations of the levels are always calculated for the unfolded spectrum.
The two procedures to calculate the average spectral density do not necessarily give the same spectral correlations at long distances. The equivalence of the two is known as spectral ergodicity. This property has been shown analytically for several random matrix ensembles . In general, ensemble averaging results in stronger level fluctuations than does spectral averaging . In practice a mixture of both methods is often useful. For example, one may calculate the average spectral density by ensemble averaging but, in order to get better statistics, calculate the correlation functions by both ensemble averaging and spectral averaging. Spectral averaging requires that the statistical properties of the eigenvalues be stationary over the spectrum, which is not the case for many systems and has to be checked each time.
### 5.2 Lattice Tests of Chiral Random Matrix Theory
We have presented a wealth of analytical evidence supporting the statement that the local spectral correlations of the Dirac operator are described by universal functions. Although the Dirac spectrum is not directly observable in experiments, we can compare the predictions of chiral RMT to numerical data obtained by lattice gauge simulations. This is the subject of this section.
This review does not discuss the global spectral density of the QCD Dirac operator. Most lattice results have large finite-size artifacts. Instanton simulations suggest, contrary to random matrix results, that the fermion determinant has a significant effect on the global Dirac spectrum .
#### 5.2.1 BRIEF INTRODUCTION TO LATTICE QCD
QCD is a renormalizable quantum field theory that must be regularized. This can be done by formulating the theory on a discrete lattice with lattice spacing $`a`$ . The largest momentum is then $`\pi /a`$. The quark fields live on the sites and the gauge fields on the links of the lattice. If the lattice is finite, one can simulate the theory on a computer. This can be done efficiently only in Euclidean space, where the gluonic weight function is $`\mathrm{exp}(S_{\mathrm{YM}})`$ with $`S_{\mathrm{YM}}`$ real. The discretized form of the Yang-Mills action $`S_{\mathrm{YM}}`$ is the Wilson action $`S_\mathrm{W}`$. The full weight function of QCD also contains the fermion determinants, which can be expressed in terms of the gauge fields. Observables are computed by generating gauge field configurations in a Monte Carlo update procedure and averaging an observable over many configurations. Since the inclusion of the fermion determinants is very time consuming, it is common to use only the gluonic part of the weight function in the Monte Carlo updates. This is called the quenched approximation, which corresponds to the limit $`N_f=0`$ or, equivalently, to the limit of infinitely heavy sea quarks. To make contact with continuum physics, the results of lattice simulations must be extrapolated to infinite lattice size (the thermodynamic limit) and to zero lattice spacing (the continuum limit).
Lattice simulations are the main source of nonperturbative information about QCD. Unfortunately, the naive discretization of fermions on a lattice leads to the so-called doubling problem: The quark propagator has poles at each corner of the Brillouin zone, which gives rise to a total of $`2^d`$ species in $`d`$ dimensions. The unwanted $`2^d1`$ species can be eliminated by adding to the Dirac operator an additional term, the Wilson term, which removes the doublers in the continuum limit. However, this term breaks chiral symmetry explicitly. Another possibility is the use of staggered (or Kogut-Susskind) fermions where one has only one spinor component per lattice site. This maintains a residual chiral symmetry but only partially reduces the number of species to $`2^{d/2}`$ in the $`a0`$ limit. A no-go theorem by Nielsen and Ninomiya states that it is not possible simultaneously to solve the doubling problem and have exact chiral symmetry on the lattice with a local action.
Fortunately, there is a way around this theorem: the remnant chiral symmetry condition of Ginsparg and Wilson ,
$$𝒟\gamma _5+\gamma _5𝒟=2a𝒟\gamma _5R𝒟,$$
(100)
where $`R`$ is a spatially local operator that is trivial in Dirac space. In contrast to $`\{𝒟,\gamma _5\}=0`$, chiral symmetry is not exact but is recovered only in the continuum limit $`a0`$. Thus, with Dirac operators satisfying Eq. (100), it is possible to get rid of the doublers even at finite lattice spacing without violating the Nielsen-Ninomiya theorem. Recently, several solutions of Eq. (100) have been found in the overlap formalism , in the domain-wall formulation in five dimensions , and in the perfect action approach .
As mentioned in Sec. 2.2.3, there may be cases in which the anti-unitary symmetries of the various lattice discretizations of the Dirac operator differ from those of the continuum operator. In particular, this is true for staggered fermions. In SU(2) color, staggered fermions are in the symmetry class of the chSE, whereas continuum fermions are in the chOE symmetry class. Staggered fermions in the adjoint representation of the gauge group (for any $`N_c`$) have the symmetries of the chOE, whereas in the continuum limit the symmetries are those of the chSE. The Wilson Dirac operator $`𝒟_\mathrm{W}`$ does not anticommute with $`\gamma _5`$ and is therefore not described by any of the chiral ensembles. However, the Hermitian operator $`\gamma _5𝒟_\mathrm{W}`$ has the same anti-unitary symmetries as the continuum Dirac operator and is described by the corresponding nonchiral ensembles (see also Ref. ). Finally, Dirac operators obeying the Ginsparg-Wilson condition, Eq. (100), have the same anti-unitary symmetries as the continuum Dirac operator, but the eigenvalues are located on the complex unit circle .
#### 5.2.2 BULK CORRELATIONS
To measure the bulk spectral correlations, one needs all eigenvalues of the Dirac operator, which is represented on the lattice by a finite sparse matrix whose dimension is proportional to the lattice volume and to $`N_c`$. The eigenvalues of such a matrix can be obtained using special algorithms, e.g. by the Cullum-Willoughby version of the Lanczos algorithm . On the lattice, this was first done by Kalkreuter . The numerical effort of this method scales with the square of the matrix dimension. For example, in SU(3) on a $`10^4`$ lattice, the Dirac matrix has 15,000 distinct positive eigenvalues, which can be computed on a typical workstation in about 40 minutes. There are exact sum rules for the sum of the squares of all Dirac eigenvalues, which can be used to check the numerical accuracy. Because the number of eigenvalues per configuration is very large and because the ensemble average can be replaced by a spectral average under the assumption of spectral ergodicity and stationarity, one needs only a few configurations to construct the bulk spectral correlations with great accuracy.
The bulk correlations have been measured in a number of lattice studies for all three symmetry classes and in instanton liquid simulations for $`\beta =2`$. As mentioned above, they are insensitive to the number of flavors and the topological charge. In all cases, excellent agreement with the predictions of the appropriate random matrix ensemble was obtained (see e.g. Fig. 3).
The agreement is perfect not only in the strong-coupling regime but also at weak coupling. In fact, the bulk spectral correlations are given by RMT even in the deconfinement phase , indicating that the gauge fields retain a sufficient degree of randomness in this region of the phase diagram.
Spectral ergodicity was investigated in Ref. , and the equivalent of a Thouless energy was found for ensemble averaging, whereas spectral averaging resulted in complete agreement with RMT correlations over distances as long as several hundred average level spacings.
#### 5.2.3 MICROSCOPIC CORRELATIONS
Because only the lowest eigenvalues contribute to the microscopic spectral correlations, a large number of statistically independent spectra are necessary. In contrast to the bulk correlations, the microscopic correlations are sensitive to the number of flavors and the topological charge.
The RMT predictions do not contain an energy scale. In order to make comparisons with lattice data, one needs to determine the energy scale $`1/V\mathrm{\Sigma }`$ to be used in the RMT expressions, see e.g. Eq. (27). This can be done by extracting $`\rho (0)`$ from the data and applying the Banks-Casher relation, Eq. (22). Because this procedure makes no reference to RMT, the comparisons of lattice data with RMT are parameter-free.
The first numerical results for the microscopic spectral density $`\rho _s`$ were obtained using instanton liquid configurations for $`N_c=2,3`$ and $`N_f=0,1,2`$ , and the expected agreement with the corresponding random matrix predictions was found. The microscopic spectral density was first observed on the lattice via the dependence of the chiral condensate on a valence quark mass $`m_v`$, as studied by the Columbia group . Figure 4
shows $`\mathrm{\Sigma }(m_v)/\mathrm{\Sigma }`$ as a function of $`m_vV\mathrm{\Sigma }`$ for different values of the coupling constant . The data for different coupling strengths in the broken phase fall on a single curve and agree with the chRMT result for $`\mathrm{\Sigma }(m_v)`$ for $`N_f=|\nu |=0`$, see Eq. (91). This figure also shows a deviation from the chRMT predictions at a scale that is of the order of the Thouless energy given in Eq. (88). The spectral mass dependence of the chiral condensate has recently been studied for different symmetry classes and different types of fermions, and a similar quality of agreement has been found .
The microscopic spectral correlations have also been investigated in great detail on the lattice , and the random matrix predictions were confirmed with very high accuracy for all three symmetry classes. A typical example is shown in Fig. 5.
Although most of these simulations have been performed in the quenched approximation, it is also possible to include dynamical fermions. Deviations from the $`N_f=0`$ results are observed only if the sea quarks are very light, with masses on the scale $`1/V\mathrm{\Sigma }`$. Otherwise, the mass term in Eq. (1) dominates the small eigenvalues in the factors of $`(i\lambda _n+m_f)`$. Analytical results for the microscopic spectral correlations in the presence of sea quarks with mass $`1/V\mathrm{\Sigma }`$ were obtained in Refs. . Lattice results in this regime agree very well with the corresponding RMT predictions . In the Schwinger model, it is numerically feasible to consider massless sea quarks, and the lattice data are again well described by RMT .
Recall that the energy scale for the small eigenvalues is $`1/V\mathrm{\Sigma }`$. Given the agreement of the microscopic spectral quantities with RMT, this means that RMT can be used to determine the infinite-volume chiral condensate $`\mathrm{\Sigma }`$ . This is most easily done by fitting the numerically determined distribution of the smallest eigenvalue to the RMT result. Since the lattice volume is known, this immediately yields $`\mathrm{\Sigma }`$.
#### 5.2.4 TOPOLOGY
The random matrix predictions for the microscopic spectral correlations depend on the number of zero modes of the Dirac operator and, thus, on the topological charge $`\nu `$ of the gauge field configurations. Thus, one should sort the configurations according to their values of $`\nu `$ and make the comparison with RMT separately in sectors of fixed topological charge.
The results presented in the previous section were obtained for the staggered Dirac operator. If one is interested in topological properties, this operator has problems. The zero modes that one would obtain in the continuum limit are shifted at finite lattice spacing $`a`$ by an amount proportional to $`a^2`$ . For reasonable simulation parameters, this amount is larger than the level spacing near zero so that the would-be zero modes are completely mixed with the nonzero modes. Thus, one should expect all spectra of the staggered Dirac operator, even if computed from gauge field configurations with nonzero topological charge, to be described by the RMT results for $`\nu =0`$. This was indeed observed in the early data and more recently confirmed in a detailed study (though only in strong coupling) . As the continuum limit is approached, the would-be zero modes should move toward $`\lambda =0`$ and eventually separate completely from the nonzero modes. This effect has recently been observed in the Schwinger model in two dimensions .
Fortunately, Dirac operators satisfying the Ginsparg-Wilson condition (100) do not have this problem. The overlap operator has exact zero modes even at finite lattice spacing, and the microscopic spectral correlations in sectors of fixed $`\nu `$ are in perfect agreement with the corresponding random matrix predictions . Figure 6 shows an example.
#### 5.2.5 CHIRAL PHASE TRANSITION
As mentioned in Sec. 3.5, the microscopic correlations are given by chiral RMT only if $`\rho (0)/V>0`$, i.e. if chiral symmetry is spontaneously broken. Chiral symmetry is restored above a critical temperature $`T_c`$. The way in which $`\rho (0)/V`$ approaches zero determines the order of the chiral phase transition and, if it is second-order, the associated critical exponents. Lattice investigations of these questions are very difficult because at finite volume there are no sharp phase transitions. Simulations are plagued by finite-size effects, critical slowing down, and other problems. It is therefore of great interest to investigate the fate of the small Dirac eigenvalues at $`T=T_c`$ analytically in random matrix models, as was discussed at the end of Sec. 3.5. There are two recent lattice studies of this question . Some of the theoretical expectations were confirmed qualitatively, e.g. the gap in the eigenvalue distribution above $`T_c`$ and the agreement of the eigenvalue distribution in this region with the soft-edge predictions of RMT. However, further studies are needed to clarify the situation at $`T=T_c`$. We are looking forward to upcoming work in this area.
### 5.3 The Thouless Energy and Beyond
As discussed above, the Dirac spectrum is described by chRMT — or, equivalently, by the zero-mode approximation of the low-energy effective theory, — only below the so-called Thouless energy $`E_c`$. The theoretical prediction for this quantity ,
$$E_c\frac{F^2}{\mathrm{\Sigma }\sqrt{V}},$$
(101)
has been confirmed quantitatively in the instanton liquid model and on the lattice . To compute the Dirac spectrum beyond the Thouless energy, the calculation must include the kinetic terms in the low-energy effective theory, as discussed in Sec. 4. The results of such an analysis can again be tested by lattice simulations. This has been done for staggered fermions . Note that in this case, the effective low-energy theory must be modified to take into account the symmetries of staggered fermions at finite lattice spacing. A convenient quantity to test the theory is the disconnected scalar susceptibility, defined by
$$\chi ^{\mathrm{disc}}(m)=\frac{1}{N}\underset{k,l=1}{\overset{N}{}}\frac{1}{(i\lambda _k+m)(i\lambda _l+m)}\frac{1}{N}\underset{k=1}{\overset{N}{}}\frac{1}{i\lambda _k+m}^2.$$
(102)
This quantity can also be computed in chRMT and in chiral perturbation theory (chPT). The result of chRMT is expected to describe the data up to $`E_c`$. The result of chPT should describe the data beyond $`E_c`$ but is expected to break down on the scale of the smallest eigenvalue, i.e. for $`m1/V\mathrm{\Sigma }`$. These expectations are confirmed by the lattice analysis (see Fig. 7 for an example).
## 6 SCHEMATIC MODELS
We wish to emphasize again that there are two different types of applications of RMT to physical problems. RMT may be applied as an exact theory for correlations of eigenvalues, as discussed in the first part of this review, or it may serve as a schematic model for chaos and disorder in physical systems. Below, we introduce a random matrix model for the chiral phase transition at nonzero chemical potential and temperature. Random matrix theories as schematic models are mostly used for a description of nonuniversal phenomena. This does not exclude that they describe universal fluctuations of the eigenvalues as well. Typically, a random matrix model describes both universal and nonuniversal properties. Well-known examples are the Anderson model for localization phenomena and the use of random matrix theory in quantum gravity .
We argued in Sec. 1.1 that chiral symmetry breaking can be understood in terms of the stiffness of the Dirac spectrum resulting from interactions of the strong color force. Because spectral stiffness is a characteristic feature of RMT, it is natural to describe the chiral phase transition in terms of a random matrix model. Several different types of schematic random matrix models have been introduced. Here we discuss models for the chiral phase transition at nonzero temperature $`T`$, for the chiral phase transition at nonzero chemical potential $`\mu `$, and for the phase diagram of QCD. We also discuss random matrix models for different types of fermions at $`\mu >0`$.
### 6.1 Chiral Random Matrix Models for the Chiral Phase Transition at Nonzero Temperature
The original motivation for introducing chiral random matrix models was to obtain a better understanding of the QCD Dirac spectrum for temperatures around the critical temperature for the chiral phase transition . The idea is to split the Dirac operator into the time derivative and a remainder that will be replaced by a chRMT,
$$𝒟=\gamma _0_0+R.$$
(103)
In a chiral basis with time dependence given by $`\mathrm{exp}[i(2n+1)\pi T\tau ]`$, where $`\tau `$ is the Euclidean time, the first term in $`𝒟`$ is diagonal and is given by a direct sum of Matsubara frequencies $`\omega _n=(2n+1)\pi T`$. The simplest model for $`𝒟`$ is obtained by replacing the diagonal matrix with positive Matsubara frequencies by the identity matrix times an effective Matsubara frequency $`t`$, and replacing the diagonal matrix with negative Matsubara frequencies by the opposite effective Matsubara frequency. After a suitable basis change, this model can be written (in the sector of zero topological charge) as
$$Z_{N_f}(m)=DW\stackrel{N_f}{det}(𝒟+m)e^{\frac{1}{2}N\mathrm{\Sigma }^2\mathrm{Tr}(W^{}W)},$$
(104)
where
$$𝒟=\left(\begin{array}{cc}0& iW+it\\ iW^{}+it& 0\end{array}\right).$$
(105)
We restrict ourselves to the unitary case ($`\beta =2`$) with a complex matrix $`W`$ of dimension $`N/2`$. The integral is over the real and imaginary parts of the matrix elements of $`W`$. For simplicity, we consider only a Gaussian probability distribution. The normalization is such that the parameter $`\mathrm{\Sigma }`$ is equal to the magnitude of the chiral condensate at zero temperature,
$$\mathrm{\Sigma }=\underset{m0}{lim}\underset{N\mathrm{}}{lim}\frac{1}{NN_f}_m\mathrm{log}Z_{N_f}(m)|_{t=0}.$$
(106)
It can be shown that the effect of the fermion determinant on the macroscopic spectral density is subleading in $`N_f/N`$, so that $`\rho (\lambda )`$ can be calculated in the quenched limit. By expanding the resolvent of the Dirac operator, defined by
$$G(z)=\frac{1}{N}\mathrm{Tr}\frac{1}{z𝒟},$$
(107)
in a geometric series of the random matrix, it can be shown that $`g=G/\mathrm{\Sigma }^2`$ satisfies a cubic equation ,
$$g^32zg^2+g(z^2t^2+1/\mathrm{\Sigma }^2)z/\mathrm{\Sigma }^2=0.$$
(108)
A variant of the method by which this equation was derived in Refs. is sometimes referred to as the Blue’s function method . The spectral density, obtained from the discontinuity of the resolvent, is shown in Fig. 8. We observe that chiral symmetry is broken up to $`t=1/\mathrm{\Sigma }`$ with a chiral condensate given by
$$G(T,z0)=\mathrm{\Sigma }\sqrt{1(\mathrm{\Sigma }t)^2}.$$
(109)
Above this temperature, the spectrum splits into two disconnected regions.
Since this model has no spacetime dependence, it should not be a surprise that all critical exponents are given by their mean field values. From Eq. (109), we see that $`\beta =1/2`$. For $`t1/\mathrm{\Sigma }`$, we find
$$Gz^{1/3},$$
(110)
resulting in another critical exponent $`\delta =3`$. The spectral density at the critical point follows by taking the discontinuity of $`G`$ across the imaginary axis, see Eq. (79), and is thus given by
$$\rho (\lambda )=\frac{N}{\pi }\lambda ^{1/3}.$$
(111)
The unfolded eigenvalues in this case are given by
$$x_k=\frac{3N}{4\pi }\lambda _k^{4/3},$$
(112)
and a nontrivial scaling limit at $`t=1/\mathrm{\Sigma }`$ is obtained by introducing the microscopic variable
$$z=\lambda N^{3/4}.$$
(113)
Starting from the exact analytical result derived in Ref. , it is possible to take the limit $`N\mathrm{}`$ with the scaling relation (113) and to obtain an analytical result for the microscopic spectral density at the critical point . We emphasize that this result is based on mean field critical exponents and is thus not applicable to QCD with nontrivial critical exponents (see also the discussions in Secs. 3.5 and 5.2.5).
As a last application of the partition function (104), we mention the explanation of the observation that the temperature at which chiral symmetry is restored is higher for gauge field configurations with a nonzero $`Z_3`$-phase than for gauge field configurations with zero $`Z_3`$-phase. The explanation in terms of the random matrix model (104) is that the lowest Matsubara frequency for the nonzero $`Z_3`$-phases is shifted to a lower value by the phase of the Polyakov loop, and thus the critical temperature is higher .
### 6.2 Chiral Random Matrix Models at Nonzero Chemical Potential
#### 6.2.1 QCD PARTITION FUNCTION AT NONZERO CHEMICAL POTENTIAL
The QCD partition function at nonzero temperature $`T`$ and chemical potential $`\mu `$ is given by
$$Z(m,\mu ,T)=\mathrm{Tr}e^{\frac{H_{\mathrm{QCD}}\mu N}{T}}=\underset{\alpha }{}e^{\frac{E_\alpha \mu N_\alpha }{T}},$$
(114)
where $`H_{\mathrm{QCD}}`$ is the Hamiltonian of QCD with eigenvalues $`E_\alpha `$ and $`N`$ is the quark number operator. At zero temperature, only the states with $`E_\alpha /N_\alpha <\mu `$ contribute to the partition function. We thus expect that the partition function is independent of $`\mu `$ below a critical chemical potential $`\mu _c`$ given by the lightest baryon mass per unit quark number. Therefore, for $`\mu <\mu _cm_N/3`$ (where $`m_N`$ is the nucleon mass), the baryon density remains zero and the chiral condensate is constant.
The quark chemical potential appears in the Lagrangian in the form $`\mu \psi ^{}\psi =\overline{\psi }(\mu \gamma _0)\psi `$ and is therefore introduced in the Dirac operator by the substitution
$$_0_0+\mu .$$
(115)
This substitution destroys the Hermiticity properties of the QCD Dirac operator, and the resulting complex phase of the fermion determinant makes Monte Carlo simulations impossible. Because of the success of the quenched approximation in lattice QCD simulations at $`\mu =0`$, it is tempting to ignore the fermion determinant in this case. However, it was shown that the critical chemical potential for quenched simulations is determined by the pion mass rather than the nucleon mass. Obviously, the basic physics of the problem is not visible in the quenched approximation. An analytical understanding of this problem was first obtained by means of a random matrix model for the Dirac operator of QCD at finite density .
#### 6.2.2 A RANDOM MATRIX MODEL FOR QCD AT FINITE DENSITY
A random matrix model for QCD at nonzero chemical potential model is obtained by writing the Dirac operator as
$$𝒟(\mu )=\mu \gamma _0+$$
(116)
and replacing $``$ by a chiral random matrix ensemble. The partition function is thus given by
$$Z_{N_f}(m)=DW\stackrel{N_f}{det}(𝒟+m)e^{\frac{1}{2}N\mathrm{\Sigma }^2\mathrm{Tr}(W^{}W)}$$
(117)
with a Dirac operator given by
$$𝒟=\left(\begin{array}{cc}0& iW+\mu \\ iW^{}+\mu & 0\end{array}\right).$$
(118)
For QCD with three or more colors (chUE) this Dirac operator has no Hermiticity properties. Therefore, the fermion determinant in Eq. (117) has a complex phase, and the partition function cannot be simulated by Monte Carlo methods.
The partition function with Dirac operator (118) can be rewritten in terms of a $`\sigma `$-model. For $`N_f=1`$, the partition function is particularly simple,
$$Z(\mu )=𝑑\sigma 𝑑\sigma ^{}[(\sigma +m)(\sigma ^{}+m)\mu ^2]^ne^{n|\sigma |^2},$$
(119)
where we have set $`\mathrm{\Sigma }=1`$ for ease of notation. For $`n=N/2\mathrm{}`$ the integrals can be evaluated by a saddle-point approximation. In the chiral limit, we find that
$$\overline{\sigma }=\{\begin{array}{ccccc}0& \mathrm{for}& \mu >\mu _c& & Z=\mu ^{2n}\hfill \\ \sqrt{1+\mu ^2}& \mathrm{for}& \mu <\mu _c& & Z=e^{n(\mu ^2+1)}\hfill \end{array}$$
(120)
with $`\mu _c`$ given by the point where the two partition functions are equal, i.e. $`\mu _c^2=\mathrm{exp}(1\mu _c^2)`$, which is solved by $`\mu _c0.53`$. The vacuum properties of this partition function depend on $`\mu `$ for $`\mu \mu _c`$. The reason for this unphysical result is discussed below in Sec. 6.2.4. As was demonstrated in Ref. for $`\beta =2`$, this model does not show diquark condensation. However, a random matrix model that does show diquark condensation was formulated in Ref. .
The QCD partition function (117) is a polynomial in $`m`$ and $`\mu `$ with coefficients that can be obtained analytically . Its properties can be analyzed by means of its zeros in the complex chemical potential plane and the complex mass plane. Figure 9 shows a first-order phase transition at $`\mu 0.53`$. In the chirally restored phase, the cut on the imaginary $`m`$ axis is no longer present. The figure also shows the points (stars) where two solutions of the saddle-point equations of the $`\sigma `$-model coincide. Indeed, the cuts end on these branch points.
In lattice QCD at finite density, the zeros of the QCD function in the complex $`\mu `$ plane can be studied by means of the Glasgow method . This method has been analyzed in terms of the above random matrix model with the conclusion that it requires an exponentially large number of gauge field configurations to obtain statistically significant results.
#### 6.2.3 ZERO TEMPERATURE LIMIT
The random matrix model discussed in the previous section has attracted a great deal of interest, and more sophisticated versions of this model have been proposed . One problem of the model (117) is that the chiral condensate is $`\mu `$-dependent below the critical value of the chemical potential. The correct zero temperature limit, with a chiral condensate that remains constant up to the critical value of the chemical potential, is obtained if all Matsubara frequencies are taken into account. We illustrate this with a one-dimensional lattice QCD model.
Writing the Dirac operator as
$$𝒟(\mu )=\gamma _0(_0+\mu )+R,$$
(121)
the $`\mu `$-dependence of the fermion determinant can be gauged away by means of a time-dependent gauge transformation, i.e.
$$\gamma _0(_0+\mu )+R=e^{\mu \tau }\left[\gamma _0_0+R\right]e^{\mu \tau }.$$
(122)
The $`\mu `$-dependence is now in the boundary conditions, but for $`T0`$ we expect that they are not important and that the partition function becomes independent of $`\mu `$. The lesson is that in order to obtain a condensate that is $`\mu `$-independent, one must treat the time derivative exactly, or, in other words, all Matsubara frequencies have to be taken into account. This was first worked out in detail for lattice QCD models at nonzero chemical potential . Let us discuss the strong coupling limit of a one-dimensional SU($`N_c`$) lattice QCD model with Kogut-Susskind fermions . In that case, the matrix elements of the discretized Dirac operator are given by
$$𝒟_{kl}=\delta _{k,l1}U_{kl}e^\mu \delta _{k,l+1}U_{kl}^{}e^\mu ,$$
(123)
where the indices are modulo $`N`$, the number of lattice sites. Antiperiodic boundary conditions result in an extra minus sign for the matrix elements $`𝒟_{1N}`$ and $`𝒟_{N1}`$. The matrices $`U_{kl}`$ are independent SU($`N_c`$) matrices. By using gauge invariance, one can reduce this partition function to a single SU($`N_c`$) integral, which can be performed analytically, resulting in the partition function
$$Z_N=2\mathrm{cosh}(N_cN\mu )+\frac{\mathrm{sinh}[(N_c+1)N\mathrm{sinh}^1m]}{\mathrm{sinh}[N\mathrm{sinh}^1m]}.$$
(124)
The number of lattice sites $`N`$ can be interpreted as the total number of Matsubara frequencies. A sharp transition is obtained in the limit $`N\mathrm{}`$. For example, the chiral condensate is given by
$$\underset{N\mathrm{}}{lim}\frac{1}{N}_m\mathrm{log}Z_N=\{\begin{array}{cc}N_c/\sqrt{1+m^2}& \text{for }\mathrm{sinh}\mu <m,\hfill \\ 0& \text{for }\mathrm{sinh}\mu >m.\hfill \end{array}$$
(125)
Similar results have been derived for a large-$`d`$ strong coupling expansion of lattice QCD at finite density (with $`d`$ the Euclidean dimensionality). The failure of the quenched approximation has been reproduced in such lattice models. The effect of including all Matsubara frequencies has also been investigated in a random matrix model obtained by replacing the remainder $`R`$ in Eq. (121) by a chiral random matrix. After ensemble averaging, one obtains a similar partition function with similar conclusions .
#### 6.2.4 QUENCHING AT NONZERO CHEMICAL POTENTIAL
The failure of the quenched approximation at nonzero chemical potential was first understood analytically in terms of chRMT by Stephanov . He showed that the quenched limit is the limit $`N_f0`$ of a partition function with the absolute value of the fermion determinant,
$$|det(𝒟(\mu )+m)|^{N_f},$$
(126)
rather than
$$[det(𝒟(\mu )+m)]^{N_f}.$$
(127)
The absolute value of the fermion determinant can be written as
$$det(𝒟+m)det(𝒟^{}+m^{})=det\left(\begin{array}{cc}iW+\mu & m\\ m& iW^{}+\mu \end{array}\right)det\left(\begin{array}{cc}iW^{}\mu & m^{}\\ m^{}& iW\mu \end{array}\right).$$
(128)
Writing the fermion determinant as a Grassmann integral, we observe that the quenched partition function can be interpreted as a partition function of quarks and conjugate antiquarks. Therefore, in addition to the usual Goldstone modes, we have Goldstone modes consisting of a quark and a conjugate antiquark . Such modes, with the same mass as the usual Goldstone bosons, have a nonzero baryon number. The critical chemical potential given by the mass of the lightest particle with nonzero baryon number is thus $`m_\pi /2`$. The failure of the quenched approximation thus has an important benefit: It allows us to write down the exact low-energy effective partition function, which is discussed in the next section. The product of the two determinants in Eq. (128) can be written as the determinant of a single Hermitian matrix. This procedure is known as Hermitization .
#### 6.2.5 QUENCHED DIRAC SPECTRA
A chRMT for quenched QCD at nonzero chemical potential is obtained by replacing the determinant in Eq. (117) by its absolute value. This model can be solved analytically in the large-$`N`$ limit . However, we will follow a different approach by analyzing the corresponding effective theory . The advantage of this approach is that the effective theory is as valid for quenched QCD as it is for the quenched chiral random matrix model.
For non-Hermitian matrices, the eigenvalues are scattered in the complex plane. Using the fact that
$$_z^{}\frac{1}{z}=\pi \delta ^2(z),$$
(129)
where the complex delta function is defined as $`\delta ^2(z)=\delta (\mathrm{Re}z)\delta (\mathrm{Im}z)`$, we find that the two-dimensional spectral density is (up to a normalization constant) given by
$$\rho (\lambda )=\frac{1}{\pi }_z^{}G(z)|_{z=\lambda }.$$
(130)
The resolvent is defined as
$$G(z)=\frac{1}{N_fV}_z\mathrm{log}Z,$$
(131)
and the quenched result is obtained in the limit $`N_f=0`$ of the partition function
$$Z=\stackrel{N_f}{det}(𝒟(\mu )+z)\stackrel{N_f}{det}(𝒟^{}(\mu )+z^{}).$$
(132)
The product in Eq. (132) can be written as
$$\stackrel{N_f}{det}(𝒟(\mu )+z)\stackrel{N_f}{det}(𝒟^{}(\mu )+z^{})=det\left(\begin{array}{cc}iW\mathrm{𝟏}_{2N_f}+B_R& \zeta \\ \zeta & iW^{}\mathrm{𝟏}_{2N_f}+B_L\end{array}\right)$$
(133)
with
$$BB_R=B_L=\left(\begin{array}{cc}\mu \mathrm{𝟏}_{N_f}& 0\\ 0& \mu \mathrm{𝟏}_{N_f}\end{array}\right)\mathrm{and}\zeta =\left(\begin{array}{cc}z\mathrm{𝟏}_{N_f}& 0\\ 0& z^{}\mathrm{𝟏}_{N_f}\end{array}\right),$$
(134)
where we have displayed the degeneracy in flavor space by means of the identity matrices $`\mathrm{𝟏}_{2N_f}`$ and $`\mathrm{𝟏}_{N_f}`$. For $`z=\mu =0`$, the quenched partition function is thus invariant under $`\mathrm{SU}_R(2N_f)\times \mathrm{SU}_L(2N_f)`$. This symmetry is broken spontaneously to the diagonal subgroup, $`\mathrm{SU}_V(2N_f)`$. At low energies, the effective partition function is therefore given by a partition function of Goldstone modes parameterized by matrices $`U\mathrm{SU}(2N_f)`$. The static effective Lagrangian is obtained from the requirement that it should have the same symmetries as the underlying microscopic partition function. The mass term was discussed in Sec. 4.1. The chemical potential term remains invariant under $`\mathrm{SU}(2N_f)\times \mathrm{SU}(2N_f)`$ transformations if at the same time $`B_R`$ and $`B_L`$, now considered as independent matrices, are transformed as
$$B_RU_RB_RU_R^1,B_LU_LB_LU_L^1.$$
(135)
The matrices $`U`$ in the Goldstone manifold transform as $`UU_RUU_L^1`$. To lowest nontrivial order in $`\mu `$, we can write two invariant combinations,
$$\mathrm{Tr}UB_LU^1B_R\mathrm{and}\mathrm{Tr}BB.$$
(136)
The coefficients follow from the conditions that the critical chemical potential is equal to one third of the mass of the lightest baryon and that the baryon density should vanish below $`\mu _c`$. For the static limit of the effective partition function, we find
$$Z=_{U\mathrm{SU}(2N_f)}DUe^{\frac{F^2V}{4}\mathrm{Tr}[U,B][U^1,B]+\frac{1}{2}\mathrm{\Sigma }V\mathrm{Tr}(MU+M^{}U^1)},$$
(137)
with the mass matrix given by
$$M=\left(\begin{array}{cc}z\mathrm{𝟏}_{N_f}& 0\\ 0& z^{}\mathrm{𝟏}_{N_f}\end{array}\right).$$
(138)
This partition function can be derived more elegantly by means of a local gauge invariance principle . For real masses $`z`$ the Dirac operator satisfies the relation $`\mathrm{det}(𝒟^{}(\mu )+z^{})=\mathrm{det}(𝒟(\mu )+z)`$, so that the partition function defined in Eq. (132) and the corresponding low-energy effective partition function given in Eq. (137) can be interpreted as the partition function of QCD at finite isospin density .
Below the critical chemical potential, $`\mu <m_\pi /2`$, only the vacuum state contributes to the partition function, so that
$$Z=e^{\mathrm{\Sigma }VN_f(z+z^{})}.$$
(139)
In terms of the effective theory, the saddle point is at $`U=\mathrm{𝟏}`$. The resolvent (131) is given by
$$G(z)=\frac{1}{N_fV}_z\mathrm{log}Z=\mathrm{\Sigma }.$$
(140)
Since $`G(z)`$ does not depend on $`z^{}`$, it follows from Eq. (130) that the spectral density is zero in the region where (140) holds. Because $`m_\pi ^2=(z+z^{})\mathrm{\Sigma }/F^2`$, see Eqs. (84) through (86), the condition $`\mu <m_\pi /2`$ means that the spectral density vanishes everywhere except in a strip
$$|\mathrm{Re}z|<\frac{2\mu ^2F^2}{\mathrm{\Sigma }}.$$
(141)
For $`z`$ inside this strip, the Goldstone modes contribute to the partition function. In terms of the effective partition function, the saddle point rotates away from $`U=\mathrm{𝟏}`$, leading to a nonvanishing diquark condensate. The rotation of the saddle point is a generic feature of low-energy effective partition functions of non-Hermitian field theories .
These results are in qualitative agreement with spectra of the quenched lattice QCD Dirac operator at $`\mu 0`$ . The correlations of the complex eigenvalues are not so well understood, but the first lattice QCD calculations seem to confirm the theoretical expectations .
### 6.3 Phase Diagram for the QCD Partition Function
The random matrix models at nonzero temperature and nonzero chemical potential introduced in the previous sections can be merged into a single schematic chRMT model for the chiral phase transition ,
$$Z=DW\stackrel{N_f}{det}\left(\begin{array}{cc}m& iW+iC\\ iW^{}+iC& im\end{array}\right)e^{\frac{1}{2}N\mathrm{\Sigma }^2\mathrm{Tr}WW^{}},$$
(142)
where $`C`$ is a diagonal matrix with $`C_{kk}=ti\mu `$ for one half of the diagonal elements and $`C_{kk}=ti\mu `$ for the other half. In the following, we set $`\mathrm{\Sigma }=1`$ for simplicity.
This random matrix partition function can be rewritten as a $`\sigma `$-model. For the case $`N_f=1`$ it is given by
$$Z(t,\mu )=𝑑\sigma e^{\frac{N}{2}\mathrm{Tr}\mathrm{\Omega }(\sigma )},$$
(143)
where
$`\mathrm{\Omega }(\sigma )`$ $`=`$ $`\sigma \sigma ^{}{\displaystyle \frac{1}{2}}\mathrm{log}[(\sigma +m)(\sigma ^{}+m)(\mu +it)^2]`$ (144)
$`{\displaystyle \frac{1}{2}}\mathrm{log}[(\sigma +m)(\sigma ^{}+m)(\mu it)^2].`$
For $`t=0`$ this is exactly the $`\sigma `$-model discussed in Sec. 6.2.2. At the saddle point, the chiral condensate is given by the expectation value of $`\sigma `$. For $`m=0`$ the saddle-point equation is of fifth order in $`\sigma `$,
$$\sigma [\sigma ^42(\mu ^2t^2+\frac{1}{2})\sigma ^2+(\mu ^2+t^2)^2+\mu ^2t^2]=0.$$
(145)
The critical points occur where one of the solutions of the quartic equation merges with the solution $`\sigma =0`$, i.e. along the curve $`(\mu ^2+t^2)^2+\mu ^2t^2=0`$. At the tricritical point, three solutions merge. This happens if in addition $`\mu ^2t^2+\frac{1}{2}=0`$.
Figure 10 shows the phase diagram in the $`\mu tm`$-space. In the $`m=0`$ plane, we observe a line of second-order phase transitions and a line of first-order phase transitions. They join at the tricritical point. Also joining at the tricritical point is a line of second-order phase transitions in the $`m`$-direction, which is the boundary of the plane of first-order transitions in $`\mu tm`$-space. This line is the collection of end points of lines of first-order phase transitions. We expect that the critical exponents for such a liquid-gas transition are given by the three-dimensional Ising model. The tricritical point was also found in a Nambu–Jona-Lasinio model .
Because the saddle-point equation (145) is of fifth order in $`\sigma `$, the critical properties of the random matrix model (142) are very similar to those of a $`\varphi ^6`$-theory. The critical dimension at the tricritical point of such theories is three, so that mean field theory, and therefore RMT, describes the correct critical behavior at this point. This random matrix model can be considered as the matrix equivalent of a Landau-Ginzburg functional. The advantage over using the standard Landau-Ginzburg theory is that in this case the spectrum of the Dirac operator is also accessible. This allows us to study the critical properties of the Dirac eigenvalues. For example, for $`\mu =0`$ we have found that the distribution of the smallest nonzero eigenvalue of the Dirac operator may serve as an order parameter for the chiral transition .
### 6.4 Random Matrix Triality at $`\mu 0`$
In previous sections, we have shown that the pattern of chiral symmetry breaking and the correlations of the Dirac eigenvalues are related to the anti-unitary symmetries of the Dirac operator. Since the chemical potential occurs only in the combination $`_0+\mu `$, the anti-unitary symmetries at $`\mu 0`$ are the same as for zero chemical potential. Thus for QCD with two colors and fermions in the fundamental representation the Dirac operator is real ($`\beta =1`$), whereas for QCD with adjoint fermions the Dirac operator is quaternion real ($`\beta =4`$). In the first case, the Dirac operator has the structure
$$𝒟=\left(\begin{array}{cc}0& W+\mu \\ W^T+\mu & 0\end{array}\right)\text{with }W\text{ real}.$$
(146)
In the second case, $`𝒟`$ is given by
$$𝒟=\left(\begin{array}{cc}0& W+\mu \\ W^{}+\mu & 0\end{array}\right)\text{with }W\text{ quaternion real}.$$
(147)
In this case, the quaternion real matrix elements of $`W`$ satisfy the reality relation $`W_{kl}^{}=\sigma _2W_{kl}\sigma _2`$. In both cases, it is easy to show that the fermion determinant is real. Furthermore, because of
$`\stackrel{N_f}{det}W\stackrel{N_f}{det}W^T=\stackrel{2N_f}{det}W`$ $`\text{for }\beta =1,`$ (148)
$`\stackrel{N_f}{det}W\stackrel{N_f}{det}W^{}=\stackrel{2N_f}{det}W`$ $`\text{for }\beta =4,`$ (149)
the flavor symmetry group is enhanced to U($`2N_f`$) in both cases (for $`\nu =0`$).<sup>3</sup><sup>3</sup>3For $`\beta =4`$ and $`\mu =0`$, the Dirac operator satisfies the relation $`(C𝒟)^T=C𝒟`$ whereas $`(C\mu \gamma _0)^T=C\mu \gamma _0`$ (with $`C`$ the charge conjugation matrix). Therefore, for $`\beta =4`$ and $`\mu 0`$, the Pfaffian of the Dirac operator is not defined and the fermion determinant $`det^{N_f/2}(𝒟+\mu \gamma _0)`$ cannot be expressed as an integral over $`N_f`$ Majorana fermions with flavor symmetry U($`N_f`$), as discussed in Sec. 2.2.3. Instead, we write $`det^{N_f}(𝒟+\mu \gamma _0)`$ as an integral over $`N_f`$ Dirac fermions with flavor symmetry U($`2N_f`$). The full flavor symmetry is broken spontaneously to Sp($`2N_f`$) and O($`2N_f`$), respectively, and is broken in the same way by the mass term. The chemical potential breaks the symmetry according to $`\mathrm{U}(2N_f)\mathrm{U}(N_f)\times \mathrm{U}(N_f)`$ in both cases. For both $`\beta =1`$ and $`\beta =4`$ the pseudoreality of the Dirac operator leads to Goldstone modes with a nonzero baryon number in the same way we have seen for the phase quenched partition function. An exact low-energy effective partition function valid to lowest order in $`m`$ and $`\mu ^2`$ can be written down . For an elaborate discussion of the symmetries, their breaking pattern, and the low-energy effective theory in both cases, see Ref. .
For an even number of flavors, dynamical Monte Carlo simulations are possible for $`\beta =1`$ and $`\beta =4`$, though not for $`\beta =2`$. However, quenched simulations have been performed for all three classes. A cut along the imaginary axis below a cloud of eigenvalues was found in instanton liquid simulations for $`N_c=2`$ at $`\mu 0`$, which corresponds to $`\beta =1`$. In lattice QCD simulations with staggered fermions for $`N_c=2`$ , a depletion of eigenvalues along the imaginary axis was observed, whereas for $`N_c=3`$ the eigenvalue distribution did not show any pronounced features .
A chiral random matrix model for the Dirac operator at $`\mu 0`$ in each symmetry class is obtained by drawing the matrix elements of $`W`$ from a Gaussian probability distribution. In the quenched approximation, the spectral properties of the random matrix Dirac operator of Eq. (118) can easily be studied numerically by diagonalizing a set of matrices with the probability distribution of Eq. (117). Figure 11 shows results for the eigenvalues of a few $`100\times 100`$ matrices for $`\mu =0.15`$ (dots). The solid curves represent the analytical result for the boundary of the domain of eigenvalues derived in Ref. for $`\beta =2`$. However, the analysis can be extended to $`\beta =1`$ and $`\beta =4`$, and with the proper scale factors, the solution is identical.
For $`\beta =1`$ and $`\beta =4`$ we observe exactly the same structure as in the previously mentioned (quenched) QCD simulations. There is an accumulation of eigenvalues on the imaginary axis for $`\beta =1`$ and a depletion of eigenvalues along this axis for $`\beta =4`$. The depletion can be understood as follows. For $`\mu =0`$ all eigenvalues are doubly degenerate. This degeneracy is broken at $`\mu 0`$, which produces the observed repulsion between the eigenvalues.
The number of purely imaginary eigenvalues for $`\beta =1`$ scales as $`\sqrt{N}`$ and is thus not visible in a leading-order saddle-point analysis. Such a $`\sqrt{N}`$ scaling is typical for the regime of weak non-Hermiticity first identified by Fyodorov et al. . Using the supersymmetric method of RMT, Efetov obtained the $`\sqrt{N}`$ dependence analytically. The case $`\beta =4`$ was also analyzed analytically , with results that are in complete agreement with the numerical simulations. Obviously, more work is needed in order to arrive at a complete characterization of universal features of the spectrum of non-Hermitian operators.
## 7 RELATED MODELS AND RELATIONS TO OTHER FIELDS
Random matrix theory has been used extensively in many different fields including nuclear physics, atomic and molecular physics, condensed matter physics, quantum chaos, quantum gravity, and mathematical physics. Before reviewing a few applications that might have some relation to QCD, we briefly discuss another class of matrix models for the QCD partition function. For a more comprehensive review of the material in this section, see Refs. .
### 7.1 One-Plaquette Models of Lattice QCD
One class of matrix models for QCD are lattice QCD partition functions on a $`2^d`$ lattice. The simplest model in this class is the one-plaquette partition function for pure gauge theory in $`d=2`$. This model is known as the Brézin-Gross-Witten model and is defined by
$$Z(J,J^{})=_{U\mathrm{U}(N)}DUe^{\mathrm{Tr}(JU^{}+J^{}U)},$$
(150)
where the integral is over the Haar measure of U($`N`$) and $`J`$ is an arbitrary complex source term. We have shown that this partition function can be expressed as a determinant of modified Bessel functions . For $`J`$ a multiple of the identity (with the interpretation of $`J1/g^2`$, where $`g`$ is the Yang-Mills coupling constant), this model was solved analytically in the large-$`N`$ limit by Gross and Witten. They found a third-order phase transition at a critical value of $`J`$. A large-$`N`$ chiral phase transition is obtained by extending this model with a fermion determinant . Despite considerable effort, the generalization of this model to more than two dimensions, which is known as the Eguchi-Kawai model , could not be solved analytically .
The Brézin-Gross-Witten model can be rewritten as a generalized Kontsevich model, which has received a great deal of attention in the theory of exactly integrable systems . There are many other relations between RMT and the theory of exactly solvable systems. Among others, the asymptotic properties of correlation functions can be derived by means of conformal field theory . An interesting overview in the context of string theory is given in Ref. .
Another class of related models is the Kazakov-Migdal model , which is also known as induced QCD. In its simplest form it is defined by
$$Z=DHDUe^{_x\mathrm{Tr}H^2(x)+_{x,\mu }\mathrm{Tr}[U(x)H(x)U^{}(x)H(x+\mu )]},$$
(151)
where the $`H(x)`$ are Hermitian matrices and the $`U(x)`$ are unitary matrices. The sum is over a lattice in $`d`$ dimensions. A zero-dimensional form of this model was proposed as a model for the transition between Poisson statistics and Wigner-Dyson statistics .
### 7.2 Universal Conductance Fluctuations in Disordered Mesoscopic Systems
In condensed matter physics, a mesoscopic system is a system whose linear size is larger than the elastic mean free path of the electrons but smaller than the phase coherence length, which is essentially the inelastic mean free path. A typical size is about 1 $`\mu `$m. The conductance $`g`$ of mesoscopic samples is closely related to their spectral properties. Using a scaling block picture, Thouless found that in the diffusive regime, the conductance is given by $`g=E_C/\mathrm{\Delta }`$, where $`E_C/\mathrm{}`$ is the inverse diffusion time of an electron through the sample and $`\mathrm{\Delta }`$ is the mean level spacing . This can be rewritten as $`g=N(E_C)`$, where $`N(E)`$ is the mean level number in an energy interval $`E`$. Thus the variance, $`\delta g^2`$, of the conductance is related to the number variance, $`\mathrm{\Sigma }^2`$, of the energy levels .
Low-temperature experiments have been performed in which the conductance of mesoscopic wires was measured as a function of an external magnetic field. The observed fluctuations in $`g`$ are of the order of $`e^2/h`$, independent of the details of the system (shape, material, etc). These are the so-called universal conductance fluctuations . One can understand this phenomenon qualitatively by estimating the number fluctuations of the electron levels using RMT results. Both the magnitude of the fluctuations and their universality can be obtained through the transfer matrix method. An interesting numerical result is that the density of eigenvalues of the transmission matrix in the Hofstadter model for universal conductance fluctuations can be described in terms of the microscopic spectral density of the chUE .
### 7.3 Anderson Localization
In more than two dimensions, a good conductor becomes an insulator when the disorder becomes sufficiently strong. This phenomenon is called Anderson localization. In the localized phase, the wave function of the electron is not described by Bloch waves, but by a localized form that decays exponentially,
$$\psi (r)e^{r/L_c}.$$
(152)
The length scale $`L_c`$ is known as the localization length. This phenomenon was first described by the Anderson model , which is a hopping model with a random potential on each lattice point. The dimensionality of the lattice plays an important role. It has been shown that in one dimension all states are localized. The critical dimension is two, whereas in three dimensions there is a delocalization transition at an energy $`E_L`$. The states below $`E_L`$ are localized whereas the states above $`E_L`$ are extended, i.e. with a wave function that scales with the size of the system. The eigenvalues of the localized states are not correlated, and their correlations are described by the Poisson distribution.
An interesting question is whether Dirac eigenfunctions can be exponentially localized. Parisi argued that localized states can only occur in quenched systems . In QCD this argument goes as follows. If the eigenfunctions were spatially localized, the eigenvalues would be uncorrelated, and there would be no repulsion between the eigenvalues. Because of the fermion determinant and the measure in Eq. (36), the eigenvalues would be repelled from the origin. Since there would be no mechanism to compensate for this repulsion, $`\rho (0)/V`$ and the chiral condensate would be zero. Therefore, if chiral symmetry is spontaneously broken, the eigenfunctions of the Dirac operator must be spatially extended. Indeed, this has been found for the wave functions of the Dirac operator for a gauge field ensemble given by a liquid of instantons as well as in lattice QCD .
In the extended domain, the situation is more complicated. An important energy scale is the Thouless energy , which is related to the diffusion time of an electron through the sample (see Sec. 7.2). With that diffusion time given by $`L^2/D`$ (the diffusion constant is denoted by $`D`$), the Thouless energy is
$$E_c=\frac{\mathrm{}D}{L^2}.$$
(153)
On time scales larger than $`\mathrm{}/E_c`$, an initially localized wave packet diffuses all over phase space. If this wave function $`\psi (t)`$ at $`t>\mathrm{}/E_c`$ is expressed as a superposition of eigenfunctions $`\varphi _i(t=0)`$, many of the overlaps $`\psi (t)|\varphi _i`$ are nonzero. Therefore we expect that for energy differences below $`E_c`$ the eigenvalues are correlated according to RMT.
The diffusive behavior of electrons is described by Goldstone modes, called diffusons. Classically, they satisfy a diffusion equation . The diffusion modes can be described in terms of a $`\sigma `$-model, which has the same structure as the low-energy effective theory for the QCD partition function. The Thouless energy corresponds to the energy scale below which the fluctuations of the zero-momentum modes dominate the fluctuations of the nonzero-momentum modes, i.e. the scale discussed in Sec. 4.4, below which the QCD Dirac spectra are given by chRMT. The analysis of classically chaotic quantum systems has made it clear that spectra are described by RMT if and only if the corresponding classical motion is chaotic . If we interpret the Dirac operator as a Hamiltonian in four Euclidean dimensions plus one artificial time dimension, we thus conclude that the classical motion of the quarks is chaotic<sup>4</sup><sup>4</sup>4There is an elaborate literature on the classical dynamics of gauge fields. For a discussion of this topic, see Ref. . and therefore diffusive in four Euclidean dimensions and one artificial time dimension, in agreement with the interpretation of Goldstone modes as diffusion modes . The equivalent of the diffusion constant can be read off immediately by comparing the two $`\sigma `$-models, or equivalently by comparing the scale of Eq. (88) with the Thouless energy, and is given by $`DF^2/\mathrm{\Sigma }`$ . Quantum mechanically, the Thouless energy can be interpreted as the spreading width of the exact eigenstates over the states of the noninteracting Hamiltonian .
### 7.4 Non-Hermitian Models with Disorder
In the past few years, several non-Hermitian models with disorder have been introduced in the literature. We mention only the simplest model, which was motivated by the study of flux-line pinning in superconductors . Its Hamiltonian is given by
$$H=\frac{1}{2m}(𝐩+i𝐡)^2+V(𝐫),$$
(154)
where $`V(𝐫)`$ is a random disorder potential, p is the momentum operator, and $`i𝐡`$ is a constant imaginary vector potential. The new feature of such models is that a localization-delocalization transition can occur even in one and two dimensions. Although the wave functions corresponding to real eigenvalues remain localized, the wave functions of complex eigenvalues can be extended. For a detailed discussion of such models, see Refs. .
The dependence on $`𝐡`$ in the Hamiltonian $`H`$ occurs only in the combination $`𝐩+i𝐡`$, i.e. in exactly the same way as the chemical potential in the QCD Dirac operator. There are two important differences from QCD: $`(a)`$ the operator $`H`$ does not have a chiral structure, and $`(b)`$ the disorder is uncorrelated Gaussian and quenched. The connection between delocalization in the Hatano-Nelson model (154) and diquark condensation in quenched QCD at nonzero chemical potential is not yet understood and deserves further attention.
After the initial work of Ginibre , non-Hermitian RMT has received a great deal of attention in the mathematics literature as well. Probably the best overview of results in this area is in the book by Girko .
### 7.5 Andreev Scattering
The term Andreev scattering (or Andreev reflection) refers to a process in which an electron hits the interface between a normal metal and a superconductor and is reflected as a hole with the opposite momentum (or vice versa) . Stated differently, two electrons can tunnel through the interface so that a Cooper pair is added to the superconducting condensate (or removed from it). This process can be described in a microscopic mean field model by the Bogoliubov-deGennes Hamiltonian, which can be written in matrix form as
$$H=\left(\begin{array}{cc}A& B\\ B^{}& A^T\end{array}\right),$$
(155)
where $`A`$ ($`A^T`$) represents the Hamiltonian for particles (holes) and $`B`$ represents the pairing field. The requirement that $`H`$ be Hermitian means that $`A`$ must be Hermitian. If the system under consideration is invariant under time reversal or spin rotation or both, there may be additional symmetries. If the system is invariant under spin rotation, $`B`$ is symmetric (in this case we consider the Hamiltonian in the space of spin-up states only). If, in addition, the system is invariant under time reversal, the elements of $`H`$ are real. Otherwise, they are complex. On the other hand, if the system is not invariant under spin rotation, $`B`$ is antisymmetric. In this case, if the system is invariant under time reversal, the elements of $`H`$ are quaternion real. Otherwise, they are complex. This classification establishes four new random matrix ensembles . The microscopic spectral correlations of these ensembles are identical to those of the chiral ensembles if the parameters (the Dyson index $`\beta `$ and the number of massless flavors $`N_f`$) are chosen appropriately.
### 7.6 Mathematical Physics and Quantum Gravity
RMT has received a great deal of attention as a problem in mathematical physics. We mention the relation between universal behavior in RMT and the asymptotic behavior of orthogonal polynomials , the relation between the classification of random matrix ensembles and the Cartan classification of symmetric spaces , and the theory of Riemannian superspaces . Various methods for the solution of random matrix problems have been proposed. We mention the mapping to a gas of noninteracting fermions, the Coulomb gas method , the Brownian motion method , the replica method , the orthogonal polynomial method , the supersymmetric method , and the operator method .
Generally, the problem for $`\beta =2`$ is much simpler than for $`\beta =1`$ and $`\beta =4`$. Nevertheless, a great deal of progress has been made for $`\beta =1`$ and $`\beta =4`$. We mention relations between the kernels of the correlation functions and the relation between massless and massive correlators . Novel mathematical methods have been developed for $`\beta =1`$ and $`\beta =4`$. We mention only the skew-orthogonal polynomial method and the extension of an operator method for $`\beta =2`$ .
There exists an elaborate literature on the random matrix formulation of 2-$`d`$ quantum gravity (see e.g. the recent reviews by Abdalla et al. and Di Franceso et al. ). The idea is that the partition function, which is a sum over all metrics, can be rewritten as a sum over random surfaces. The sum over the dual graphs of a discretization of these random surfaces can be rewritten in terms of a random matrix partition function that can be analyzed with standard random matrix methods.
## 8 CONCLUSIONS
Chiral random matrix theory is the archetypal model for the spontaneous breaking of chiral symmetry. In chRMT, chiral symmetry is spontaneously broken for any finite number of massless flavors, whereas in QCD, breaking is believed to occur only below a certain number of massless flavors, perhaps as few as four. In this review, we have studied chiral symmetry breaking from the perspective of the eigenvalues of the QCD Dirac operator, where broken chiral symmetry implies that the smallest eigenvalues are spaced as $`1/V`$. The robustness of chiral symmetry breaking in chRMT can then be understood as the crystallization of the eigenvalues by the long-range random interactions. The number of eigenvalues in a sequence containing $`N`$ eigenvalues fluctuates on average only by $`\sqrt{\mathrm{log}N}/\pi `$ as opposed to $`\sqrt{N}`$ for uncorrelated eigenvalues. For uncorrelated eigenvalues, we do not expect chiral symmetry breaking for any number of massless flavors. The conclusion is that chiral symmetry breaking in QCD requires strongly correlated Dirac eigenvalues but not quite as strongly correlated as in chRMT. To make this conclusion more quantitative, we have formulated an effective theory for the QCD Dirac spectrum, which, in addition to the usual Goldstone modes, contains Goldstone bosons of spectral quarks that are “ghost” quarks introduced to probe the Dirac spectrum. Like the usual chiral Lagrangian, this Lagrangian consists of a kinetic term and a mass term. An important scale is the spectral mass for which the long-wavelength fluctuations of these two terms are of equal order of magnitude. In the theory of disordered condensed matter systems, this scale is known as the Thouless energy. For spectral masses below the Thouless energy, the mass dependence is given by the zero-momentum sector of the effective theory, and it coincides with that of chRMT in the limit of large matrices. The deviations from chRMT, given by the contributions of the nonzero-momentum modes, increase the fluctuations of the eigenvalues. For an increasing number of flavors, the position of the smallest eigenvalues moves to the right, whereas the slope of the spectral density increases with $`N_f`$. However, chiral symmetry remains broken for any value of $`N_f`$. The chiral condensate is simply a parameter of the effective partition function.
The statistical properties of the QCD Dirac eigenvalues have been investigated by numerous lattice QCD simulations, and the chRMT predictions have been verified in great detail. The behavior of Dirac spectra in the domain beyond the Thouless energy has been less well studied, but in all cases, the predictions of the effective theory for the Dirac spectrum agree well with lattice simulations.
Inspired by the successes of chRMT as an exact description of the statistical properties of the QCD Dirac spectrum, we have constructed a schematic model for the chiral phase transition. Although this approach does not provide rigorous results, it has been useful in advancing a qualitative understanding of the chiral phase transition at both zero and nonzero chemical potential. Examples are the properties of the Dirac spectrum in the approach to the critical temperature, the failure of the quenched approximation at $`\mu >0`$, and the phase diagram of the QCD partition function in the $`\mu Tm`$ plane.
Finally, RMT encapsulates a duality between order and chaos. Complete mixing of all degrees of freedom can be described by a single effective degree of freedom. This is the progress we have made toward an understanding of the complexity of the QCD vacuum.
## Acknowledgments
This work was supported in part by the US Department of Energy under contracts DE-FG-88ER40388, DE-FG02-91ER40608, and DE-AC02-98CH10886. We would like to acknowledge discussions and collaborations on the subject of this review with G. Akemann, I.L. Aleiner, Y. Alhassid, A. Altland, B.L. Altshuler, M.E. Berbenni-Bitsch, B.A. Berg, D. Dalmazi, P.H. Damgaard, Y.V. Fyodorov, A.M. Garcia-Garcia, M. Göckeler, T. Guhr, M.Á. Halász, H. Hehl, A.D. Jackson, N. Kaiser, B. Klein, J.-Z. Ma, H. Markum, S. Meyer, S.M. Nishigaki, J.C. Osborn, R. Pullirsch, P.E.L. Rakow, A. Schäfer, M. Schnabel, A. Schwenk, B. Seif, T.H. Seligman, M.K. Şener, B.D. Simons, R.E. Shrock, E.V. Shuryak, A.V. Smilga, M.A. Stephanov, D. Toublan, H.A. Weidenmüller, T. Wilke, and M. Zirnbauer. We also thank I.L. Aleiner for a critical reading of Sec. 7.
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# 1 Introduction
## 1 Introduction
The sine-Gordon S-matrix for the quantum scattering of solitons was solved initially in the literature in , by imposing conditions on the S-matrix called crossing and unitarity. These conditions were previously observed in other S-matrices found using perturbative techniques in quantum field theory , and were subsequently taken to be axioms for the non-perturbative result.
The bootstrap methods use the $`U_q(\widehat{sl_2})`$ Hopf algebra , and its spin half representation $`W`$, which is a space of functions taking values in $`C^2`$, where the first component corresponds to ‘solitons’ and the second to ‘anti-solitons’. Then $`WW`$ corresponds to a two soliton system, which can interact by collision. There is an initial two soliton quantum state in $`WW`$, and after the collision process we have a final quantum state in $`WW`$. The scattering matrix gives a map $`WWWW`$, which sends the initial to the final state. In terms of the Hopf algebra, this map is a braiding .
This scattering matrix is fairly easy to find up to a multiplication by a scalar function, but the scalar function itself is more difficult. We denote this scalar function by $`a(z)`$ where $`zC`$. The crossing condition in terms of this function becomes
$$a(z)=a\left(\frac{q}{z}\right)\frac{(zz^1)}{(zq^1z^1q)}$$
(1)
and the unitarity condition is $`a(z)a(z^1)=1`$.
The crossing condition arises from physics by equating a scattering process with the same scattering process after rotating the space-time diagram of the collision by a right-angle. This rotation involves a time reversal of one of the incoming and one of the outgoing solitons, which implies that these are turned into anti-solitons. The rotation of the diagram also implies that $`z`$, which corresponds roughly to the relative momentum of the colliding solitons (conserved in the collision), is transformed to $`z^1q`$. The interested reader can refer to the comprehensive book for a complete discussion. The article is also an excellent review.
Zamolodchikov-Zamolodchikov solved these equations (1) to get a formula for $`a(z)`$ in terms of a double infinite product of gamma functions. It was thought that this product probably converged for all physical values of $`z`$, and all values of $`q`$ where $`|q|=1`$. However no proof of this was provided. The convergence of this function has also not been treated in subsequent papers in the literature.
Here, in this paper, we find an alternative formula which is more amenable to a convergence analysis. We find that the convergence is highly delicate, and the function converges with probability one, for $`q`$ on the unit circle, i.e. $`q=e^{2\pi i\tau }`$, including convergence for all irrational algebraic values of $`\tau `$, and diverges for certain transcendental values of $`\tau `$.
Another alternative formula for $`a(z)`$ was found in Johnson , as an integral, or as a combination of ‘regularised’ quantum dilogarithms. This was done to make contact with semi-classical results for the scattering which involve integrating classical time delays. However, the convergence of this formula, or of the individual quantum dilogarithms, was also difficult to analyse because the contour integrals which one has to do are difficult to perform, involving sums over an infinite double set of poles.
In this paper we shall consider the problem purely in terms of braidings of the representations of a Hopf algebra, rather than invoking the crossing and unitarity conditions of S-matrix theory. We begin with the loop group of analytic functions used in the classical inverse scattering procedure for sine-Gordon , and deform this by inclusion of a parameter $`q`$ .
## 2 The universal enveloping algebra $``$
We begin with the loop group of analytic functions from $`C^{}`$ to $`SL_2(C)`$ which obey the symmetry condition $`U\varphi (z)U^1=\varphi (z)`$, where
$$U=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right).$$
The Lie algebra for this group has generators $`X_{\pm 1}`$, $`X_{\pm 2}`$ and $`H`$, given by
$`X_{+1}(z)=\left(\begin{array}{cc}0& 0\\ z& 0\end{array}\right)`$ , $`X_1(z)=\left(\begin{array}{cc}0& 1/z\\ 0& 0\end{array}\right),H(z)=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),`$ (2)
$`X_{+2}(z)=\left(\begin{array}{cc}0& z\\ 0& 0\end{array}\right)`$ , $`X_2(z)=\left(\begin{array}{cc}0& 0\\ 1/z& 0\end{array}\right).`$ (3)
These generators obey the usual coproduct rule for the universal enveloping algebra of a Lie algebra, namely $`\mathrm{\Delta }(X_{\pm 1})=X_{\pm 1}1+1X_{\pm 1}`$ etc. This rule can be deformed by the inclusion of a parameter $`qC`$ to give
$`\mathrm{\Delta }H=1H+H1`$ , $`\mathrm{\Delta }X_{\pm 1}=X_{\pm 1}q^{H/2}+q^{H/2}X_{\pm 1},`$ (4)
$`\mathrm{\Delta }X_{\pm 2}`$ $`=`$ $`X_{\pm 2}q^{H/2}+q^{H/2}X_{\pm 2}.`$ (5)
The Lie algebra structure remains the same. If in addition we define a counit $`ϵ`$ (which kills all the generators) and an antipode $`S`$ (which has $`S(H)=H`$ and $`S(X_{\pm n})=q^{\pm 1}X_{\pm n}`$) we can make the universal enveloping algebra into a Hopf algebra, which we denote $``$. This is the so-called ‘principal gradation’ of $`U_q(\widehat{\mathrm{𝑠𝑙}_2})`$, which is actually a subalgebra of $`U_q(\widehat{\mathrm{𝑠𝑙}_2})`$. We use the convention that $`r=\sqrt{q}`$ in the terms $`q^{\pm H/2}`$.
## 3 Rigid braided tensor categories
Here we shall give a highly abbreveiated, specialised and incomplete account of rigid braided tensor categories, for a full account see .
The representations of a Hopf algebra $``$ form a category, with objects the representations, and the morphisms $`\rho :VW`$ are intertwining maps for the representations $`V`$ and $`W`$. This means that $`\rho `$ is linear and that $`\rho (h(v))=h(\rho (v))`$ for all $`h`$.
The tensor product of two representations is also a representation. The action on $`VW`$ is given by the coproduct $`\mathrm{\Delta }:`$. If we write $`\mathrm{\Delta }h=h_{(1)}h_{(2)}`$, then $`h(vw)=h_{(1)}(v)h_{(2)}(w)`$. The category contains an ‘identity object’, the representation $`C`$ with all generators having zero action. This means that category of representations forms a tensor or monoidal category. (Technically we should also say that the associator is trivial.)
If the category is rigid, for an object $`V`$ there is a dual object $`V^{}`$, and an ‘evaluation’ morphism $`\mathrm{eval}:V^{}VC`$ given by $`\mathrm{eval}(\alpha ,v)=\alpha (v)`$.
If the category is braided, for two objects $`V`$ and $`W`$ there is a morphism $`\mathrm{\Psi }_{VW}:VWWV`$. The braiding is functorial, which means that if there is a morphism $`\theta :VX`$, then the maps $`(I\theta )\mathrm{\Psi }_{VW}:VWWX`$ and $`\mathrm{\Psi }_{XW}(\theta I):VWWX`$ are the same.
Figure 1 shows the standard diagramatic notation used for braided categories. Elements of representations are denoted by vertical lines. For evaluation (a), note that $`C`$ is traditionally denoted by an invisible line. The braiding $`\mathrm{\Psi }_{VW}:VWWV`$ is shown in (b), and the rule for the functoriality of the braiding in (c).
The finite dimensional representations of a quasitriangular Hopf algebra form a rigid braided tensor category. We shall assume that the representations of our Hopf algebra $``$ also form a rigid braided tensor category, and this will allow us to explicitly calculate the braiding.
## 4 The ‘standard’ representation of $``$
Take $`W`$ to be a vector space of analytic functions $`:C^{}C^2`$ (or at least analytic in a neighbourhood of zero and a neighbourhood of infinity), which obeys the condition $`Uw(z)=w(z)`$ for all $`wW`$. The algebra $``$ acts on $`W`$ using matrix multiplication, $`(hw)(z)=h(z)w(z)`$, for the five generators listed in (3).
We can consider a dual space $`W^{}`$ to $`W`$, which shall consist of analytic functions $`:C^{}C^2`$ which are defined for $`|z|`$ sufficiently small and sufficiently large. Now define an evaluation map $`\mathrm{eval}:W^{}WC`$ by
$`\mathrm{eval}\left(\left(\begin{array}{c}f\\ g\end{array}\right)\left(\begin{array}{c}u\\ v\end{array}\right)\right)`$ $`=`$ $`{\displaystyle \frac{1}{4\pi i}}{\displaystyle _\gamma }\left(f(z)u(z)+g(z)v(z)\right){\displaystyle \frac{dz}{z}},`$ (6)
where $`\gamma `$ consists of two anticlockwise circular contours about $`0`$, one of large radius, and one of small radius. To find the action of $``$ on $`W^{}`$ we use the fact that the action commutes with $`\mathrm{eval}:W^{}WC`$, and that the action of the generators is zero on $`C`$. This means that
$`0`$ $`=`$ $`H\left(\mathrm{eval}\left(\left(\begin{array}{c}f\\ g\end{array}\right)\left(\begin{array}{c}u\\ v\end{array}\right)\right)\right)`$
$`=`$ $`\mathrm{eval}\left(H\left(\left(\begin{array}{c}f\\ g\end{array}\right)\left(\begin{array}{c}u\\ v\end{array}\right)\right)\right)=\mathrm{eval}\left(H\left(\begin{array}{c}f\\ g\end{array}\right)\left(\begin{array}{c}u\\ v\end{array}\right)+\left(\begin{array}{c}f\\ g\end{array}\right)H\left(\begin{array}{c}u\\ v\end{array}\right)\right),`$
so we get
$`\mathrm{eval}\left(H\left(\begin{array}{c}f\\ g\end{array}\right)\left(\begin{array}{c}u\\ v\end{array}\right)\right)`$ $`=`$ $`\mathrm{eval}\left(\left(\begin{array}{c}f\\ g\end{array}\right)\left(\begin{array}{c}u\\ v\end{array}\right)\right)`$
$`=`$ $`{\displaystyle \frac{1}{4\pi i}}{\displaystyle \left(f(z)u(z)g(z)v(z)\right)\frac{dz}{z}},`$
so we deduce that
$$H\left(\begin{array}{c}f\\ g\end{array}\right)(z)=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)\left(\begin{array}{c}f(z)\\ g(z)\end{array}\right).$$
Now we continue with the generator $`X_{+1}`$,
$`0`$ $`=`$ $`\mathrm{eval}\left(X_{+1}\left(\left(\begin{array}{c}f\\ g\end{array}\right)\left(\begin{array}{c}u\\ v\end{array}\right)\right)\right)`$
$`=`$ $`\mathrm{eval}\left(X_{+1}\left(\begin{array}{c}f\\ g\end{array}\right)\left(\begin{array}{c}u/r\\ vr\end{array}\right)+\left(\begin{array}{c}f/r\\ gr\end{array}\right)X_{+1}\left(\begin{array}{c}u\\ v\end{array}\right)\right),`$
from which we get
$$\mathrm{eval}\left(X_{+1}\left(\begin{array}{c}f\\ g\end{array}\right)\left(\begin{array}{c}u/r\\ vr\end{array}\right)\right)=\mathrm{eval}\left(\left(\begin{array}{c}f/r\\ gr\end{array}\right)\left(\begin{array}{c}0\\ zu\end{array}\right)\right)=\frac{r}{4\pi i}zg(z)u(z)\frac{dz}{z}.$$
From this, and the corresponding calculations for the other generators, we see that the action on $`W^{}`$ is given by the matrix multiplication
$`\left(h\left(\begin{array}{c}f\\ g\end{array}\right)\right)(z)`$ $`=`$ $`\stackrel{~}{h}(z)\left(\begin{array}{c}f(z)\\ g(z)\end{array}\right),`$ (7)
where the matrices $`\stackrel{~}{h}(z)`$ are given by
$`\stackrel{~}{X}_{+1}(z)=q\left(\begin{array}{cc}0& z\\ 0& 0\end{array}\right)`$ , $`\stackrel{~}{X}_1(z)=q^1\left(\begin{array}{cc}0& 0\\ 1/z& 0\end{array}\right),\stackrel{~}{H}(z)=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),`$ (8)
$`\stackrel{~}{X}_{+2}(z)=q\left(\begin{array}{cc}0& 0\\ z& 0\end{array}\right)`$ , $`\stackrel{~}{X}_2(z)=q^1\left(\begin{array}{cc}0& 1/z\\ 0& 0\end{array}\right).`$ (9)
There is a morphism $`\theta :WW^{}`$ given by
$`\theta (k)(z)`$ $`=`$ $`\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)k(zq).`$ (10)
To verify this we check that it commutes with the actions of the generators, for example;
$`(X_{+1}\theta (\left(\begin{array}{c}u\\ v\end{array}\right)))(z)`$ $`=`$ $`\left(\begin{array}{cc}0& qz\\ 0& 0\end{array}\right)\theta (\left(\begin{array}{c}u\\ v\end{array}\right))(z)=\left(\begin{array}{cc}0& qz\\ 0& 0\end{array}\right)\left(\begin{array}{c}v(zq)\\ u(zq)\end{array}\right),`$
$`\theta (X_{+1}\left(\begin{array}{c}u\\ v\end{array}\right))(z)`$ $`=`$ $`\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)(X_{+1}\left(\begin{array}{c}u\\ v\end{array}\right))(zq)=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)\left(\begin{array}{cc}0& 0\\ qz& 0\end{array}\right)\left(\begin{array}{c}u(zq)\\ v(zq)\end{array}\right),`$
which are equal as required.
## 5 The braidings
We use the convention for tensor products that $`C^2C^2C^4`$ and $`M_2M_2M_4`$, where
$`\left(\begin{array}{c}u\\ v\end{array}\right)\left(\begin{array}{c}u^{}\\ v^{}\end{array}\right)\left(\begin{array}{c}uu^{}\\ uv^{}\\ vu^{}\\ vv^{}\end{array}\right),\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)\left(\begin{array}{cc}a^{}& b^{}\\ c^{}& d^{}\end{array}\right)\left(\begin{array}{cccc}aa^{}& ab^{}& ba^{}& bb^{}\\ ac^{}& ad^{}& bc^{}& bd^{}\\ ca^{}& cb^{}& da^{}& db^{}\\ cc^{}& cd^{}& dc^{}& dd^{}\end{array}\right).`$ (11)
Then we can consider $`WW`$ or $`W^{}W`$ as a space of analytic maps from a subset of $`C^{}\times C^{}`$ to $`C^4`$. Now if $`kW^{}W`$ we have
$$\mathrm{eval}(k)=\frac{1}{4\pi i}\left(\begin{array}{cccc}1& 0& 0& 1\end{array}\right)k(x,x)\frac{dx}{x}.$$
By the coproduct rule $`H`$ and (for example) $`X_{+1}`$ act on $`WW`$ by matrix multiplication
$`(Hk)(x,y)`$ $`=`$ $`(H(x)I_2+I_2H(y))k(x,y),`$
$`(X_{+1}k)(x,y)`$ $`=`$ $`(X_{+1}(x)q^{H(y)/2}+q^{H(x)/2}X_{+1}(y))k(x,y),`$
and they act on $`W^{}W`$ by
$`(Hk)(x,y)`$ $`=`$ $`(\stackrel{~}{H}(x)I_2+I_2H(y))k(x,y),`$
$`(X_{+1}k)(x,y)`$ $`=`$ $`(\stackrel{~}{X}_{+1}(x)q^{H(y)/2}+q^{\stackrel{~}{H}(x)/2}X_{+1}(y))k(x,y).`$
The braiding $`\mathrm{\Psi }_{WW}:WWWW`$ will be assumed to have the form $`(\mathrm{\Psi }_{WW}k)(x,y)=M(x,y)k(y,x)`$, where $`M(x,y)`$ is a $`4\times 4`$ matrix. Since the braiding is a morphism we must have $`\mathrm{\Psi }_{WW}(hk)=h(\mathrm{\Psi }_{WW}k)`$ for the five generators $`h`$ and all $`kWW`$. The cases for $`h=H`$ and $`X_{+1}`$ are given below:
$`M(x,y)(H(y)I_2+I_2H(x))`$ $`=`$ $`(H(x)I_2+I_2H(y))M(x,y),`$
$`M(x,y)(X_{+1}(y)q^{H(x)/2}+q^{H(y)/2}X_{+1}(x))`$ $`=`$ $`(X_{+1}(x)q^{H(y)/2}+q^{H(x)/2}X_{+1}(y))M(x,y).`$
A simple calculation will show that these five conditions determine the matrix $`M(x,y)`$ up to a complex multiple, and we find
$$(\mathrm{\Psi }_{WW}k)(x,y)=a(x,y)\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& \frac{xy(q^21)}{q^2x^2y^2}& \frac{q(x^2y^2)}{q^2x^2y^2}& 0\\ 0& \frac{q(x^2y^2)}{q^2x^2y^2}& \frac{xy(q^21)}{q^2x^2y^2}& 0\\ 0& 0& 0& 1\end{array}\right)k(y,x),$$
where $`a(x,y)`$ is complex valued. Note that $`(\mathrm{\Psi }_{WW}^2k)(x,y)=a(x,y)a(y,x)k(x,y)`$. In the same manner we can determine the braiding $`\mathrm{\Psi }_{W^{}W}:W^{}WWW^{}`$ to be $`(\mathrm{\Psi }_{W^{}W}k)(x,y)=N(x,y)k(y,x)`$, where the matrix $`N(x,y)`$ is given by
$$(\mathrm{\Psi }_{W^{}W}k)(x,y)=\frac{c(x,y)}{q^2y^2x^2}\left(\begin{array}{cccc}q^2y^2x^2& 0& 0& (q^21)xy\\ 0& 0& q(y^2x^2)& 0\\ 0& q(y^2x^2)& 0& 0\\ (q^21)xy& 0& 0& q^2y^2x^2\end{array}\right)k(y,x),$$
where $`c(x,y)`$ is another complex valued function.
Now we use the fact that as the braiding is functorial, it must commute with the evaluation morphism. We see that the maps $`(I\mathrm{eval})(\mathrm{\Psi }_{W^{}W}I)(I\mathrm{\Psi }_{WW})`$ and $`\mathrm{eval}I:W^{}WWW`$ are the same. In terms of the standard pictures, this is fig. 2.
Now, identifying $`W^{}WW`$ with maps from subsets of $`(C^{})^3`$ to $`C^2C^2C^2`$, we get
$`((\mathrm{\Psi }_{W^{}W}I)(I\mathrm{\Psi }_{WW})k)(x,y,z)`$ $`=`$ $`(N(x,y)I_2)((I\mathrm{\Psi }_{WW})k)(y,x,z)`$ (12)
$`=`$ $`(N(x,y)I_2)(I_2M(x,z))k(y,z,x),`$ (13)
and applying $`I\mathrm{eval}`$ to this gives
$$\frac{1}{4\pi i}(I_2\left(\begin{array}{cccc}1& 0& 0& 1\end{array}\right))(N(x,z)I_2)(I_2M(x,z))k(z,z,x)\frac{dz}{z},$$
and some matrix multiplication shows that this is
$$\frac{1}{4\pi i}a(x,z)c(x,z)(\left(\begin{array}{cccc}1& 0& 0& 1\end{array}\right)I_2)k(z,z,x)\frac{dz}{z}.$$
Just applying $`\mathrm{eval}I`$ to $`k`$ gives
$$\frac{1}{4\pi i}(\left(\begin{array}{cccc}1& 0& 0& 1\end{array}\right)I_2)k(z,z,x)\frac{dz}{z},$$
and since these must be the same for all choices of $`k`$ we deduce that $`c(x,z)=1/a(x,z)`$.
Now use the fact that the braiding commutes with the morphism $`\theta :WW^{}`$ in (10), so $`(I\theta )\mathrm{\Psi }_{WW}`$ and $`\mathrm{\Psi }_{W^{}W}(\theta I):WWWW^{}`$ are the same. Then
$`((I\theta )\mathrm{\Psi }_{WW}k)(x,y)`$ $`=`$ $`(I_2\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right))(\mathrm{\Psi }_{WW}k)(x,qy)`$ (14)
$`=`$ $`(I_2\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right))M(x,qy)k(qy,x),`$ (15)
$`(\mathrm{\Psi }_{W^{}W}(\theta I)k)(x,y)`$ $`=`$ $`N(x,y)((\theta I)k)(y,x),`$ (16)
$`=`$ $`N(x,y)(\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)I_2)k(qy,x).`$ (17)
From some more matrix multiplication, this is true if
$`{\displaystyle \frac{1}{a(x,y)}}`$ $`=`$ $`a(x,qy){\displaystyle \frac{x^2q^2y^2}{q(x^2y^2)}}.`$ (18)
## 6 The normalisation of the braiding, $`|q|=1`$.
In this section we find the normalisation $`a(x,y)`$. For the moment we shall suppose that the value of $`x`$ is fixed. Suppose that a solution $`a(x,y)`$ of (18) is an analytic function of $`y`$ (except for isolated singularities) in some annulus centered on zero contained in the region $`|y|>|x|`$. We can narrow the annulus down until it no longer contains any isolated singularities or zeros. For convenience we set $`z=x/y`$ (so $`|z|<1`$), and $`f(z)=a(x,x/z)`$. Then $`f(z)`$ satisfies the equation
$`f(z)f(z/q)`$ $`=`$ $`q{\displaystyle \frac{z^21}{z^2q^2}}.`$ (19)
Now $`f(z)`$ will have a winding number $`\omega Z`$ around zero as $`z`$ winds once around zero. The function $`c(z)=z^\omega f(z)`$ will have zero winding number, so its log, $`b(z)=\mathrm{log}(c(z))`$, will be analytic (and single valued) in the annulus. Now $`c(z)`$ obeys the equation $`c(z)c(z/q)=q(z^2/q)^\omega (z^21)/(z^2q^2)`$, so we must have
$$b(z)+b(z/q)=\mathrm{log}(1z^2)\mathrm{log}(1z^2/q^2)\mathrm{log}(q)\omega \mathrm{log}(z^2/q),$$
as $`z/q`$ is in the annulus if $`z`$ is, because $`|q|=1`$. Since all the other functions are single valued in the annulus, we must have $`\omega =0`$. Now we can take the Laurent expansion of both sides in the annulus, and compare coefficients, to get the unique solution (up to the addition of a multiple of $`\pi i`$)
$`b(z)={\displaystyle \frac{1}{2}}\mathrm{log}(q)+{\displaystyle \underset{n>0}{}}{\displaystyle \frac{1}{n}}{\displaystyle \frac{1q^{2n}}{1+q^{2n}}}z^{2n},|z|<1,z\mathrm{annulus}.`$ (20)
By the uniqueness of analytic continuation, the exponential of this formula must coincide with $`f(z)`$ in a disk from zero up to the radius of convergence of the series. We also conclude that $`f(z)`$ did in fact not have any zeros or isolated singularities in this disk.
In the same manner, if $`f(z)`$ were analytic (except for isolated singularities) in some annulus centered on zero outside the unit disk, we could conclude that on that annulus (again up to the addition of a multiple of $`\pi i`$)
$`b(z)={\displaystyle \frac{1}{2}}\mathrm{log}(q){\displaystyle \underset{n>0}{}}{\displaystyle \frac{1}{n}}{\displaystyle \frac{1q^{2n}}{1+q^{2n}}}z^{2n},|z|>1.`$ (21)
We can now see explicitly from the series (on the assumption that the series converge anywhere) that $`b(z)=b(1/z)`$ plus a multiple of $`\pi i`$, i.e. that $`1/f(z)=\pm f(1/z)`$.
Now we can look at different values of $`x`$. As a consequence of the analyticity condition we assumed for $`a(x,y)`$, it can be seen that $`a(x,y)`$ in fact only depends on $`z=x/y`$. We shall abuse our previous notation by referring to $`a(z)`$.
## 7 The convergence of the series by number theory, $`|q|=1`$.
The radius of convergence $`R`$ of the series (20) is given by
$`R^1=\underset{n+\mathrm{}}{lim\; sup}\left|{\displaystyle \frac{1}{n}}{\displaystyle \frac{1q^{2n}}{1+q^{2n}}}\right|^{1/(2n)}=\underset{n+\mathrm{}}{lim\; sup}\left|{\displaystyle \frac{1q^{2n}}{1+q^{2n}}}\right|^{1/(2n)}.`$ (22)
We see that the series is not even defined if $`q`$ is an even root of $`1`$. For any other root of unity the series has radius of convergence $`1`$, except for $`q=\pm 1`$, when the series terminates. Now consider the case $`q=e^{2\pi i\tau }`$, where $`\tau `$ is irrational.
Call $`nN`$ type 1 if $`|1q^{2n}|<1`$, which implies that $`|1+q^{2n}|1`$. Then
$$\underset{n+\mathrm{},\mathrm{type}\mathrm{\hspace{0.17em}1}}{lim\; sup}\left|\frac{1q^{2n}}{1+q^{2n}}\right|^{1/(2n)}\underset{n+\mathrm{},\mathrm{type}\mathrm{\hspace{0.17em}1}}{lim\; sup}1=1.$$
Call $`nN`$ type 2 if $`|1q^{2n}|1`$, in which case
$$\underset{n+\mathrm{},\mathrm{type}\mathrm{\hspace{0.17em}2}}{lim\; sup}\left|\frac{1q^{2n}}{1+q^{2n}}\right|^{1/(2n)}\underset{n+\mathrm{},\mathrm{type}\mathrm{\hspace{0.17em}2}}{lim\; sup}\left|\frac{1}{1+q^{2n}}\right|^{1/(2n)}1,$$
and we deduce that
$$R^1=\underset{n+\mathrm{},\mathrm{type}\mathrm{\hspace{0.17em}2}}{lim\; sup}\left|\frac{1q^{2n}}{1+q^{2n}}\right|^{1/(2n)}.$$
Then we find
$$\underset{n+\mathrm{},\mathrm{type}\mathrm{\hspace{0.17em}2}}{lim\; sup}\left|\frac{1}{1+q^{2n}}\right|^{1/(2n)}R^1\underset{n+\mathrm{},\mathrm{type}\mathrm{\hspace{0.17em}2}}{lim\; sup}\left|\frac{2}{1+q^{2n}}\right|^{1/(2n)},$$
which implies
$`R^1=\underset{n+\mathrm{},\mathrm{type}\mathrm{\hspace{0.17em}2}}{lim\; sup}\left|{\displaystyle \frac{1}{1+q^{2n}}}\right|^{1/(2n)}=\underset{n+\mathrm{}}{lim\; sup}\left|{\displaystyle \frac{1}{1+q^{2n}}}\right|^{1/(2n)}.`$ (23)
If we let $`d(\tau ,n)`$ be the minimum distance from $`4\pi in\tau `$ to an odd multiple of $`\pi i`$, then $`d(\tau ,n)|1+q^{2n}|\frac{2}{\pi }d(\tau ,n)`$, so
$`R=\underset{n+\mathrm{}}{lim\; inf}\left(\underset{p\mathrm{odd}}{\mathrm{min}}|4\pi in\tau p\pi i|\right)^{1/(2n)}=\underset{n+\mathrm{}}{lim\; inf}\left(\underset{p\mathrm{odd}}{\mathrm{min}}|\tau {\displaystyle \frac{p}{4n}}|\right)^{1/(2n)}1.`$ (24)
We see that the radius of convergence of the power series is dependent on how well $`\tau `$ can be approximated by rational numbers. Fortunately many results are known in this area , and we shall use one of these now.
Suppose that the irrational number $`\tau R`$ is algebraic of degree $`k>1`$ (this means that it is a root of a polynomial of degree $`k`$ with integer coefficients). Then there is a constant $`K>0`$ so that for all $`n`$ and all $`pZ`$,
$$|\tau \frac{p}{4n}|\frac{K}{(4n)^k}.$$
From the formula for $`R`$ above (24),
$$R\underset{n+\mathrm{}}{lim\; inf}\left(\frac{K}{2^k(2n)^k}\right)^{1/(2n)}=1,$$
so we conclude that for any irrational algebraic number $`\tau `$, the radius of convergence is 1.
To get a radius of convergence less than 1, we shall have to create an irrational number with very good rational approximations. Let
$`\tau ={\displaystyle \underset{s1}{}}{\displaystyle \frac{1}{4m_s}},`$ (25)
where the strictly positive integers $`m_s`$ have the property that $`4m_s`$ divides $`m_{s+1}`$ for all $`s1`$. If we set $`n=m_t`$, then
$$\frac{n}{m_{t+1}}|l4n\tau |\frac{2n}{m_{t+1}}$$
for $`l=_{ts1}m_t/m_s`$ an odd integer, so
$$\frac{1}{4m_{t+1}}\underset{p\mathrm{odd}}{\mathrm{min}}|\tau \frac{p}{4n}|\frac{1}{2m_{t+1}}.$$
Then by (24),
$$R\underset{t\mathrm{}}{lim\; inf}\left(\frac{1}{m_{t+1}}\right)^{1/(2m_t)}.$$
If we set $`m_1=1`$ and $`m_{s+1}=2^{2sm_s}`$ for all $`s1`$, then the radius of convergence of the series (20) for $`\tau `$ given by (25) is zero. Also $`\tau `$ is too closely approximated by non-equal rational numbers to be rational itself.
## 8 The convergence of the series by probability, $`|q|=1`$.
We consider the probability that the radius of convergence $`R`$ of the series (20) is 1, given that $`q`$ has a uniform distribution on the circle (equivalently, $`\tau `$ has a uniform distribution on $`[0,1]`$). From (24), for $`0<s<1`$:
$`P[R>s]=P[\underset{n\mathrm{}}{lim\; inf}\underset{p\mathrm{odd}}{\mathrm{min}}|\tau {\displaystyle \frac{p}{4n}}|^{1/(2n)}>s]=\underset{m\mathrm{}}{lim}P[\underset{nm}{inf}\underset{p\mathrm{odd}}{\mathrm{min}}|\tau {\displaystyle \frac{p}{4n}}|^{1/(2n)}>s].`$ (26)
For any random variable $`X_n`$ and $`s<t<1`$,
$$(nmX_n>t)\underset{nm}{inf}X_n>s,$$
or alternatively
$$P[X_n>tnm]P[\underset{nm}{inf}X_n>s].$$
Then from (26),
$$P[R>s]\underset{m\mathrm{}}{lim}P[\underset{p\mathrm{odd}}{\mathrm{min}}|\tau \frac{p}{4n}|^{1/(2n)}>tnm].$$
If we define
$$A_n(x)=\{\tau [0,1]:\underset{p\mathrm{odd}}{\mathrm{min}}|\tau \frac{p}{4n}|x\},$$
then (with superscript $`c`$ denoting complement)
$$P[R>s]\underset{m\mathrm{}}{lim}P[_{nm}A_n(t^{2n})^c].$$
Then by taking complements
$$P[Rs]\underset{m\mathrm{}}{lim}P[_{nm}A_n(t^{2n})]\underset{m\mathrm{}}{lim}\underset{nm}{}P[A_n(t^{2n})],$$
so if the series
$`{\displaystyle \underset{n1}{}}P[A_n(t^{2n})]`$ (27)
converges, then $`R>s`$ with probability one. If we write
$$A_n(x)=\{\tau [0,1]:\underset{p\mathrm{odd}}{\mathrm{min}}|4n\tau p|4nx\},$$
then, if $`4nx1`$, in the interval $`[0,4n]`$ there are $`2n`$ odd integers $`p`$, and each has a disjoint interval of length $`8nx`$ about it satisfying the inequality above. From this we find $`P[A_n(x)]=4nx`$, so the sum (27) becomes (with the exception of a finite number of terms at the beginning),
$$\underset{n1}{}4nt^{2n},$$
which converges since $`|t|<1`$. We conclude that $`P[R>s]=1`$, and so $`R=1`$ with probability one. Since the algebraic numbers have measure zero, the series must have $`R=1`$ for many transcendental (non-algebraic) $`\tau `$.
## 9 The case $`|q|1`$.
In this case we would get the same unique series if the annulus in which $`a(z)`$ was analytic and free of zeros or singularities was sufficiently wide. We would need to have both $`z`$ and $`z/q`$ in the same annulus, so the ratio of the outer and inner radii of the annulus would have to be greater than the larger of $`|q|`$ and $`1/|q|`$. If this condition was satisfied, we would have the solutions (20) or (21). However there may be other solutions to the normalisation which would always have zeros or singularities in such wide annuli.
Let us examine the series (20) for $`|q|1`$. Here $`lim_{n+\mathrm{}}|(1q^{2n})/(1+q^{2n})|=1`$, so the series has radius of convergence 1. However now we can make an analytic continuation of the series. In the case $`|q|<1`$ we write, for any integer $`k1`$,
$`{\displaystyle \frac{1q^{2n}}{1+q^{2n}}}=1+\mathrm{\hspace{0.17em}2}{\displaystyle \underset{m=1}{\overset{k}{}}}(1)^mq^{2mn}{\displaystyle \frac{2(1)^kq^{2(k+1)n}}{1+q^{2n}}},`$ (28)
and substituting this into (20) gives
$$b(z)=\frac{1}{2}\mathrm{log}(q)\mathrm{log}(1z^2)\mathrm{\hspace{0.17em}2}\underset{m=1}{\overset{k}{}}(1)^m\mathrm{log}(1z^2q^{2m})\underset{n1}{}\frac{2(1)^kq^{2(k+1)n}z^{2n}}{n(1+q^{2n})}.$$
The last term in the formula tends to zero uniformly on any bounded set as $`k\mathrm{}`$, allowing us to take the limit as $`k\mathrm{}`$ of the other terms. We then get the infinite product expansion, valid everywhere in $`C^{}`$,
$`a_+(z)`$ $`=`$ $`{\displaystyle \frac{1}{r}}{\displaystyle \frac{1}{(1z^2)}}{\displaystyle \underset{\mathrm{odd}m1}{}}{\displaystyle \frac{(1z^2q^{2m})^2}{(1z^2q^{2(m+1)})^2}},|q|<1.`$ (29)
A similar rearrangement to (28) for $`|q|>1`$ would give
$`a_+(z)`$ $`=`$ $`{\displaystyle \frac{(1z^2)}{r}}{\displaystyle \underset{\mathrm{odd}m1}{}}{\displaystyle \frac{(1z^2q^{2(m+1)})^2}{(1z^2q^{2m})^2}},|q|>1.`$ (30)
These are the unique (up to a sign) solutions of (19) which are analytic for $`z`$ in a neighbourhood of zero in $`C^{}`$. In the same way we can analytically extend the series (21) to obtain $`a_{}(z)=\pm 1/a_+(1/z)`$, the unique solutions of (18) which are analytic for $`z`$ in a neighbourhood of infinity in $`C^{}`$.
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# Comparison between Geometric Quantisation and Covariant Quantum Mechanics
## 1 Introduction
Some years ago A. Jadczyk and M. M. proposed a covariant formulation of Quantum Mechanics for a scalar particle on a curved spacetime with absolute time, based on non standard methods such as fibred manifolds, jet spaces, non–linear connections, systems of connections, cosymplectic structures and Froelicher smooth spaces . This theory has been extended to spin particles in cooperation with D. Canarutto , further developed in cooperation with J. Janyška, D. Saller, C. Tejero Prieto and R. Vitolo and partially extended to a Lorentz manifold in cooperation with J. Janyška and R. Vitolo .
In the proceedings of the previous session of the meeting on Lie Theory, we have accounted for a summary of this theory . In order to capture the non standard methods and results of this theory it would be useful to compare it with the more standard Geometric Quantisation. This is the goal of the present paper.
For the sake of simplicity, our theory in the Galileian case will be conventionally referred to as Covariant Quantum Mechanics (CQM). Moreover, we shall be concerned with the main thread of Geometric Quantisation (GQ) and omit to consider special approaches dealing with quantisation of cosymplectic structures and so on. We have no pretension at all of analysing extensively the wide literature of Geometric Quantisation; such a task would require a much larger space than a short note. Here, we just try to discuss some basic items concerning the comparison of the above theories.
Acknowledgements: The first author thanks the organizers of the Meeting for the invitation and for the warm hospitality. This work has been partially supported by University of Florence, Italian MURST and Italian GNFM of CNR.
## 2 Covariant Quantum Mechanics
Let us start with a very brief sketch of the skeleton of the theory. For further details the reader can refer, for instance, to .
The covariance of the theory includes also independence from the choice of units of measurements. For this reason, we need a rigorous treatment of this feature and assume the following “positive 1–dimensional semi–vector spaces” over $`IR^+`$ as fundamental unit spaces (roughly speaking they have the same algebraic structure of $`IR^+`$, but no distinguished generator over $`IR^+`$): the space $`𝕋`$ of time intervals, the space $`𝕃`$ of lengths, the space $`𝕄`$ of masses.
Moreover, we assume the Planck constant to be an element $`\mathrm{}𝕋^{}𝕃^2𝕄`$.
We refer to a particle with mass $`m𝕄`$ and charge $`q𝕋^{}𝕃^{3/2}𝕄^{1/2}`$.
### 2.1 Classical theory
The classical framework is described in the following way.
The spacetime is an oriented $`(n+1)`$–dimensional manifold $`𝑬`$ (in the standard case $`n=3`$), the absolute time is an affine space associated with the vector space $`IR𝕋`$, the absolute time map is a fibring $`t:𝑬𝑻`$. We denote fibred charts of spacetime by $`(x^\lambda )(x^0,x^i)`$. The tangent space and the vertical space of $`𝑬`$ are denoted by $`T𝑬`$ and $`V𝑬`$.
A motion is a section $`s:𝑻𝑬`$. The phase space is the first jet space of motions $`J_1𝑬`$ . We denote fibred charts of phase space by $`(x^0,x^i;x_0^i)`$. The absolute velocity of a motion $`s`$ is its first jet prolongation $`j_1s:𝑻J_1𝑬`$. An observer is a section $`o:𝑬J_1𝑬`$ and the observed velocity of a motion $`s`$ is the map $`[o]s:=j_1sos:𝑻𝕋^{}V𝑬`$.
The spacelike metric is a scaled Riemannian metric of the fibres of spacetime $`g:𝑬𝕃^2(V^{}𝑬\underset{𝑬}{}V^{}𝑬)`$. Given a particle of mass $`m`$, it is convenient to consider the re–scaled spacelike metric $`G:=\frac{m}{\mathrm{}}g:𝑬𝕋(V^{}𝑬\underset{𝑬}{}V^{}𝑬)`$.
The gravitational field is a time preserving torsion free linear connection of the tangent space of spacetime $`K^{\mathrm{}}:T𝑬T^{}𝑬\underset{T𝑬}{}TT𝑬`$, such that $`[K^{\mathrm{}}]g=0`$ and the curvature tensor $`R[K^{\mathrm{}}]`$ fulfills the condition $`R^{\mathrm{}}{}_{\lambda }{}^{i}{}_{\mu }{}^{}{}_{}{}^{j}=R^{\mathrm{}}{}_{\mu }{}^{j}{}_{\lambda }{}^{}^i`$.
The electromagnetic field is a scaled 2–form $`f:𝑬(𝕃^{1/2}𝕄^{1/2})\mathrm{\Lambda }^2T^{}𝑬`$, such that $`df=0`$. Given a particle of charge $`q`$, it is convenient to consider the re–scaled electromagnetic field $`F:=\frac{q}{\mathrm{}}f:𝑬\mathrm{\Lambda }^2T^{}𝑬`$.
The electromagnetic field $`F`$ can be “added”, in a covariant way, to the gravitational connection $`K^{\mathrm{}}`$ yielding a (total) spacetime connection $`K`$, with coordinate expression
$`K_i{}_{}{}^{h}{}_{j}{}^{}=K^{\mathrm{}}{}_{i}{}^{h}{}_{j}{}^{},K_j{}_{}{}^{h}{}_{0}{}^{}=K_0{}_{}{}^{h}{}_{j}{}^{}=K^{\mathrm{}}{}_{0}{}^{h}{}_{j}{}^{}+\frac{1}{2}F^h{}_{j}{}^{},K_0{}_{}{}^{h}{}_{0}{}^{}=K^{\mathrm{}}{}_{0}{}^{h}{}_{0}{}^{}+\frac{1}{2}F^h{}_{0}{}^{}.`$
This turns out to be a time preserving torsion free linear connection of the tangent space of spacetime, which still fulfills the properties that we have assumed for $`K^{\mathrm{}}`$.
The fibring of spacetime, the total spacetime connection and the spacelike metric yield, in a covariant way, a 2–form $`\mathrm{\Omega }:J_1𝑬\mathrm{\Lambda }^2T^{}J_1𝑬`$ of phase space, with coordinate expression
$`\mathrm{\Omega }=G_{ij}^0(dx_0^i(K_\lambda {}_{}{}^{i}{}_{0}{}^{}+K_\lambda {}_{}{}^{i}{}_{h}{}^{}x_{0}^{h})dx^\lambda )(dx^jx_0^jdx^0).`$
This is a cosymplectic form , i.e. it fulfills the following properties: 1) $`d\mathrm{\Omega }=0`$, 2) $`dt\mathrm{\Omega }^n:J_1𝑬𝕋\mathrm{\Lambda }^nT^{}J_1𝑬`$ is a scaled volume form of $`J_1𝑬`$. Conversely, the cosymplectic form $`\mathrm{\Omega }`$ characterises the spacelike metric and the total spacetime connection. Moreover, the closure of $`\mathrm{\Omega }`$ is equivalent to the conditions that we have assumed on $`K`$.
There is a unique second order connection $`\gamma :J_1𝑬𝕋^{}TJ_1𝑬`$, such that $`i_\gamma \mathrm{\Omega }=0`$. We assume the generalised Newton’s equation $`[\gamma ]j_1s=0`$ as the equation of motion for classical dynamics .
We can also obtain this equation by a Lagrangian formalism according to a cohomological procedure in the following way . The cosymplectic form $`\mathrm{\Omega }`$ admits locally potentials of the type $`\mathrm{\Theta }:J_1𝑬T^{}𝑬`$, defined up to a closed form of the type $`\alpha :𝑬T^{}𝑬`$, which are called Poincaré–Cartan forms . Each Poincaré–Cartan form $`\mathrm{\Theta }`$ splits, according to the splitting of $`T^{}𝑬`$ induced by $`J_1𝑬`$, into the horizontal component $`:J_1𝑬T^{}𝑻`$, called Lagrangian, and the vertical component $`𝒫:J_1𝑬V^{}𝑬`$, called momentum. These components are observer independent, but depend on the chosen gauge of the starting Poincaré–Cartan form. On the other hand, given an observer $`o`$, each Poincaré–Cartan form $`\mathrm{\Theta }`$ splits, according to the splitting of $`T^{}𝑬`$ induced by $`o`$, into the horizontal component $`[o]:J_1𝑬T^{}𝑻`$, called observed Hamiltonian, and the vertical component $`𝒫[o]:J_1𝑬V^{}𝑬`$, called observed momentum. Moreover, the horizontal component of $`\mathrm{\Omega }`$, according to the splitting of $`T^{}J_1𝑬`$ induced by $`J_2𝑬`$, is the map $`=G^{\mathrm{}}([\gamma ]):J_2𝑬𝕋^{}V^{}𝑬`$, which turns out to be the Euler–Lagrange operator associated with $``$. We have the coordinate expressions
$`=(\frac{1}{2}G_{ij}^0x_0^ix_0^j+A_ix_0^i+A_0)dx^0,𝒫=(G_{ij}^0x_0^j+A_i)(dx^ix_0^idx^0),`$
and, in a chart adapted to $`o`$,
$`[o]=(\frac{1}{2}G_{ij}^0x_0^ix_0^jA_0)dx^0,𝒫[o]=(G_{ij}^0x_0^j+A_i)dx^i,`$
where $`Ao^{}\mathrm{\Theta }`$.
The cosymplectic form $`\mathrm{\Omega }`$ yields in a covariant way the Hamiltonian lift of functions $`f:J_1𝑬IR`$ to vertical vector fields $`H[f]:J_1𝑬VJ_1𝑬`$; consequently, we obtain the Poisson bracket $`\{f,g\}`$ between functions of phase space. Given an observer, the law of motion can be expressed, in a non covariant way, in terms of the Poisson bracket and the Hamiltonian.
More generally, chosen a time scale $`\tau :J_1𝑬T𝑻`$, the cosymplectic form $`\mathrm{\Omega }`$ yields, in a covariant way, the Hamiltonian lift of functions $`f`$ of phase space to vector fields $`H_\tau [f]:J_1𝑬TJ_1𝑬`$, whose time component is $`\tau `$. In view of our developments in the Quantum Theory, we prove that $`H_\tau [f]`$ is projectable on a vector field $`X[f]:𝑬T𝑬`$ if and only if the following conditions hold: i) the function $`f`$ is quadratic with respect to the affine fibres of $`J_1𝑬𝑬`$ with second fibre derivative $`f^{\prime \prime }G`$, where $`f^{\prime \prime }:𝑬T𝑻`$, ii) $`\tau =f^{\prime \prime }`$. A function of this type is called special quadratic and has coordinate expression of the type
$`f=\frac{1}{2}f^0G_{ij}^0x_0^ix_0^j+f_i^0x_0^i+\stackrel{𝑜}{f},\mathrm{with}f^0,f_i^0,\stackrel{𝑜}{f}:𝑬IR.`$
The vector space of special quadratic functions is not closed under the Poisson bracket, but it turns out to be an $`IR`$–Lie algebra through the covariant special bracket
$`[[f,g]]=\{f,g\}+\gamma (f^{\prime \prime }).g\gamma (g^{\prime \prime }).f.`$
We have the subalgebra of quantisable functions whose time component factorises through $`𝑻`$, the subalgebra of functions whose time component is constant, the subalgebra of affine functions whose time component vanishes and the abelian subalgebra of spacetime functions which factorise through $`𝑬`$. In particular, the Hamiltonian is a quantisable function, the components of the momentum are affine functions and the spacetime coordinates are spacetime functions.
Moreover, the map $`fX[f]`$ turns out to be a morphism of Lie algebras. The coordinate expression of the tangent lift of $`f`$ is $`X[f]=f^0_0f^i_i`$.
### 2.2 Quantum theory
The quantum framework is described in the following way.
A quantum bundle is defined to be a 1–dimensional complex vector bundle over spacetime $`𝑸𝑬`$ equipped with a Hermitian metric $`h:𝑸\underset{𝑬}{\times }𝑸\mathrm{\Lambda }^nV^{}𝑬`$ with values in the complexified volume forms of the fibres of spacetime. We shall refer to normalised local bases $`b`$ of $`𝑸`$ and to the associated complex coordinates $`z`$; accordingly, the coordinate expression of a quantum section is of the type $`\mathrm{\Psi }=\psi b`$, with $`\psi :𝑬`$.
We consider also the extended quantum bundle $`𝑸^{}J_1𝑬`$, by taking the pullback of $`𝑸𝑬`$, with respect to the map $`J_1𝑬𝑬`$. A system of connections of $`𝑸`$ parametrised by the sections of $`J_1𝑬𝑬`$ induces, in a covariant way, a connection of $`𝑸^{}`$, which is called universal . A characteristic property of the universal connection is that its contraction with any vertical vector field of the bundle $`J_1𝑬𝑬`$ vanishes.
A quantum connection is defined to be a connection $`q`$ of the extended quantum bundle, which is Hermitian, universal and whose curvature is $`R[q]=\mathrm{i}\mathrm{\Omega }`$. We stress that $`\frac{1}{\mathrm{}}`$ has been incorporated in $`\mathrm{\Omega }`$ through the re–scaled metric $`G`$. In a chart adapted to the observer $`o`$, the coordinate expression of a quantum connection is locally of the type
$`q_0=[o],q_i=𝒫[o],q_i^0=0,`$
where the choice of the potential $`A[o]`$ is locally determined by $`q`$ and by the quantum base $`b`$. A quantum connection exists if and only if the cohomology class of $`\mathrm{\Omega }`$ is integer; the equivalence classes of quantum bundles equipped with a quantum connection are classified by the cohomology group $`H^1(𝑬,U(1))`$ .
We stress the minimality of our quantum bundle and quantum connection.
Let us assume a quantum bundle equipped with a quantum connection.
Any other quantum object is obtained, in a covariant way, from this quantum structure. The quantum connection lives on the extended quantum bundle, while we are looking for further quantum objects living on the original quantum bundle. This goal is successfully achieved by a method of projectability: namely, we look for objects of the extended quantum bundle which are projectable to the quantum bundle and then we take their projections. Indeed, our method of projectability turns out to be our way of implementing the covariance of the theory; in fact, it allows us to get rid of the family of all observers, which is encoded in the quantum connection (through $`J_1𝑬`$).
J. Janyška has proved that all covariant quantum Lagrangians of the quantum bundle are proportional to
$`\mathrm{L}[\mathrm{\Psi }]=\mathrm{dt}\left(\mathrm{h}(\mathrm{\Psi },\mathrm{i}\overline{}\mathrm{\Psi })+\mathrm{h}(\mathrm{i}\overline{}\mathrm{\Psi },\mathrm{\Psi })(\overline{\mathrm{G}}\mathrm{h})(^{^{^{}}}\mathrm{\Psi },^{^{^{}}}\mathrm{\Psi })+\mathrm{k}\mathrm{r}\mathrm{h}(\mathrm{\Psi },\mathrm{\Psi })\right),`$
where $`k`$ is an arbitrary real factor, $`\overline{}`$ denotes the covariant differential with respect to time induced by the phase space, $`^{^{^{}}}`$ denotes the vertical covariant differential and $`r:𝑬IR𝕋^{}`$ is the scalar curvature of the spacelike metric $`G`$. Thus, $`k`$ remains undetermined in our scheme. Several authors have tried to determine this factor via Feynmann’s path integral approach, but they found different results, according to different ways to perform the integral .
The standard Lagrangian formalism yields, from the above covariant quantum Lagrangian, the covariant quantum $`(n+1)`$–momentum, the covariant Euler–Lagrange equation and the covariant conserved probability current. These objects can also be obtained directly in terms of covariant differentials through the quantum connection, by means of the projectability method. The coordinate expression of the Euler–Lagrange equation is
$`\mathrm{S}.\psi \left(\stackrel{o}{}_0+\frac{1}{2}{\displaystyle \frac{_0\sqrt{|\mathrm{g}|}}{\sqrt{|\mathrm{g}|}}}\right)\psi \mathrm{i}\frac{1}{2}\left(\stackrel{o}{\mathrm{\Delta }}_0\mathrm{k}\mathrm{r}_0\right)\psi =0,`$
where
$`\stackrel{𝑜}{\mathrm{\Delta }}_0G_0^{hk}\stackrel{𝑜}{}_h\stackrel{𝑜}{}_k+K_h{}_{}{}^{k}{}_{}{}^{h}\stackrel{𝑜}{}_{k}^{},\stackrel{𝑜}{}_\lambda _\lambda \mathrm{i}\mathrm{A}_\lambda ,`$
denote the Laplacian and the covariant differential induced by the connections $`q,K`$ and by the observer attached to the spacetime chart. J. Janyška has proved that any covariant Schrödinger equation is of the above type, hence it is the Euler–Lagrange equation associated with a covariant Lagrangian. We assume the above equation as the quantum dynamical equation.
Next, we classify the Hermitian vector fields of the extended quantum bundle, which are projectable to the quantum bundle. We find that the projected vector fields $`Y:𝑸T𝑸`$ of the quantum bundle, called quantum vector fields, constitute a Lie algebra naturally isomorphic to the Lie algebra of quantisable functions. The coordinate expression of the quantum vector field associated with the quantisable function $`f`$ is
$`Y[f]=f^0_0f^j_j+\left(\mathrm{i}(\mathrm{f}^0\mathrm{A}_0\mathrm{f}^\mathrm{h}\mathrm{A}_\mathrm{h}+\stackrel{o}{\mathrm{f}})\frac{1}{2}\mathrm{div}\mathrm{X}[\mathrm{f}]\right)zz,`$
where $`\mathrm{div}X[f]=f^0\frac{_0\sqrt{|g|}}{\sqrt{|g|}}\frac{_j(f^j\sqrt{|g|})}{\sqrt{|g|}}`$.
The quantum vector field $`Y[f]`$ acts on the sections $`\mathrm{\Psi }`$ of the quantum bundle via the associated Lie derivative $`Z[f]:=\mathrm{i}\mathrm{Y}[\mathrm{f}]_{}`$. In particular, we obtain
$`Z[_0](\mathrm{\Psi })=\mathrm{i}(_0+\frac{1}{2}{\displaystyle \frac{_0\sqrt{|\mathrm{g}|}}{\sqrt{|\mathrm{g}|}}})\psi \mathrm{b},\mathrm{Z}[𝒫_\mathrm{j}](\mathrm{\Psi })=\mathrm{i}(_\mathrm{j}+\frac{1}{2}{\displaystyle \frac{_\mathrm{j}\sqrt{|\mathrm{g}|}}{\sqrt{|\mathrm{g}|}}})\psi \mathrm{b}.`$
Next, we consider the pre–Hilbert functional quantum bundle $`𝑯𝑻`$ over time, whose infinite dimensional fibres are constituted by the sections of the quantum bundle at a given time and with compact support. The quantum dynamical operator $`\mathrm{S}`$ can be regarded as a covariant differential $`[\chi ]`$ of the functional quantum bundle; hence, the quantum Lagrangian yields a lift of the quantum connection $`q`$ of the extended quantum bundle to a connection $`\chi `$ of the functional quantum bundle.
Moreover, we can see that, if $`f`$ is a quantisable function, then
$`\widehat{f}=\mathrm{i}(\mathrm{Y}[\mathrm{f}]_{}\mathrm{f}^{\prime \prime }\mathrm{}[\chi ])`$
is the unique combination of $`Z[f]`$ and $`[\chi ]`$, which yields an operator acting on the fibres of the functional quantum bundle. We have the following coordinate expression
$`\widehat{f}(\mathrm{\Psi })=(\frac{1}{2}f^0\stackrel{𝑜}{\mathrm{\Delta }}{}_{0}{}^{}\mathrm{i}\mathrm{f}^\mathrm{j}\stackrel{o}{}_\mathrm{j}+\stackrel{o}{\mathrm{f}}+\frac{1}{2}\mathrm{k}\mathrm{f}^0\mathrm{r}_0\mathrm{i}\frac{1}{2}{\displaystyle \frac{_\mathrm{j}(\mathrm{f}^\mathrm{j}\sqrt{|\mathrm{g}|})}{\sqrt{|\mathrm{g}|}}})\psi \mathrm{b}.`$
The map $`f\widehat{f}`$ is injective. Moreover, $`\widehat{f}`$ is Hermitian.
We assume $`\widehat{f}`$ to be the Hermitian quantum operator associated with the quantisable function $`f`$. This is our correspondence principle.
We define the commutator of Hermitian fibred operators $`h,k`$ of the functional quantum bundle by $`[h,k]:=\mathrm{i}(\mathrm{hk}\mathrm{kh})`$. Then, for each quantisable functions $`f,g`$, we obtain the formula
$`[\widehat{f},\widehat{g}]=\widehat{[f,g]}+[(g^{\prime \prime }Y[f]_{}f^{\prime \prime }Y[g]_{}),\mathrm{S}].`$
The second term in the above formula is the obstruction for the map $`\widehat{}:f\widehat{f}`$ to be a morphism of Lie algebras. There is any substantial physical reason by which the map $`\widehat{}`$ should be a morphism of Lie algebras? On the other hand, the restriction of the map $`\widehat{}`$ to the subalgebra of affine functions yields an injective morphism of Lie algebras.
The Feynmann path integral formulation of Quantum Mechanics can be naturally expressed in our formalism; in particular, the Feynmann amplitudes arise naturally via parallel transport with respect to the quantum connection . So the Feynmann path integral can be regarded as a further way to lift the quantum connection $`q`$ to a functional quantum connection.
In the particular case when spacetime is flat, our quantum dynamical equations turns out to be the standard Schrödinger equation and our quantum operators associated with spacetime coordinates, momenta and energy coincide with the standard operators.
Therefore, all usual examples of standard Quantum Mechanics are automatically recovered in our covariant scheme.
The above procedure can be easily extended to classical and quantum multi–body systems.
The above covariant theory can be extended to particles with spin; in this way, we obtain a generalised Pauli equation and all that.
Several techniques of the above theory (including the Lie algebra of quantisable functions and the corresponding Lie algebra of quantum vector fields) can be reproduced on a Lorentz manifold in a covariant way in the sense of Einstein. However, we do not know so far how to achieve a Hilbert stuff in this contest. It is possible that this problem has no solution (out of the Quantum Field Theory), as it is commonly believed.
## 3 Comparison
Covariant Quantum Mechanics (CQM) has several points in common with Geometric Quantisation (GQ) . The differences between the two theories arise from the fact that their basic goals are different: quantisation procedure and covariant formulation, respectively.
### 3.1 Quantisation
GQ can be regarded as a general programme aimed at “quantising” a classical system. More precisely, GQ is aimed at establishing a procedure in order to represent an algebra of functions of a classic symplectic manifold into a Hilbert space, according to some reasonable rules.
Perhaps, the original notion of “canonical quantisation” goes back to P. M. Dirac . The first rigorous mathematical formulation of the notion of “geometric pre-quantisation” was due to I. E. Segal ; later J. M. Soriau and B. Kostant founded the Geometric Quantisation. This theory has been refined by several authors, see and . For instance, see , a “full quantisation” of a symplectic manifold $`(M,\omega )`$ is defined to be a pair $`(,\delta )`$ where $``$ is a separable complex Hilbert space and $`\delta `$ is a map taking functions $`fC^{\mathrm{}}(M)`$ to self adjoint operators $`\delta _f`$ of $``$ such that
1. $`\delta `$ is $`IR`$-linear,
2. $`\delta _1=\mathrm{id}_{}`$,
3. $`[\delta _f,\delta _g]=\mathrm{i}\mathrm{}\delta _{\{\mathrm{f},\mathrm{g}\}}`$,
4. if $`𝒜C^{\mathrm{}}(M)`$ is a complete subalgebra, i.e. if its centraliser with respect to the Poisson bracket is $`IR`$, then $`\delta _𝒜`$ acts irreducibly on $``$.
There are no go theorems stating that there are no such quantisations in several situations ; among them we mention the famous Groenewold - Van Hove theorem for the symplectic manifold $`(IR^{2n},\omega )`$. If there is no full quantisation, then one looks for a subalgebra $`𝒪C^{\mathrm{}}(M)`$ to be quantised.
Since the requirement of a quantisation is too restrictive, one defines, as first step, a pre-quantisation by requiring just properties $`(1,2,3)`$. A pre-quantisation, called the Dirac problem, exists for every symplectic manifold, whose symplectic form defines an integer cohomology class.
In some respects, the aim of CQM is not the quantisation of a classical system. More precisely, we are just looking for a covariant formulation of the standard Quantum Theory . On the other hand, any quantum measurement is eventually constituted by classical observations, so we need to consider a classical spacetime as background of the Quantum Theory. Then, this background structure plays an important role in the Quantum Theory. But our heuristic geometric techniques are partially different from those of representations of Lie algebras. We observe also that in our formulation the classical spacetime and its structures, rather than the classical dynamics, determine the Quantum Theory.
Eventually, we do obtain a correspondence principle, which is a consequence of a classification theorem and not a postulate. But, this can be regarded as a quantisation only partially.
### 3.2 Covariance
The standard literature on GQ is not concerned with the special or general relativistic covariance of the theory. Indeed, the language of GQ is geometric, hence coordinate free; but this is not sufficient for attaining the covariance. In fact, in the standard literature on GQ, a given frame of reference is implicitly assumed.
On the other hand, CQM looks for a formulation of standard Quantum Mechanics, which be manifestly covariant (with respect to all frames of references, including accelerated frames), in the spirit of General Relativity.
Indeed, it would be natural to take a curved Lorentz manifold as classical background spacetime for such a theory. However, it is well known that there are serious physical difficulties to formulate the Special or General Relativistic Quantum Mechanics in this framework; actually, these difficulties led to the Quantum Field Theory. On the other hand, we realised that it is possible to keep the framework of Quantum Mechanics (Schrödinger and Pauli equations and all that) and formulate our covariant theory in a curved spacetime with absolute time and spacelike Riemannian metric. This approach stands in between the standard non relativistic Quantum Mechanics and a possible general relativistic Quantum Mechanics. In fact, our classical spacetime supports accelerated frames and several features of General Relativity (including the geometric interpretation of the gravitational field), but misses all features strictly related to the Lorentz metric (including the finite speed of signals).
In the flat case, this setting allows us to recover the standard non relativistic Quantum Mechanics. In the curved case, it suggests several new interpretations and techniques which might be possibly useful for a “true” general relativistic theory.
Thus, the covariance is the leading principle of our theory. Indeed, it turns out to be a powerful heuristic guide, as in all general relativistic theories. All main differences between our theory and GQ are related to the covariance.
Our covariant approach can be compared with a large literature dealing with Galileian General Relativity and covariant formulations of Quantum Mechanics .
### 3.3 Generality
The starting programme of GQ is quite general and is based on weak assumptions. In fact, GQ deals just with a symplectic manifold without further structure. On the other hand, strong symmetries of the framework are usually assumed and specified case by case.
CQM starts with a fibred manifold equipped with a spacelike metric and a fibre preserving linear connection, which fulfill a natural condition. In our opinion, the above geometric structure well reflects the physical features occurring in all examples of interest for Quantum Mechanics and yields in a functorial way any further object which is needed for the development of the classical theory (including the cosymplectic structure).
Thus, the CQM deals with a type of model more specific than that of GQ. On the other hand, the large generality of GQ is a beautiful mathematical feature, which, in practice, cannot be physically implemented in full extent. Actually, perhaps all concrete examples of physical interest that can be treated in the framework of GQ can be regarded as particular cases of our model.
### 3.4 Role of time
In GQ, time is essentially an exterior parameter. This theory basically deals with classical and quantum systems which do not depend explicitly on time. So, the starting classical configuration space is a manifold $`𝑺`$ which does not “include” time. If the theory needs to consider time, then it refers to the product manifold $`𝑬𝑻\times 𝑺`$; the fact that spacetime is a product manifold means that a global observer has been implicitly chosen.
In CQM, the requirement that the theory be observer independent imposes that spacetime “includes” time but be not naturally split into space and time.
In Einsteinian General Relativity, spacetime $`𝑬`$ yields no observer independent time $`𝑻`$ and space $`𝑺`$, hence we have no observer independent projections $`𝑬𝑻`$ and $`𝑬𝑺`$. In our Galileian General Relativity, spacetime $`𝑬`$ is equipped with an observer independent time $`𝑻`$ and projection $`𝑬𝑻`$, but we have no observer independent space $`𝑺`$ and projection $`𝑬𝑺`$. In non relativistic GQ, spacetime $`𝑬`$ is equipped with time $`𝑻`$, space $`𝑺`$, and projections $`𝑬𝑻`$ and $`𝑬𝑺`$.
The further developments of our theory respect the starting assumption on the existence of absolute time without a preferred splitting of spacetime. Thus, all peculiar features of our theory follow from the covariance through the role of time. In particular, the cosymplectic structure of our phase space, the universality of the quantum connection, the method of projectability, the absence of the problem of polarisations and the construction of classical and quantum Hamiltonians are related to the role of time.
### 3.5 Phase space
In GQ the phase space is, in principle, any manifold supporting a symplectic form. Usually, the cotangent manifold of a manifold plays the role of phase space in virtue of the fact that it carries a canonical symplectic form. Thus, phase space has even dimension.
In CQM phase space is constituted by the first jet of sections of the spacetime fibred over time . This choice is essential for the covariance of the theory. Thus, the phase space has odd dimension. Actually, the techniques related to even and odd phase spaces, respectively, present important differences.
On the other hand, any observer induces an affine fibred isomorphism of the first jet bundle with the vertical tangent bundle of spacetime (up to a time–scale factor); moreover, the spacelike metric induces a linear fibred isomorphism of the vertical tangent bundle with the vertical cotangent bundle of spacetime. Thus, breaking the covariance, the choice of an observer and the reference to the spacelike metric allow us to compare our phase space with that of GQ.
### 3.6 Symplectic and cosymplectic structures
In GQ the basic geometric structure of Classical Mechanics is constituted by a symplectic form and a Hamiltonian function of phase space. In principle, nothing else is necessary; in practice, one adds the fibring of the phase space over the configuration space and a suitable group of symmetry.
In CQM the geometric structure of Classical Mechanics is constituted by the spacetime fibred over time, the spacelike metric and the spacetime connection. These objects yield a cosymplectic form in a covariant way.
Any observer yields, by pullback and vertical restriction, a Riemannian symplectic form of the fibres of the vertical tangent bundle of spacetime. In this way we recover the analogue of the symplectic form of GQ. However, we stress that this symplectic form is not covariant and carries less information than the original cosymplectic form.
Furthermore, the cosymplectic form and the choice of an observer yield a classical Hamiltonian function. Thus, in CQM the cosymplectic form encodes the Hamiltonian (but an observer is needed to extract it).
### 3.7 Classical Lie brackets
In the original programme of GQ the classical Lie algebra to be quantised is the Poisson Lie algebra of all functions of the phase space, whose bracket is associated with the symplectic form. Actually, we have mentioned before that some obstructions to this quantisation programme occur, but there is no subalgebra $`𝒪`$ that can be consistently considered for all cases.
In CQM we do define the Poisson Lie algebra of all functions of phase space, whose bracket is associated with the cosymplectic form. However, we are only partially interested in this algebra. Indeed, we exhibit a new Lie algebra of special quadratic functions, which is involved in our Quantum Theory. This Lie algebra includes all functions which are usually quantised in the standard approaches.
### 3.8 Symmetries
In GQ the conserved quantities associated with the group of symmetries of the classical system are not necessarily quantisable. For a system whose phase space is a co-adjoint orbit for some group, one possibly imposes that the generators of the group are quantisable and act irreducibly on the Hilbert space. This is done in order to establish a correspondence between “elementary systems” at the classical and quantum levels, and can be considered as a sort of irreducibility condition.
In CQM there is no need for any specific group of symmetries acting on the classical system. Actually, the procedure of quantisation does not depend on such a group. On the other hand, the possible classical symmetries yield interesting consequences, including the momentum map . In particular, all symmetries of the classical structure yield conserved quantisable functions .
### 3.9 Quantum structure
In GQ the quantum bundle is assumed to be a Hermitian line bundle over phase space. Moreover, the quantum bundle is assumed to be equipped with a Hermitian connection whose curvature is proportional to the classical symplectic form.
In CQM the quantum bundle is assumed to be a Hermitian line bundle over spacetime. Moreover, the extended quantum bundle is assumed to be equipped with a Hermitian connection whose curvature is universal and proportional to the classical cosymplectic form.
Thus the novelties of CQM consist in the following minimal assumptions: the quantum bundle lives on spacetime and not on the phase space, the quantum connection is universal.
### 3.10 Polarisation and projection method
In GQ one realises that the base space of the quantum bundle is too big in order to obtain an irreducible representation of the classical Poisson Lie algebra and to fulfill the uncertainty principle. Then, one looks for a polarisation $`P`$, that is for a Lagrangian subbundle $`P`$ of the complexified tangent bundle of the phase space $`M`$, such that $`D_{}=P\overline{P}`$ has constant rank and $`P`$, $`P+\overline{P}`$ are closed under the Lie bracket; moreover, the polarisation is said to be reducible if the quotient of the phase space $`M`$ by the distribution $`D`$ exists and the canonical projection $`\pi :MM/D`$ is a submersion. Once a polarisation is chosen, we can consider the polarised sections of the quantum bundle, that is the sections whose covariant derivative with respect to every vector field of the polarisation vanish. The polarised sections should yield the Hilbert space with the correct size. Actually, the problem of finding polarisations is very hard in practice, should be faced case by case and leads to a lot of complications and ambiguities, where the beauty of the original programme misses over considerably.
In CQM we have, in a covariant way, an implicit natural polarisation, namely the vertical polarisation.
The quantum connection is the only source of all further quantum objects, such as quantum dynamical equation, probability current, quantum operators and so on.
On the other hand, the quantum connection lives on the extended quantum bundle, while we are looking for further quantum objects living on the original quantum bundle. This goal is successfully achieved by the method of projectability.
Thus, in simple words, the difficult search for the inclusion of a polarisation (which should be performed case by case) is substituted by the easy search for projectable objects (which is successful and can be performed in general).
### 3.11 Half–densities and half–forms
In GQ, once a polarisation $`P`$ is chosen, we take the polarised sections of the quantum bundle. The problem is how to build a Hilbert space. For instance, if $`P`$ is reducible, then the Hermitian product $`h(s_1,s_2)`$ of two polarised sections can be understood as a function on the quotient space $`M/D`$; but a problem arises because $`M/D`$ has no natural volume element, which is necessary to define the Hilbert space of $`L^2`$–polarised sections. One way to remedy this problem is to tensor out the sections by the half–densities associated to $`D`$ and to use the natural partial flat connection of $`D`$ to build the Hilbert space. Unfortunately, even for the simplest physical systems, the results do not agree with those found in Quantum Mechanics.
Therefore, a further modification of the theory is needed: here is where half-forms come into play. They are defined through a metaplectic structure for the bundle of frames of the polarisation $`P`$; this imposes new conditions, since metaplectic structures do not exist in general. Even further problems arise, since it may happen that the tensor product of the quantum bundle with the bundle of half–forms has no polarised section at all.
In CQM we do not see any trouble concerning half–densities and all that. It seems that this convenient feature depends on the fact that CQM gets rid of all problems concerning polarisations.
In the first version of CQM the theory assumed a $``$–valued Hermitian metric of the quantum bundle. In this case the theory, in order to prove that quantum operators are Hermitian, needed to use half–densities.
In the recent version of CQM the theory assumes a $`(\mathrm{\Lambda }^nV^{}𝑬)`$–valued Hermitian metric of the quantum bundle. In this case the theory uses just sections of the quantum bundle and does not need half–densities at all.
### 3.12 Schroedinger equation
In GQ there is no clear Schrödinger equation in an explicit formulation ready to be directly compared with the Schrödinger equation of standard Quantum Mechanics. The Schrödinger equation in this framework is understood as the infinitesimal generator of the flow of the Hamiltonian acting on the Hilbert space. The problem is that in general the Hamiltonian does not preserve the chosen polarisation and therefore we are led to compare different polarisations, this is handled through the Blattner-Kostant- Sternberg kernel machinery, see , which is extremely cumbersome. Even here one encounters new problems since there are obstructions discovered by Blattner in order to construct the BKS kernel.
In CQM we have an explicit, covariant and intrinsic Schrödinger equation, which is immediately comparable with the standard Schrödinger equation (see also ). Moreover, we prove that it comes from a quantum Lagrangian.
### 3.13 Hilbert space and Hilbert bundle
GQ and standard Quantum Mechanics deal with a Hilbert space, which usually consists of $`L^2`$ sections of the quantum bundle.
CQM deals with a Hilbert bundle over time. This formulation is unusual but does not seem to be really in contrast with standard Quantum Mechanics.
Once more, this novelty is related to the explicit role of time in the CQM.
### 3.14 Feynmann path integral
In CQM the Feynmann amplitudes appear very naturally in terms of parallel transport with respect to the quantum connection. In fact, the classical Lagrangian turns out to play the role of local symbol of the quantum connection of the extended quantum bundle over time.
The proof of the equivalence of the Feynmann path integral with the covariant Schrödinger equation has not yet been worked out in detail.
### 3.15 Energy
Most of the practical difficulties of GQ run around the quantisation of energy. Actually, the classical Hamiltonian function has to be explicitly postulated and is not encoded in the basic structure of spacetime. Moreover, the energy requires a treatment quite different from the simpler approach required by other observables such as spacetime coordinates and momentum.
Even in CQM the energy has a special role, but no hard problems arise in this respect. First of all, we stress that energy is encoded in the basic geometric structure of spacetime and that it appears explicitly, in a non covariant way, by means of the choice of an observer.
If one accepts the point of view of Covariant Quantum Mechanics, maybe one can understand the difficulties of GQ from this perspective. In fact, in practice what is done in GQ is to take the vertical tangent (or cotangent) space of spacetime as phase space, instead of the first jet space; accordingly, GQ tries to formulate and quantise energy by methods related to vertical subspace. On the other hand, in CQM, it is very clear that energy is related to the horizontal aspect of phase space. We could roughly say that the vertical aspect of phase space is essentially related to the static geometry of spacetime, while the horizontal aspect is related to the dynamics. This observation can also be analysed in terms of the Lie algebra of quantisable functions and their tangent lift; actually, the quantisable functions dealing with the vertical aspects of spacetime constitute the subalgebra of affine functions, while energy is a quadratic function. In CQM, all quantisable functions can quantised on the same footing.
### 3.16 Examples
In GQ it is not granted that every reasonable physical example can be worked out. Actually, there are few examples that have been successfully solved.
CQM reduces to standard Quantum Mechanics in the flat case. Hence, in CQM, all standard physical examples can be formulated; moreover, the Schrödinger equation and quantum operators corresponding to the quantisable functions can be explicitly and immediately computed. Of course, the integration of the Schrödinger equation and the computation of the energy spectrum is an analytical question which should be faced case by case.
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# Stability of the Smectic Quantum Hall State: A Quantitative Study
\[
## Abstract
We present an effective elastic theory which quantitatively describes the stripe phase of the two-dimensional electron gas in high Landau levels ($`N2`$). The dynamical matrix is obtained with remarkably high precision from the density-density correlation function in the time-dependent Hartree-Fock approximation. A renormalization group analysis shows that at $`T=0`$, as the partial filling factor $`\mathrm{\Delta }\nu \nu \nu `$ moves away from $`1/2`$, the anisotropic conducting state may undergo quantum phase transitions: stripes may get pinned along their conducting direction by disorder, or may lock into one another to form a two-dimensional crystal. The model predicts values of $`\mathrm{\Delta }\nu `$ for each transition. The transitions should be reflected in the temperature dependence of the dissipative conductivity.
\]
In a strong magnetic field perpendicular to the plane of a two-dimensional (2D) electron or hole system, the energy spectrum is characterized by discrete Landau levels (LLs) separated by the cyclotron energy $`\mathrm{}\omega _c=\mathrm{}eB/mc`$ and the Zeeman energy $`g\mu _BB`$. Since the degeneracy of each LL, $`𝒩_\varphi `$, is proportional to the magnetic field, for sufficiently strong magnetic fields, only a small number of low-lying LLs are occupied. In this situation, the kinetic energy is practically quenched, sometimes leading to interesting strongly correlated liquid ground states . For example, it is now very well established that the ground state of a sufficiently clean system is given by the fractional quantum Hall fluid if the filling factor $`\nu =𝒩/𝒩_\varphi `$ ($`𝒩=\text{number of electrons or holes}`$), which measures how many LLs are filled, is close to some rational numbers such as $`1/3`$, $`1/5`$, $`2/3`$, $`2/5`$, etc.
However, recent experiments on high-mobility 2D electron and hole systems have revealed that the transport properties are qualitatively different for higher filling factors ($`\nu >4`$ for electrons and $`\nu >2`$ for holes.) One of the most remarkable findings is that the longitudinal resistivities are highly anisotropic near half integer filling factors. As a natural explanation, it has been suggested that the electrons form unidirectional charge density waves or “stripes”, predicted earlier by Koulakov et al. and Moessner and Chalker . An important theoretical development was the observation that the states of this system may be classified by their symmetries, which are highly analogous to those of liquid crystals. The possible states include stripe crystals, smectic, and nematic phases . Of these, the smectic shows promise of explaining the anisotropic transport data at very low temperature .
While some progress has been made in understanding the temperature dependence of transport in this system , the true groundstate when quantum fluctuations are included remains a subject of debate. It has been shown that continuous quantum phase transitions may occur among different possible states when the system parameters are varied. To assess whether this happens requires a knowledge of the parameters entering the effective theory from microscopic calculations. In one such study, MacDonald and Fisher, using an elastic edge state model, found that both interstripe locking interactions and pinning by weak random disorder are relevant perturbations at all $`\mathrm{\Delta }\nu `$ . However, they argued that both perturbations are extremely small in any experimentally accessible systems, and computed the anisotropic resistivities using a semiclassical Boltzmann transport theory.
In this work, we focus on filling factors away from $`\mathrm{\Delta }\nu =1/2`$. In recent work , it was found that in the Hartree-Fock approximation, uniform stripe states are unstable to the formation of density modulations along the stripes — i.e., to the formation of a stripe crystal. Interestingly, the collective modes of the modulated stripe state obtained from a time-dependent Hartree-Fock approximation (TDHFA) shows that the motion of density deformations in the low energy modes looks strikingly similar to that of phonon modes in a highly anisotropic lattice (see Fig. 1). Motivated by this observation, we develop an effective elastic lattice model which provides a quantitative description of the low temperature behavior of the stripe phase. As described below, from the density-density correlation function in the TDHFA, we may numerically obtain a dynamical matrix which reproduces the low-energy microscopic behavior of the system with remarkably high precision. Using the result, we perform a renormalization group (RG) analysis on two important perturbations: interactions among stripe modulations, which presumably stabilize a stripe crystal, and disorder, which may pin the electrons and render the smectic insulating. In contrast to uniform stripes , we find for modulated stripes that, both these perturbations are irrelevant close to $`\mathrm{\Delta }\nu =1/2`$, and that they become relevant at different critical values as $`\mathrm{\Delta }\nu `$ moves away from 1/2. Either of these Kosterlitz-Thouless (KT) transitions represent a metal-insulator transition, and could be identified experimentally by measuring the filling factor dependence of the activation gap on the insulating (i.e., quantized Hall effect) side of the transition.
In our simple lattice model, we assume that each lattice site $`𝐑`$ is occupied by an object with a local density profile or a “form factor” $`f(𝐫)`$, and that the only dynamic variable is the displacement $`𝐮(𝐑)`$ of the object center. The elastic energy associated with the displacements may be described by a Hamiltonian in Fourier space as
$$H_0=\frac{1}{2}\underset{𝐪}{}\underset{\alpha \beta }{}u_\alpha ^{}(𝐪)D_{\alpha \beta }(𝐪)u_\beta (𝐪),(\alpha ,\beta =x,y).$$
(1)
Since the Hamiltonian is Hermitian and invariant under a rotation by $`\pi `$, all dynamical density matrix elements are real and $`D_{xy}(𝐪)=D_{yx}(𝐪)`$, leaving us only three independent real parameters per wave vector. Below, we will treat $`D_{\alpha \beta }`$ as a fitting parameter, chosen to match the density-density correlation function obtained from the TDHFA of the underlying microscopic fermion model in the partially filled uppermost LL.
Since the density-density correlation function is computed at the one-phonon level in the TDHFA , a proper matching procedure requires the corresponding quantity in the lattice theory to be defined also within a harmonic approximation. More specifically, assuming $`𝐪𝐮`$ is small, the density fluctuation operator in the lattice model may be written as
$`\delta n(𝐪)n(𝐪)n(𝐪)`$ (2)
$`f(𝐪)\left[i𝐪𝐮(𝐪)+{\displaystyle \underset{𝐆0}{}}{\displaystyle 𝑑𝐫e^{i𝐆[𝐫𝐮(𝐫)]i𝐪𝐫}}\right],`$ (3)
where $`𝐆`$ is a reciprocal lattice vector and $`\mathrm{}`$ denotes an average over the above elastic Hamiltonian. Ignoring the higher-order second term, the density-density correlation function may be written as
$`\chi (𝐪,𝐪^{},\tau )𝒯\delta n(𝐪,\tau )\delta n(𝐪^{},0)`$ (4)
$`={\displaystyle \frac{f(𝐪)f^{}(𝐪^{})}{(𝒩ab)^2}}{\displaystyle \underset{\alpha \beta }{}}q_\alpha q_\beta ^{}𝒯u_\alpha (𝐪,\tau )u_\beta ^{}(𝐪^{},0),`$ (5)
where $`𝒯`$ is the time ordering operator and $`a`$ and $`b`$ are the lattice constants as in Fig. 1. Due to the lattice symmetry, $`u_\alpha (𝐪)`$ is defined only within the first Brillouin zone, and the above expression vanishes unless $`𝐪^{}=𝐪+𝐆`$ for any reciprocal lattice vector $`𝐆`$.
Using the fact that the single-LL-projected displacement operators obey $`[u_x(𝐪),u_y(𝐪^{})]=il_B^2\delta _{𝐪+𝐪^{},0}`$ ($`l_B=\sqrt{\mathrm{}c/eB}`$ is the magnetic length), the above correlation function is easily computed, and may be written in terms of the Matsubara frequency $`i\omega _n`$ as
$`\chi (𝐪+𝐆,𝐪+𝐆^{},i\omega _n)=`$ (6)
$`{\displaystyle \frac{W(𝐪;𝐆,𝐆^{})}{i\omega _n\omega _𝐪}}{\displaystyle \frac{W^{}(𝐪;𝐆,𝐆^{})}{i\omega _n+\omega _𝐪}},`$ (7)
where $`W(𝐪;𝐆,𝐆^{})`$ are weights that depend on $`D_{\alpha \beta }`$, and $`\omega _𝐪=l_B^2\sqrt{D_{xx}D_{yy}D_{xy}^2}`$ is the eigenmode frequency.
For each wave vector $`𝐪`$ in the first Brillouin zone, we need to compute $`\chi `$ for at least three reciprocal lattice vectors to fit all unknown parameters in the problem. Assuming the stripes are along the $`x`$-axis as in Fig. 1, we use $`𝐆,𝐆^{}=0,\pm 𝐆_0`$ ($`𝐆_0\widehat{𝐲}2\pi /a`$), to find the three dynamical matrix elements $`D_{xx}`$, $`D_{xy}`$, and $`D_{yy}`$, as well as the three form factors $`f(𝐪)`$, and $`f(𝐪\pm 𝐆_0)`$. (Due to the inversion symmetry, the $`f`$’s are real.) For concreteness, we set $`\nu =6`$ (the third lowest LL) in our calculations. In our matching procedure, we focus on the $`3\times 3`$ Hermitian matrix $`W(𝐪;𝐆,𝐆^{})`$ with $`𝐆,𝐆^{}=0,\pm 𝐆_0`$ provided to us by the TDHFA. Together with the mode frequency $`\omega _𝐪`$, this gives ten parameters, of which we use six to produce the dynamical matrix and form factors. (In practice, it is easiest to use the off-diagonal elements of the Hermitian matrix.) The remaining four are used to estimate an errorbar and to examine the validity of the fitting. For example, Fig. 2 shows $`\omega _𝐪`$ obtained from the TDHFA and the lattice theory. For all the quantities, the agreement is quite impressive. The relative error is typically $`10^8`$ unless $`\omega _𝐪`$ is very small, and always remains below $`10^3`$.
The numerical solution of $`D_{\alpha \beta }`$ agrees remarkably well with a harmonic theory of a 2D charged smectic system with long-range Coulomb interaction such as discussed in Ref. . Specifically, the functional forms of the dynamical matrix elements may be written for small $`𝐪`$ as
$`D_{xx}`$ $``$ $`q_x^2+q_x^2/q,`$ (8)
$`D_{xy}`$ $``$ $`q_xq_y+q_xq_y/q,`$ (9)
$`D_{yy}`$ $``$ $`q_y^2+q_y^2/q+q_x^4,`$ (10)
$`\omega _𝐪`$ $``$ $`q_x\sqrt{(q_y^2+q_x^4)/q},`$ (11)
where the terms proportional to $`1/q`$ are from the long-range Coulomb interaction, the $`q_x^4`$ term from bending energy of each stripe, and all the rest from effective short-range interactions.
Note that $`\omega _𝐪=0`$ at $`q_x=0`$. In fact, the only true gapless mode of the stripe crystal is at $`𝐪=0`$. However, the gaps are extremely small if $`q_x=0`$, for any $`q_y`$ . (In these modes, the stripes slide with respect to one another .) We are thus able to model the results of the TDHFA as a charged smectic, despite the formal presence of gaps away from $`𝐪=0`$. Obviously, these nearly gapless modes greatly affect the low temperature behavior; the RG exponents depend on all $`q_y`$’s up to the first Brillouin zone boundary. Thus, a reliable elastic Hamiltonian must match the microscopic system at short wavelengths perpendicular to the stripes. Along the stripes, however, only the long-wavelength behavior is important.
Using our quantitative elastic theory, we examine some important perturbations and their effect at low temperatures by computing the RG exponents. We first consider locking between stripes, which controls the transition between the smectic and the stripe crystal phases. The perturbative action may be written as
$`S_\lambda =`$ (12)
$`\lambda {\displaystyle \underset{j}{}}{\displaystyle 𝑑x𝑑\tau \mathrm{cos}\frac{2\pi }{b}\left[u_x(x,j,\tau )u_x(x,j1,\tau )\right]}.`$ (13)
Here we represent the displacements in real space, with the $`j`$ labeling the stripes. For convenience, we have adopted a continuum approximation in the $`x`$-axis, which should not affect the small $`q_x`$ behavior. A standard RG analysis yields the flow equation
$$\frac{d\lambda }{d\mathrm{}}=\left(2\frac{x_\kappa }{2}\right)\lambda ,$$
(14)
where
$$x_\kappa \frac{2al_B^4}{b}𝑑q_y(1\mathrm{cos}q_ya)\underset{q_x0}{lim}\frac{D_{yy}(𝐪)q_x}{\omega _𝐪}.$$
(15)
The exponent is plotted as a function of $`\mathrm{\Delta }\nu `$ in Fig. 3. Note that if $`\mathrm{\Delta }\nu 0.38`$, then $`2x_\kappa /2`$ becomes negative and locking is irrelevant .
Another important perturbation is that of weak disorder, which may pin the stripes. We model this with, perturbative action for white noise disorder,
$$S_V=𝑑𝐫𝑑\tau V(𝐫)n(𝐫),$$
(16)
where the disorder average is given by $`\overline{V(𝐫)V(𝐫^{})}=V_0^2ab\delta (𝐫𝐫^{})`$. Using the density in Eq. (3), the RG analysis is straightforward. If the $`𝐆`$ in the second term of Eq. (3) is parallel to the stripes, the flow equation for the most relevant operator $`G_x=2\pi /b`$ is given by
$$\frac{dV_0}{d\mathrm{}}=\left(\frac{3}{2}x_V\right)V_0,$$
(17)
with
$$x_V=\frac{al_B^4}{2b}𝑑q_y\underset{q_x0}{lim}\frac{D_{yy}(𝐪)q_x}{\omega _𝐪}.$$
(18)
As plotted in Fig. 3, if $`\mathrm{\Delta }\nu 0.43`$, the exponent is negative and pinning along the stripes is irrelevant.
The RG analysis is more complicated if $`𝐆`$ is perpendicular to the stripes. We find, for any $`𝐆=(0,G_y)`$, that the free energy $`F\mathrm{ln}𝒟𝐮\mathrm{exp}(S_0S_V)`$ computed using Eq. (11) diverges at low temperatures as $`T^{2/5}`$ for any $`\mathrm{\Delta }\nu `$. This indicates that pinning across the stripes is always relevant. Our interpretation of this is that the stripes will be trapped in channels; however, they are still free to move along the channels so that this would not spoil the phase transitions we found above.
Using particle-hole symmetry for $`\mathrm{\Delta }\nu >1/2`$, the above results imply that for $`0.43\mathrm{\Delta }\nu 0.57`$, the stripes are pinned only in the the direction perpendicular to the stripes at $`T=0`$, supporting the existence of an anisotropic metallic state . This is the main result of this Letter, and is in qualitative agreement with experiment .
Our results suggest that the anisotropic transport properties observed in experiment may well represent a new and unusual metallic state of the quantum Hall system, separated from insulating (quantized Hall) states by quantum phase transitions as a function of $`\nu `$, on either side of $`\mathrm{\Delta }\nu =1/2`$. One direct probe of this is the activation energy $`\mathrm{\Delta }`$ of the diagonal transport coefficients $`\rho _{xx},\rho _{yy}`$ in the quantized Hall state. As the transition is approached from below, the gap is controlled by the growing correlation length $`\xi `$, which for a KT transition has the characteristic form $`\mathrm{\Delta }\xi ^1\mathrm{exp}\{C/\sqrt{|\nu _c\nu |}\}`$, where $`C`$ is a non-universal constant. The observation of such behavior in finite temperature studies away from $`\mathrm{\Delta }\nu =1/2`$ would constitute direct evidence of a new type of phase transition for these systems.
We conclude with a few final remarks. Our simple model was analyzed only in the third lowest LL and it included neither finite thickness of the 2D layer nor more complex effects such as spin-orbit coupling. We expect that the qualitative picture presented will be robust, although the precise location of the phase transitions will surely be affected. Although in our specific model, pinning preempted locking, the order of transitions may well be reversed when these effects are included. We note also that the method we have developed here works very well away from $`\mathrm{\Delta }\nu =1/2`$, but breaks down close to this value. The breakdown occurs for $`|\mathrm{\Delta }\nu 1/2|0.01`$, and is reflected in rapidly growing error bars for our matching procedure in that small range. This is in part caused by the charge motion of the collective modes found in TDHFA becoming more complicated near half-filling, taking on a mixed character of both edge states and lattice phonons. Apparently our simple elastic model does not capture this complicated behavior.
In summary, we have developed a method for generating a quantitative elastic model of quantum Hall stripes from microscopic TDHFA calculations. An RG analysis shows that an anisotropic conducting (smectic) state may be stable against crystallization and pinning near filling factor 1/2, and may be destabilized by either of these away from this filling in a continuous quantum phase transition.
The authors would like to thank Ganpathy Murthy, Allan MacDonald, and Jim Eisenstein for helpful suggestions and discussions. This work was supported by NSF Grant No. DMR-9870681.
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# Parallel Quantum Computation, the Library of Babel and Quantum Measurement as the Efficient Librarian
## Abstract
The complementary roles played by parallel quantum computation and quantum measurement in originating the quantum speed-up are illustrated through an analogy with a famous metaphor by J.L. Borges.
Why there is a quantum speed-up is considered to be an interesting open problem. This letter is an excerpt from the paper titled “Performing Quantum Measurement in Suitably Entangled States Originates the Quantum Computation Speed Up” . That paper is rather lengthy, due to the need of checking that all types of quantum algorithms found so far obey the speed-up mechanism propounded. We shall summarize herebelow the justification of the speed-up provided in .
Let us enter “in media res”. One unnecessary but clarifying step of Simon’s algorithm is to measure the content of register $`v`$ (designated by $`\left[v\right]`$ in the following) in the entangled state
$$\frac{1}{\sqrt{2^n}}\underset{x=0}{\overset{2^n1}{}}|x_a|f\left(x\right)_v.$$
(1)
State (1) is the result of reversibly computing the function $`f\left(x\right)`$ for all possible values of $`x`$, in quantum parallelism; $`a`$ and $`v`$ are two n-qubit registers containing the argument $`x`$ and the respective function $`f\left(x\right)`$; $`|x_a`$ denotes an eigenstate of register $`a`$ (in the measurement basis), etc.. Given $`B=\{0,1\}`$, as well known $`f\left(x\right)`$ is a function from $`B^n`$ to $`B^n`$ with the following properties:
* for any $`x`$, there is one and only one $`x^{^{}}x`$ such that $`f\left(x\right)=f\left(x^{^{}}\right)`$;
* all such $`x`$ and $`x^{^{}}`$ are evenly spaced by a constant $`r`$; namely, for all $`x`$, $`\left|xx^{^{}}\right|=r`$ (we are following the simplified version of Simon’s algorithm);
* given $`x`$, computing $`f\left(x\right)`$ requires poly(n) time, whereas given a value $`\stackrel{\mathrm{\_}}{f}`$ of $`f\left(x\right)`$, finding the two values of $`x`$ such that $`f\left(x\right)=f\left(x^{^{}}\right)=\stackrel{\mathrm{\_}}{f}`$ requires exp(n) time by classical computation – the function is hard to reverse.
The problem is to efficiently find $`r`$ by using a quantum computer that, given $`x`$, computes $`f\left(x\right)`$ in poly(n) time. In fact, this computer has already been used to reach state (1). Because of the character of $`f\left(x\right)`$, measuring $`\left[v\right]`$ in (1) yields an outcome of the form
$$\frac{1}{\sqrt{2}}\left(\right|\stackrel{\mathrm{\_}}{x}_a+|\stackrel{\mathrm{\_}}{x}+r_a)|\stackrel{\mathrm{\_}}{f}_v,$$
(2)
where $`f\left(\stackrel{\mathrm{\_}}{x}\right)=f\left(\stackrel{\mathrm{\_}}{x}+r\right)=\stackrel{\mathrm{\_}}{f}.`$
Without entering into detail, the subsequent part of Simon’s algorithm consists in measuring $`\left[a\right]`$ after performing the Hadamard transform on the state of register $`a`$ in (2). The measurement outcome $`z`$ is such that $`rz=0`$; $`rz`$ denotes the module 2 inner product of the two numbers in binary notation (seen as row matrices). By repeating the overall process a poly(n) number of times, a number of constraints $`rz_i`$ sufficient to identify $`r`$ is gathered with any desired probability of success.
The role played by quantum measurement in originating the speed-up will be discussed by using a special way of comparing quantum and classical efficiency: the quantum computational cost of going from quantum state (1) to quantum state (2) is benchmarked with the classical computational cost of going through their (symbolic) descriptions. Such descriptions can be visualized as the print-outs of (1) and (2), provided that $`x,`$ $`f\left(x\right)`$, $`\stackrel{\mathrm{\_}}{x}`$, etc. are substituted by proper numerical values. This criterion is instrumental to achieving an a-posteriori self evident result.
We shall instrumentally use the following way of thinking (opposite to our view):
> quantum computation can produce a number of parallel outputs exponential in register size, at the cost of producing one output, but this “exponential wealth” is easily spoiled by the fact that quantum measurement reads only one output.
Let us examine the cost of classically deriving description (2) from description (1). This latter can be visualized as the print-out of the sum of 2<sup>n</sup> tensor products. Loosely speaking, two values of $`x`$ such that $`f\left(x_1\right)=f\left(x_2\right)`$, must be exp$`\left(n\right)`$ spaced. Otherwise such a pair of values could be found in poly$`\left(n\right)`$ time by classical “trial and error”.
The point is that the print-out would create a Babel Library<sup>*</sup><sup>*</sup>*From the story “The Library of Babel” by J.L. Borges. effect. Even for a small $`n`$, it would fill the entire known universe with, say, $`\mathrm{}`$ $`|x_1_a|f\left(x_1\right)_v`$ $`\mathrm{}`$ here, and $`\mathrm{}`$ $`|x_2_a|f\left(x_2\right)_v`$ $`\mathrm{}`$ \[such that $`f\left(x_1\right)=f\left(x_2\right)`$\] in Alpha Centauri. Finding such a pair of print-outs would still require exp$`\left(n\right)`$ time. The capability of directly accessing that “exponential wealth” would be frustrated by its “exponential dilution”. This seems to be in match with the baffling feeling inspired by Borges’ story.
The quantum measurement of $`\left[v\right]`$ instead, distills the desired pair of arguments in a time linear in $`n`$, namely in the number of qubits of register $`v`$.It is a basic axiom of quantum measurement theory that the time required to measure $`\left[v\right]`$ is independent of state (1) entanglement – entanglement is interaction free. In fact, it does more than randomly selecting one measurement outcome; by selecting one outcome, it performs a logical operation (selecting the two values of $`x`$ associated with the value of that outcome) crucial for solving the problem. The active role played by quantum measurement in originating the speed-up, complementary to the production of the parallel computation outputs, appears to be self evident.
Ref. also shows that performing or skipping $`\left[v\right]`$ measurement in (1) is equivalent. It also formalizes the active role played by quantum measurement. Given a suitably entangled state before measurement, the constraint that there is a single measurement outcome becomes a set of logical-mathematical constraints that represent the problem to be solved (or the hard part thereof). Satisfaction of such constraints, by the measurement outcome, amounts to having solved the problem. The computational complexity of satisfying these constraints comes from entanglement and is completely transparent to measurement time, which justifies the speed-up.Having found that the speed-up is an observable consequence of the transition from the quantum to the classical world, naturally revamps the quantum measurement problem. After having – so to speak – buried it, its unexpected return might be coldly welcomed; while acknowledging this, we should note that this time we are clearly dealing with a new and striking fact, not with a debatable interpretation of quantum measurement.
Acknowledging the active role played by quantum measurement yields a more realistic vision of what quantum computation is and is not.
It is not, as commonly believed, the quantum transposition of reversible Turing-machine computation, where quantum measurement would only be needed to read the output of a sequential computation process. In fact, we have shown that quantum measurement plays a crucial role in efficiently creating that output. A quantum computation yielding a speed-up is not a reversible computation, although reversibility is of course essential to prepare the state before measurement.
Today, people is looking for non-sequential (e.g. topological) forms of quantum computation. It is therefore a matter of some importance to understand that the current “quantum algorithms” are already non-sequential in character.
It is reasonable to think that detaching the notion of “quantum algorithm” from that of sequential computation – a classical vestige – is a precondition for pursuing further developments at a fundamental level.
For example, let us consider the possibility of exploiting particle statistics symmetrizations to achieve a quantum speed-up. Such symmetrizations can be seen as projections on symmetrical (constrained) Hilbert subspaces. There is no relation between a projection and sequential reversible computation, namely a unitary evolution. If instead quantum computation is (properly) seen as a projection on a constrained Hilbert subspace, which amounts to solving a problem, then we have an analogy with particle statistics symmetrizations to work with.
Thanks are due to the co-Authors of Ref. for their consent to issue an excerpt.
References
1. G. Castagnoli, D. Monti, A. Sergienko, “Performing Quantum Measurement in Suitably Entangled States Originates the Quantum Computation Speed Up”, arXiv: quant-ph/9908015 v2 14 Feb. 2000.
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# The Effect of the Environment on 𝛼-Al2O3 (0001) Surface Structures
\[
## Abstract
We report that calculating the Gibbs free energy of the $`\alpha `$-Al<sub>2</sub>O<sub>3</sub> (0001) surfaces in equilibrium with a realistic environment containing both oxygen and hydrogen species is essential for obtaining theoretical predictions consistent with experimental observations. Using density-functional theory we find that even under conditions of high oxygen partial pressure, the metal terminated surface is surprisingly stable. An oxygen terminated $`\alpha `$-Al<sub>2</sub>O<sub>3</sub> (0001) surface becomes stable only if hydrogen is present on the surface. In addition, including hydrogen on the surface resolves discrepancies between previous theoretical work and experimental results with respect to the magnitude and direction of surface relaxations.
PACS numbers:68.35.Bs,68.35.Md,71.15.Ap
\]
The nature of the corundum surface ($`\alpha `$-Al<sub>2</sub>O<sub>3</sub>) is of considerable importance in a wide variety of technological applications. These range from catalytic supports and thin-film substrates to corrosion and wear protection in mechanical systems. Yet, despite considerable experimental and theoretical efforts over the years, the surface structure, and even the surface stoichiometry, is a matter of strong controversy. The principal reason for this is that it is difficult to prepare clean, uniform surfaces with specific, well-defined structures and stoichiometries. The typical heterogeneity of the surface makes the interpretation of low energy electron diffraction (LEED) and other surface spectroscopic data difficult. In addition, corundum is an insulator which makes the application of STM and other techniques based on electron spectroscopy problematic.
Knowing the structure and stoichiometry of the corundum surface, however, is essential for understanding the electronic, mechanical, and chemical properties which determine its reactivity and performance in various applications. Yet it has only been within the past two years that new experimental techniques have enabled the chemical identification of the surface terminations for these systems. Renaud has reported a $`(1\times 1)`$ structure, prepared by heating in an oxygen-rich atmosphere, which could reasonably be interpreted as being O-terminated . Toofan and Watson found in a tensor LEED experiment both Al- and O-terminated domains in a 2:1 ratio, respectively . Ahn and Rabelais annealed the surface under UHV conditions and determined the detailed structure of the Al-terminated surface using time-of-flight scattering and recoiling spectrometry (TOF-SARS). The sensitivity of TOF-SARS also enabled the detection of hydrogen randomly distributed on the surface which was stable even at an annealing temperature of 1100 C.
Despite recent advances in experimental techniques, however, many questions and inconsistencies remain that theoretical calculations performed to date have not been able to resolve. Corundum ($`\alpha `$-Al<sub>2</sub>O<sub>3</sub>) crystallizes in a structure which can be described by a primitive rhombohedral unit cell with two Al<sub>2</sub>O<sub>3</sub> formula units (10 atoms) or by a conventional hexagonal unit cell with six Al<sub>2</sub>O<sub>3</sub> formula units (30 atoms). For the hexagonal unit cell, the atoms are stacked along the (0001) direction according to the sequence $`R`$-AlAlO<sub>3</sub>-AlAlO<sub>3</sub>-$`R`$, where $`R`$ represents the continuing sequence in the bulk. The oxygen atoms form hcp layers, and the metal atoms fill two thirds of the octahedral sites between these layers.
In the $`R`$-Al-Al-O<sub>3</sub>-$`R`$ corundum structure (0001) stacking sequence, there are three unique stoichiometric slice planes. Yet theoretical methods ranging from the empirical to the ab initio have so far identified only one stable termination stoichiometry for the $`\alpha `$-Al<sub>2</sub>O<sub>3</sub> (0001) surface: an Al-monolayer indicated by AlO<sub>3</sub>Al-$`R`$ in Fig. 1. No stable O-terminated surface has been found, so how can the experimentally observed O-terminations be explained? In addition to the surface stoichiometry, there is considerable disagreement with respect to the magnitude and direction of the surface relaxations. For the Al-terminated surface, experimentally observed relaxations range from -51% to -63%, whereas the density-functional theory (DFT) pseudopotential calculations predict a relaxation of about -87%, and Hartree-Fock gives -40%. Toofan and Watson reported an outward relaxation for their two domain system, which is the opposite of what other experimental and theoretical investigations had concluded. These discrepancies are significant and need to be resolved.
What has not been previously considered in the theoretical treatment of the $`\alpha `$-Al<sub>2</sub>O<sub>3</sub> (0001) surface is the effect of the environment on the surface structures and stoichiometry. Under realistic conditions, a surface will exchange atoms with its surroundings. Hence in this work we present the first analysis of the Gibbs free energy of the system with respect to its dependence on the chemical potentials of the components present in the material and the environment at 0K and 1000K. For a metal oxide such as $`\alpha `$-Al<sub>2</sub>O<sub>3</sub>, the O<sub>2</sub> partial pressure is obviously the most important factor in the analysis, as well as temperature. In addition, the presence of stable hydrogen on the surface also needs to be addressed, as it can be incorporated into the bulk structure during growth, remain from calcination of Al(OH)<sub>3</sub> in the synthesis of $`\alpha `$-Al<sub>2</sub>O<sub>3</sub>, or may result from exposure of the surface to water vapor prior to placement in the UHV chamber. Yet to our knowledge the presence of hydrogen on the corundum surface has not yet been investigated in theoretical studies. In this paper we report that the surface stoichiometries, structures, and properties change significantly depending upon the chemical potential of O<sub>2</sub>, H<sub>2</sub>, and H<sub>2</sub>O gases in equilibrium with the surface at different temperatures. We use the full-potential linearized augmented planewave (FP-LAPW) method to solve the Kohn-Sham equations and calculate the total energies, forces, and chemical potentials for all reasonable $`(1\times 1)`$ corundum (0001) surface geometries. We find that the metal-terminated surface is surprisingly stable across the range of a physically realistic oxygen chemical potential. An oxygen terminated surface becomes stable only if hydrogen is present on the surface, even at partial O<sub>2</sub> pressures which are too high for standard UHV equipment. In addition, including hydrogen on the surface of both the aluminum and oxygen terminated surfaces results in calculated relaxations which agree with the latest experimental results, in both magnitude and direction.
To systematically investigate the $`(1\times 1)`$ (0001) surface which is observed at annealing temperatures below 1250 C, we have generated what we believe are all the possible ($`1\times 1`$) (0001) Al- and/or O-surface terminations of Al<sub>2</sub>O<sub>3</sub>, plus several additional surfaces containing hydroxyl groups. Previous studies which investigated possible ($`1\times 1`$) (0001) surface terminations examined only the three types of structures obtained by simply cleaving the unit cell at unique positions: AlO<sub>3</sub>Al-$`R`$, O<sub>3</sub>AlAl-$`R`$, and AlAlO<sub>3</sub>-$`R`$. For the single Al-terminated layer there are three other possible locations for the Al atom in addition to the corresponding bulk site, as shown in Fig. 1. For the O-terminated surface, other investigations have not considered the possibility of oxygen vacancies which could occur and still maintain a perfect ($`1\times 1`$) surface periodicity. Two additional O-terminated structures can be created by introducing one or two oxygen vacancies per unit cell, indicated by O<sub>2</sub>AlAl-$`R`$ and O<sub>1</sub>AlAl-$`R`$, respectively. Hence, there are a total of at least eight possible ($`1\times 1`$) geometries of the (0001) surface which we investigate to determine their relative stability.
In our calculations, the $`(1\times 1)`$ $`\alpha `$-Al<sub>2</sub>O<sub>3</sub> surface is modeled by a slab which consists of a finite number of layers and is infinite in the plane of the surface. The slabs are repeated periodically along the direction and separated by 10Å of vacuum. The slab contains six oxygen O<sub>3</sub> layers and from ten to fourteen aluminum layers, depending upon the specific surface studied. We carefully tested that the thickness of the vacuum as well as that of the slabs are sufficiently large to ensure that surface-surface interactions through both the vacuum and the slab are negligible. The two surfaces of the slab are identical and inversion symmetry is maintained.
For our total-energy calculations we use the generalized gradient approximation of Perdew et al. for the exchange-correlation potential and the FP-LAPW method as implemented within the WIEN97 program to solve the Kohn-Sham equations. A uniform k-point mesh with ten points is used for the entire surface Brillouin zone. An identical mesh is used for the bulk calculations to ensure consistency. The calculated equilibrium lattice constants for the hexagonal unit cell and the non-symmetry fixed positions of the atoms agree with experimental values to within 0.7%. The calculations give the heat of formation at 0K of the bulk Al<sub>2</sub>O<sub>3</sub>, $`\mathrm{\Delta }H_\mathrm{f}^\mathrm{o}`$=17.37 eV, which is in good agreement with the experimental value of 17.24 eV .
To compare the stability of surfaces with different stoichiometries in chemical and thermal equilibrium with the gas phase and with the bulk, we calculate the Gibbs free energy $`\mathrm{\Omega }`$ of the surface relative to the chemical potential of oxygen. The Gibbs free surface energy $`\mathrm{\Omega }`$ of a slab at temperature $`T`$ and partial pressure $`p`$ is given by $`\mathrm{\Omega }=E^{\mathrm{total}}+\mathrm{\Delta }G^{\mathrm{vib}}N_i\mu _i(T,p)`$, where $`E^{\mathrm{total}}`$ is the scf energy of the slab, $`\mathrm{\Delta }G^{\mathrm{vib}}`$ is the vibrational contribution to the Gibbs free energy, $`N_i`$ the number of $`i`$th type of atoms in the slab, and $`\mu _i(T,p)`$ is the chemical potential of the $`i`$th type of atom at a given temperature and pressure. The sum is over all the types of the atoms in the slab. To calculate the values of the chemical potential for all species we use $`\mu _i(T,p)`$ = $`\mu _i^0`$ \+ $`\mathrm{\Delta }\mu _i(T,p)`$, where 0K is taken as the reference state, and $`\mathrm{\Delta }\mu _i(T,p)`$ is the change in free energy from that reference state to the system at a given temperature and pressure. The 0K values for surfaces and for bulk Al and Al<sub>2</sub>O<sub>3</sub> are taken from our total energy calculations. For the dissociation energies of the H<sub>2</sub>, O<sub>2</sub>, and H<sub>2</sub>O molecules, we use experimental values. The $`\mathrm{\Delta }\mu _i(T,p)`$ values for all crystalline and gas phase species are taken from the JANAF Thermochemical Tables .
As each surface is considered here to be in chemical and thermal equilibrium with the bulk and the environment, the chemical potentials are constrained, so for the clean surface: $`2\mu _{\mathrm{Al}}+3\mu _\mathrm{O}=E_{\mathrm{Al}_2\mathrm{O}_3}^{\mathrm{bulk}}`$, $`\mu _{\mathrm{Al}}<E_{\mathrm{Al}}^{\mathrm{bulk}}`$, and $`\mu _\mathrm{O}<\frac{1}{2}E_{\mathrm{O}_2}^{\mathrm{molecule}}`$, where E$`{}_{}{}^{\mathrm{bulk}}{}_{\mathrm{Al}_2\mathrm{O}_3}{}^{}`$ is the total energy per bulk Al<sub>2</sub>O<sub>3</sub> formula unit. The limits of the chemical potential for oxygen are determined by conditions of equilibrium with relevant oxygen-containing species in the system. The maximum oxygen chemical potential is that of a maximum concentration of O<sub>2</sub> molecules at 0K, which corresponds to O<sub>2</sub> condensing on the surface. This is the reference value of the oxygen chemical potential, which is the zero on the right hand side of the graph in Fig. 2. Since the system is in equilibrium with bulk Al<sub>2</sub>O<sub>3</sub> the minimum $`\mu _\mathrm{O}`$ occurs when $`\mu _{\mathrm{Al}}`$ is at a maximum and $`\mu _\mathrm{O}=\frac{1}{3}E_{\mathrm{Al}_2\mathrm{O}_3}^{\mathrm{bulk}}\frac{2}{3}E_{\mathrm{Al}}^{\mathrm{bulk}}`$. Below this range, aluminum metal would condense on the surface.
For the surfaces where hydrogen is adsorbed, the free energy of the surface is calculated relative to the chemical potential of hydrogen $`\mu _\mathrm{H}`$ in both H<sub>2</sub> and H<sub>2</sub>O, as desorption of hydrogen as either molecular species is possible. The lines for the HO<sub>3</sub>AlAl-$`R`$ and H<sub>3</sub>O<sub>3</sub>AlAl-$`R`$ surface energies calculated with respect to $`\mu _\mathrm{H}`$ in equilibrium with H<sub>2</sub> and H<sub>2</sub>O cross at $`\mu _\mathrm{O}`$ equal to -2.7 eV at 0K. Below -2.7 eV, $`\mu _\mathrm{H}`$ is determined by the equilibrium with H<sub>2</sub>, and above with H<sub>2</sub>O. Below these lines either H<sub>2</sub> or H<sub>2</sub>O would condense on the surface. It should be noted that the only independent variables in our system is $`\mu _\mathrm{O}`$ in equilibrium with bulk Al<sub>2</sub>O<sub>3</sub>, and $`\mu _\mathrm{H}`$. From these the dependent Al, O<sub>2</sub>, H<sub>2</sub> and H<sub>2</sub>O chemical potentials follow.
The results are shown in Fig. 2. The Gibbs free energies per surface area for various $`(1\times 1)`$ geometries of the (0001) surface are displayed as a function of $`\mu _\mathrm{O}`$. The dashed black vertical lines bracket the allowed range of $`\mu _\mathrm{O}`$. For the systems containing just aluminum and oxygen, there is clearly only one overwhelmingly low energy surface stoichiometry: the AlO<sub>3</sub>Al-$`R`$ structure. This is consistent with the previous theoretical results. Our results show, however, that this stability extends across the entire range of the oxygen chemical potential, even up to the limit where an O<sub>2</sub> condensate would form on the surface. The empirical models explain this remarkable stability on the basis of simple ionic considerations, as it is the only structure without a significant surface dipole. The actual situation is a bit more complicated, of course, and the quantum mechanical calculations performed by us and others indicate that the surface aluminum atom relaxes inward until it is practically coplanar with the oxygen layer, and rehybridizes to an $`sp^2`$ orbital configuration which is dramatically stabilized by autocompensation and bonding considerations. That is, the structure relaxes and charge is transferred so that the $`sp^2`$ bonds between the aluminum atoms and the three oxygens are filled, the aluminum $`3p_z`$ orbital perpendicular to the surface is empty, and there are no partial occupancies of dangling bonds at the surface. In our calculations the lowest empty surface state, which is primarily the Al $`3p_z`$ with some Al $`3s`$ and O $`2p`$ character, is 4.5 eV above the Fermi level. This is very high and contributes to the relative chemical inertness of the relaxed Al-terminated surface. Moving the topmost aluminum atom of the Al-terminated surface from its position to one of the three non-bulk positions indicated in Fig. 1, increases the surface energy dramatically from 133 meV/Å<sup>2</sup> to 355, 282, and 228 meV/Å<sup>2</sup> for the three sites, respectively.
The O<sub>3</sub>-terminated surface has a very unfavorable energy, although this decreases sharply as $`\mu _\mathrm{O}`$ increases. This stoichiometry has the highest dipole at the surface, which is reflected in a very large work function of 9.66 eV, compared to 5.97 eV for the Al-terminated surface above. In addition, there are partially filled open valencies on the surface oxygens which are not relieved by relaxation. Our results indicate, however, that this high energy surface can become more stable through oxygen evaporation. The O<sub>2</sub>-terminated surface has a consistently lower free energy than the O<sub>3</sub>AlAl-$`R`$, and the O<sub>1</sub>AlAl-$`R`$ surface is more stable than O<sub>2</sub>AlAl-$`R`$, except at the highest range of the oxygen chemical potential. At 1000K the metal-terminated surfaces become the lowest energy surface under the most oxygen-deficient conditions. This is consistent with what happens at annealing temperatures above 1250 C, where oxygen evaporation leaves behind a metallic aluminum overlayer . The figure does not convey the kinetic barriers required to move from one surface stoichiometry to another, but the relative thermodynamic stability is clear and is consistent with experimental observations.
The large surface dipole and dangling bonds of the O<sub>3</sub>-terminated surface can also be compensated by the addition of hydrogen to the surface. Adding one hydrogen per unit cell to the surface dramatically lowers the free energy and the work function from 9.66 to 7.06 eV. This surface is much more stable than the O<sub>3</sub>-terminated surface at all oxygen partial pressures. Saturating all surface oxygens with hydrogen, indicated by the H<sub>3</sub>O<sub>3</sub>AlAl-$`R`$ line in Fig. 2, results in the lowest free energy and greatest stability across the entire range of physically realistic conditions. The oxygen open valencies are filled, charge transfer occurs from the H to the topmost O layer, and the work function is reduced to 3.44 eV, the lowest of all surfaces examined. The region of negative surface energy of the H<sub>3</sub>O<sub>3</sub>-terminated surface indicates that the hydroxylated surface is lower in energy than bulk $`\alpha `$-Al<sub>2</sub>O<sub>3</sub> plus free H<sub>2</sub>O. This negative surface energy reflects the strength of the H-OAl bond and is consistent with the negative heat of the $`\alpha `$-Al<sub>2</sub>O<sub>3</sub> \+ 3H<sub>2</sub>O $``$ 2Al(OH)<sub>3</sub> reaction. This reaction is exothermic by 0.19 eV per Al<sub>2</sub>O<sub>3</sub> formula unit relative to liquid water at 298.15K and 1 atm pressure .
We find that in addition to explaining the relative stabilities of the (0001) $`(1\times 1)`$ surfaces observed experimentally, considering the presence of hydrogen on the surface is essential to resolve the controversy surrounding the surface relaxations. Our results are shown in Table I. In the absence of hydrogen, our full potential results find a large inward relaxation of 86% for the first layer, in very close agreement with previous pseudopotential calculations, although there is less agreement for subsequent layers. When a hydrogen atom occupies site 1, however, the contraction between the first Al and O layers is $``$ 69$`\%`$, in very close agreement with the experimental result of Ahn and Rabalais . Several other locations for the hydrogen atom were investigated, with similar results.
For the oxygen-only terminated surface, theoretical predictions have indicated an inward relaxation of the top layer of 0.05Å. Yet the experimental results of Toofan and Watson, in which one of the two domains observed was O-terminated, clearly showed an outward relaxation of 0.12Å. Placing a hydrogen at site 1, however, results in a predicted outward relaxation of 0.11Å, in excellent agreement with experiment. Saturating the surface with hydrogen, the H<sub>3</sub>O<sub>3</sub>-structure, increases the predicted relaxation to 0.19Å. Hence even though H was not detectable with the instrumentation used by Toofan and Watson, our free energy calculations and predicted relaxations provide strong evidence that H was most likely present on the surface.
Real surfaces are usually much more complicated and heterogeneous than the model $`(1\times 1)`$ system studied here, but encorporating the effects of chemical and thermal equilibrium with the environment into the model provides a sound framework for interpretation of experimental results on these real systems.
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# QUARK CONDENSATE IN THE DEUTERON
## I introduction
The QCD vacuum has a complex structure, with condensates of quarks and gluons, that can be disturbed by the presence of hadronic matter. In the case of nucleons, for instance, valence quarks give rise to an anti-screening interaction, which reduces the magnitude of the condensate. This gives rise to the nucleon sigma-term ($`\sigma _N`$), that can be extracted from pion-nucleon scattering.
In the case of nuclei, in first approximation the effects of idependent nucleons add up . But as nucleons are interacting, there also exist modifications of the condensate due to the nucleon-nucleon potential. It is reasonable to believe that the influence of this potential is more important in large nuclei, but the study of these systems is complicated and requires simplifying approximations. Therefore it is interesting to look for effects of the NN interaction over the condensate in light nuclei. The deuteron, in particular, has been extensively explored and allows calculations with little theoretical uncertainties.
In principle, one should use QCD to study the reaction of the quark condensate to the presence of hadronic matter. However, as this is beyond our present capabilities, we use effective interactions of colourless hadrons in place of the fundamental ones. Effective theories should be as close as possible to QCD and, in particular, share its symmetries. The interactions of quarks and gluons are approximately invariant under chiral transformations and broken, in the $`SU(2)`$ sector, by the very small quark masses. Therefore, at the hadron level, one requires the effective theory to possess approximate chiral symmetry, now broken by $`\mu `$, the pion mass.
In the case of NN interactions, most of the dynamics relevant at large and intermediate energies can be described, in the framework of effective theories, by exchanges of one and two pions . For the short distance region, on the other hand, neither meson nor quark models produce precise quantitative predictions and realistic potentials must rely on free parameters. In the case of the deuteron, these short distance uncertainties are minimized, for it is heavily dominated by the one pion exchange potential (OPEP) .
In this work we discuss the disturbances of the QCD vacuum produced by the deuteron. In sect. II, we concentrate on the dependence of its binding energy on the quark mass, to derive the quark condensate using the Feynman-Hellmann theorem. The changes induced in the quark condensate by the nuclear force can be related to exchange currents, as probed by means of both photons and pions. Thus, in sect. III we discuss the case of electromagnetic probes and in sect. IV we study $`\pi `$d scattering. Finally, in sect. V we present our results and discuss how they are related to measurable quantities.
## II Feynman-Hellmann
The deuteron mass is written as $`M=2mϵ`$, where m is the nucleon mass and $`ϵ`$ is the binding energy, which we take as positive. The part of $`M`$ due to chiral symmetry breaking corresponds to the deuteron sigma-term, given by
$$\sigma _d=d^3r\left(d\left|_{SB}\right|d0\left|_{SB}\right|0\right),$$
(1)
where $`_{SB}`$ is the symmetry breaking term of the QCD Lagrangian. In the symmetric isospin limit it is given by $`_{SB}=\widehat{m}\overline{q}q`$, where $`q`$ is the SU(2) quark field and $`\widehat{m}`$ is the average quark mass: $`\widehat{m}=(m_u+m_d)/2`$. At leading order in the chiral expansion the effective and fundamental symmetry breaking parameters are related by a constant, denoted by $`B`$: $`\mu ^2=2B\widehat{m}`$. As $`\widehat{m}`$ and $`\mu ^2`$ are small, we have
$$\sigma _d=\widehat{m}\frac{dM}{d\widehat{m}}=\mu ^2\frac{dM}{d\mu ^2}$$
(2)
and write $`\sigma _d=2\sigma _N+\sigma _ϵ`$, where $`\sigma _N=\mu ^2dm/d\mu ^2`$ and $`\sigma _ϵ`$ describes the changes in the condensate as compared to an assembly of static non-interacting nucleons.
In the framework of the Schrödinger equation, the binding energy is
$$ϵ=d^3r\psi ^{}\left(\frac{^2}{m}+V\right)\psi ,$$
(3)
where $`\psi `$ is the deuteron wave function. Thus
$`{\displaystyle \frac{dϵ}{d\mu ^2}}`$ $`=`$ $`{\displaystyle d^3r\left[\psi ^{}\left(\frac{\sigma _N}{\mu ^2}\frac{^2}{m^2}+\frac{dV}{d\mu ^2}\right)\psi +\frac{d\psi ^{}}{d\mu ^2}\left(\frac{^2}{m}+V\right)\psi +\psi \left(\frac{^2}{m}+V\right)\frac{d\psi ^{}}{d\mu ^2}\right]}`$ (4)
$`=`$ $`{\displaystyle d^3𝒓\left[\psi ^{}\left(\frac{\sigma _N}{\mu ^2}\frac{^2}{m^2}+\frac{dV}{d\mu ^2}\right)\psi ϵ\frac{d}{d\mu ^2}\left(\psi ^{}\psi \right)\right]}.`$ (5)
The term proportional to $`ϵ`$ in this result does not contribute when the deuteron wave function is kept properly normalized and we write
$$\sigma _ϵ=d^3r\psi ^{}\left(\sigma _N\frac{^2}{m^2}+\mu ^2\frac{dV}{d\mu ^2}\right)\psi .$$
(6)
The first term on the r.h.s. of this equation is the effect of the scalar nucleon number and reduces the sigma commutator by a factor $`(1T/m)`$, where $`T`$ is the nucleon kinetic energy, as compared to the additive assumption. Using the equation of motion, we have
$$\sigma _ϵ=d^3r\psi ^{}\left[\frac{\sigma _N}{m}(V+ϵ)+\mu ^2\frac{dV}{d\mu ^2}\right]\psi .$$
(7)
The contribution proportional to $`ϵ`$ is tiny and will not be considered in the sequence. The deuteron is heavily dominated by the one pion exchange potential ($`V_\pi `$) and we write the full NN interaction as
$$V=\overline{V}_\pi +W,$$
(8)
where $`\overline{V}_\pi `$ is the OPEP regularized at small distances and $`W`$ represents other short and medium range effects, associated with either meson or quark dynamics. In the absence of a theory for the influence of chiral symmetry breaking over both $`W`$ and the regularizing potential, we assume that these functions do not depend explicitly on $`\mu `$.
For the deuteron channel one has $`𝝉^{(1)}𝝉^{(2)}=3`$ and the OPEP reads
$$V_\pi =\left(\frac{g_A}{f_\pi }\right)^2\frac{\mu ^3}{16\pi }\left[𝝈^{(1)}𝝈^{(2)}\left(U_CG\right)+S_{12}U_T\right],$$
(9)
where
$`U_C`$ $`=`$ $`{\displaystyle \frac{e^{\mu r}}{\mu r}},`$ (10)
$`U_T`$ $`=`$ $`\left(1+{\displaystyle \frac{3}{\mu r}}+{\displaystyle \frac{3}{\mu ^2r^2}}\right){\displaystyle \frac{e^{\mu r}}{\mu r}}`$ (11)
and G is proportional to a delta-function: $`G=4\pi /\mu ^3\delta ^3(r)`$. The effects of this last term are cancelled by the regularization procedure and we skip them in the sequence.
The derivative of $`V_\pi `$ with respect to $`\mu ^2`$ is
$`{\displaystyle \frac{dV_\pi }{d\mu ^2}}`$ $`=`$ $`2{\displaystyle \frac{f_\pi }{g_A}}\left({\displaystyle \frac{d}{d\mu ^2}}{\displaystyle \frac{g_A}{f_\pi }}\right)V_\pi +{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{g_A}{f_\pi }}\right)^2{\displaystyle \frac{\mu }{16\pi }}\left[𝝈^{(1)}𝝈^{(2)}\left(1{\displaystyle \frac{2}{\mu r}}\right)+S_{12}\left(1+{\displaystyle \frac{1}{\mu r}}\right)\right]e^{\mu r}`$ (12)
$``$ $`2{\displaystyle \frac{f_\pi }{g_A}}\left({\displaystyle \frac{d}{d\mu ^2}}{\displaystyle \frac{g_A}{f_\pi }}\right)V_\pi +\left({\displaystyle \frac{dV_\pi }{d\mu ^2}}\right)_{\frac{g_A}{f_\pi }}.`$ (13)
This allows eq.(7) to be written as
$$\sigma _ϵ=\mu ^2\frac{d\overline{V}_\pi }{d\mu ^2}_{\frac{g_A}{f_\pi }}+c\overline{V}_\pi .$$
(14)
with
$`\overline{V}_\pi {\displaystyle d^3r\psi ^{}\overline{V}_\pi \psi },`$ (15)
$`\mu ^2{\displaystyle \frac{d\overline{V}_\pi }{d\mu ^2}}_{\frac{g_A}{f_\pi }}{\displaystyle d^3r\psi ^{}\mu ^2\left(\frac{d\overline{V}_\pi }{d\mu ^2}\right)_{\frac{g_A}{f_\pi }}\psi },`$ (16)
and
$$c=\frac{\sigma _N}{m}+2\mu ^2\left(\frac{1}{g_A}\frac{dg_A}{d\mu ^2}\frac{1}{f_\pi }\frac{df_\pi }{d\mu ^2}\right).$$
(17)
The quantity $`\sigma _ϵ`$ represents the part of the deuteron $`\sigma `$-term due to NN intraction and may be probed by scalar sources. In practice, these sources may be associated with either photons or pions, as we discuss in the next sections. In order to interpret eq.(14), one notes that the coefficient $`c`$, given by eq.(17), vanishes in the chiral limit: $`\mu ^2=0c=0`$. Hence, at tree level, only the first term contributes, which represents the interaction of the scalar source with the pion exchanged between the nucleons. The coefficient $`c`$, on the other hand, receives contributions from the kinetic energy term and from the derivative of the $`\pi `$NN coupling constant. The latter, as we show in the sequence, corresponds to the interaction of the scalar source with the pion cloud that dresses the $`\pi `$N vextex.
In order to estimate the derivative of $`f_\pi `$, we use the result produced by Gasser and Leutwyler and write:
$`{\displaystyle \frac{df_\pi }{d\mu ^2}}`$ $`=`$ $`{\displaystyle \frac{d}{d\mu ^2}}\left\{F\left[1+{\displaystyle \frac{\mu ^2}{F^2}}\left(\mathrm{}_4^r(\lambda ){\displaystyle \frac{1}{16\pi ^2}}\mathrm{ln}{\displaystyle \frac{\mu ^2}{\lambda ^2}}\right)\right]\right\}`$ (18)
$`=`$ $`{\displaystyle \frac{1}{F}}\left[\mathrm{}_4^r(\lambda ){\displaystyle \frac{1}{16\pi ^2}}\left(1+\mathrm{ln}{\displaystyle \frac{\mu ^2}{\lambda ^2}}\right)\right],`$ (19)
where $`F`$ is the value of $`f_\pi `$ for $`\mu =0`$, $`\mathrm{}_4^r(\lambda )`$ is a renormalization constant and $`\lambda `$ is the renormalization scale. As far as the derivative of $`g_A`$ is concerned, we use the expression derived by Mojžiš and by Fearing, Lewis, Mobed and Scherer and have
$`{\displaystyle \frac{dg_A}{d\mu ^2}}`$ $`=`$ $`{\displaystyle \frac{d}{d\mu ^2}}\left\{G_A\left[1+{\displaystyle \frac{4\mu ^2}{m^2}}a_3{\displaystyle \frac{\mu ^2G_A^2}{16\pi ^2F^2}}{\displaystyle \frac{\mu ^2}{16\pi ^2F^2}}\left(1+2G_A^2\right)\mathrm{ln}{\displaystyle \frac{\mu ^2}{\lambda ^2}}\right]+{\displaystyle \frac{\mu ^2}{4\pi ^2F^2}}b_{17}^r(\lambda )\right\}`$ (20)
$`=`$ $`G_A\left[{\displaystyle \frac{4a_3}{m^2}}{\displaystyle \frac{G_A^2}{16\pi ^2F^2}}{\displaystyle \frac{1}{16\pi ^2F^2}}\left(1+2G_A^2\right)\left(1+\mathrm{ln}{\displaystyle \frac{\mu ^2}{\lambda ^2}}\right)\right]+{\displaystyle \frac{1}{4\pi ^2F^2}}b_{17}^r(\lambda ),`$ (21)
where $`G_A`$ is the value of $`g_A`$ in the limit $`\mu 0`$ and $`b_{17}^r(\lambda )`$ is a constant. Note that the expression adopted for $`g_A`$, within curly brackets, is slightly different from that obtained earlier by Bernard, Kaiser and Meissner and consistent with that produced by Gasser, Sainio and Švarc.
For future purposes, we write down the following results
$`V_\pi `$ $`=`$ $`\left({\displaystyle \frac{g_A}{f_\pi }}\right)^2{\displaystyle \frac{\mu ^3}{16\pi }}{\displaystyle }dr[u^2`$ (22)
$`+`$ $`2\sqrt{8}(1+{\displaystyle \frac{3}{\mu r}}+{\displaystyle \frac{3}{\mu ^2r^2}})uw(1+{\displaystyle \frac{6}{\mu r}}+{\displaystyle \frac{6}{\mu ^2r^2}})w^2]{\displaystyle \frac{e^{\mu r}}{\mu r}},`$ (23)
$`\mu ^2{\displaystyle \frac{dV_\pi }{d\mu ^2}}_{\frac{g_A}{f_\pi }}`$ $`=`$ $`\left({\displaystyle \frac{g_A}{f_\pi }}\right)^2{\displaystyle \frac{\mu ^3}{32\pi }}{\displaystyle }dr[(\mu r2)u^2`$ (24)
$`+`$ $`2\sqrt{8}(\mu r+1)uw(\mu r+4)w^2]{\displaystyle \frac{e^{\mu r}}{\mu r}},`$ (25)
where $`u`$ and $`w`$ are the standard S and D components of the deuteron wave function. These expressions contain negative powers of r, but this does not pose problems for the integration, even in the case of unregularized potentials, since $`u`$ and $`w`$ vanish at the origin. The numerical implications of the results presented here will be explored in sect. V. We now discuss some possible ways of probing the many-body effects of the condensate.
## III electromagnetic probes
A probe which couples locally to the pion field $`\mathit{\varphi }`$ is sensitive to the quantity $`A|\mathit{\varphi }^2|A`$, i.e., to the nuclear condensate. In particular, when a nucleus A is probed by electromagnetic interactions, the many body effects of the condensate correspond to meson exchange contributions to the forward Compton amplitude $`F_{mec}^A(0)`$, for soft photons. This relationship was established by Chanfray and Ericson , using the static approximation, but it is more general and its derivation does not require this assumption. Indeed, in their work on the extension of the Bethe-Levinger sum rule, Ericson, Rosa-Clot and Kulagin have shown that $`F_{mec}^A(0)`$ contains a pion exchange term, which is the seagull represented in fig. 1(a) and can be expressed as:
$$F_{mec}^A(0)=\frac{2}{3}e^2𝑑𝒓\left(A|\mathit{\varphi }^2|AAN|\mathit{\varphi }^2|N\right),.$$
(26)
The second term in the r.h.s. of eq.(26) represents the expectation value of $`\mathit{\varphi }^2`$ for an assembly of free nucleons, which has to be subtracted to obtain the exchange piece. On the other hand, the matrix element $`A|\mathit{\varphi }^2|A`$ is related to the quark condensate by $`_{SB}`$, the chiral symmetry breaking term in the Lagrangian for the $`SU(2)`$ sector, as discussed by Chanfray and Ericson . In the case of QCD one has $`_{SB}=\widehat{m}\overline{q}q`$, assuming $`m_u=m_d=\widehat{m}`$. This symmetry breaking term transforms according to the $`(\frac{1}{2},\frac{1}{2})`$ representation of $`SU(2)\times SU(2)`$ and one requires the same to happen with the effective counterpart. In the case of non-linear realizations of the symmetry, this corresponds to the choice
$$_{SB}=\mu ^2f_\pi \sqrt{f_\pi ^2\mathit{\varphi }^2}.$$
(27)
Imposing the equivalence of the fundamental and effective descriptions, we obtain
$`A|_{SB}|A`$ $`=`$ $`\widehat{m}A|\overline{q}q|A`$ (28)
$`=`$ $`\mu ^2f_\pi ^2{\displaystyle \frac{1}{2}}\mu ^2A|\mathit{\varphi }^2|A+\mathrm{}`$ (29)
In the case of the vacuum, it yields the Gell-Mann-Oakes and Renner relation: $`\widehat{m}0|\overline{q}q|0=\mu ^2f_\pi ^2`$. We apply this relation to both nuclei and free nucleons. Using these results in eq.(26), we obtain the following relation between the condensate and the meson exchange Compton amplitude
$$F_A^{exch}(0)=\frac{4}{3}e^2f_\pi ^2𝑑𝒓\left(\frac{A|\overline{q}q|AAN|\overline{q}q|N}{0|\overline{q}q|0}\right),$$
(30)
which is the same result of ref., but now obtained without the use of the static approximation. In the case of the deuteron, the exchange amplitude is related to the $`\sigma _ϵ`$ calculated in the previous section through
$$F_A^{exch}(0)=\frac{4e^2}{3\mu ^2}\sigma _ϵ.$$
(31)
Two comments on formula (30) are in order. The soft photon amplitude on deuteron is given by the Thomson limit: $`F_d(0)=e^2/M`$. The exchange part $`F_A^{exch}(0)`$ is hidden in this term together with other contributions and they all add up to the Thomson value. The second remark concerns the composition of $`\sigma _ϵ`$, built of three terms: the kinetic energy term, the derivative of the $`\pi `$NN coupling constant, and the derivative of the pion propagator. When transposed into the Compton amplitude, the third part gives rise to the usual meson exchange term of fig. 1(a), where two photons interact with an exchanged pion. The derivative of the $`\pi NN`$ coupling attaches the two photons to the $`\pi NN`$ vertex, fig.1(b). As far as the kinetic energy term is concerned, the fact that $`\mathit{\varphi }^2`$ is a scalar object means that its expectaton value involves a $`\overline{\psi }\psi `$ combination of nucleon fields, which displays the same reduction factor $`(1T/m)`$ as the sigma commutator, with respect to the ordinary nucleon density. Similar remarks apply to pion rescattering. Numerical values will be discussed in Sect.V.
## IV pion probes
Pions exchanged between nucleons may also be probed by means of external pions. In this section we consider $`a_{mec}`$, the MEC contribution to the pion-deuteron scattering lenght. The quadri-momenta for pions at rest are $`k=k^{}=(\omega ,\mathrm{𝟎})`$, where $`\omega =\mu `$ or $`0`$ depending on whether the pions are physical or soft. The $`\pi d`$ scattering length is generically given by
$$a\left(\omega \right)=\frac{\mu }{2\pi \left(1+\mu /M\right)}𝑑𝒓\psi ^{}\left(𝒓\right)t(𝒓;\omega )\psi \left(𝒓\right),$$
(32)
where $`t`$ is the part of the amplitude for the process $`\pi NN\pi NN`$ which does not contain two positive energy nucleons propagating forward in time.
When PCAC holds, the sigma commutator is related to the soft pion scattering amplitude. Hence the value of $`\sigma _ϵ`$ is associated with many body effects in the soft pion PCAC amplitude, since $`a_{mec}^{PCAC}(0)\alpha \sigma _ϵ`$ . We confront this relation with the direct evaluation of $`a_{mec}(\mu )`$, the quantity accessible to experiment. The structure of this amplitude was already discussed in ref. and here we are interested in its relationship with $`\sigma _ϵ`$. This question is important because it concerns the possibility of obtaining empirical information about $`\sigma _ϵ`$ from measurements of the $`\pi `$d scattering length.
For soft pions, the operator $`t_{mec}`$ is completely dominated by processes involving only pions and nucleons, whereas for physical pions there are other contributions, mainly due to $`\mathrm{\Delta }`$ excitations.
In the $`\pi `$N sector, the basic interactions are obtained from the following non-linear Lagrangian, approximately invariant under SU(2)$`\times `$SU(2)
$`_{\pi N}^{int}`$ $`=`$ $`{\displaystyle \frac{1}{8f_\pi ^2}}\left[^\mu \mathit{\varphi }^2_\mu \mathit{\varphi }^2\mu ^2\mathit{\varphi }^4\right]+{\displaystyle \frac{g_A}{2f_\pi }}\overline{N}\gamma ^\mu \gamma _5𝝉N_\mu \mathit{\varphi }`$ (33)
$``$ $`{\displaystyle \frac{1}{4f_\pi ^2}}\overline{N}\gamma ^\mu 𝝉N\mathit{\varphi }\times _\mu \mathit{\varphi }+{\displaystyle \frac{g_A}{8f_\pi ^3}}\overline{N}\gamma ^\mu \gamma _5𝝉N\mathit{\varphi }_\mu \mathit{\varphi }^2+\mathrm{},`$ (34)
designed to be used in the tree approximation.
The meson exchange currents are given by the diagrams shown in fig.2, which contain pion propagators coupled to nucleons. Hence it is useful to parametrize the non relativistic MEC contribution to $`t`$ in the nucleon sector as
$$t_{mec}^N(𝒒;\omega )=\frac{1}{2\mu }\left(\frac{g_A}{2f_\pi }\right)^2\left\{\left[\alpha _n(\omega )\right]\frac{𝝈^{(1)}𝒒𝝈^{(2)}𝒒}{(𝒒^2+\mu ^2)}+\alpha _1^{}(\omega )\mu ^2\frac{𝝈^{(1)}𝒒𝝈^{(2)}𝒒}{(𝒒^2+\mu ^2)^2}\right\},$$
(35)
where $`𝒒`$ is the momentum exchanged between the nucleons and the coefficients $`\alpha _n`$ are determined dynamically, from the graphs of fig.2.
The evaluation of the diagrams 1-10 of fig.2 in the non relativistic tree approximation yields
$`\alpha _1={\displaystyle \frac{2}{f_\pi ^2}},`$ (36)
$`\alpha _1^{}={\displaystyle \frac{1}{f_\pi ^2}}\left(32{\displaystyle \frac{\omega ^2}{\mu ^2}}\right),`$ (37)
$`\alpha _2={\displaystyle \frac{2}{f_\pi ^2}},`$ (38)
$`\alpha _3+\alpha _4+\alpha _5+\alpha _6={\displaystyle \frac{2}{f_\pi ^2}}{\displaystyle \frac{\omega ^2}{m^2}},`$ (39)
$`\alpha _7+\alpha _8=\left({\displaystyle \frac{g_A}{2f_\pi }}\right)^2{\displaystyle \frac{\omega ^2}{m^2}},`$ (40)
$`\alpha _9+\alpha _{10}=0.`$ (41)
As discussed in ref., there is a cancellation between $`\alpha _1`$ and $`\alpha _2`$, required by chiral symmetry. The results for $`\alpha _3+\alpha _4+\alpha _5+\alpha _6`$ and $`\alpha _7+\alpha _8`$ disagree with those of ref. by factors ($`1`$) and (-$`\frac{3}{2}`$) respectively, due to algebraic mistakes in that work, but this has little influence over numerical results.
The MEC amplitude in configuration space is
$$t_{mec}^N(𝒓;\omega )=\frac{1}{2\mu }\frac{1}{3}\left\{\left[\alpha _n(\omega )\right]V_\pi (r)\alpha _1^{}(\omega )\mu ^2\left(\frac{dV_\pi (r)}{d\mu ^2}\right)_{\frac{g_A}{f_\pi }}\right\}.$$
(42)
We now consider the contributions of the $`\mathrm{\Delta }`$ and $`\sigma _N`$ to $`t_{mec}`$. The former were studied in ref. and its efect can be incorporated into eq.(35) by means of the global coefficient $`\alpha _\mathrm{\Delta }\omega ^2/\mu ^2`$, with $`\alpha _\mathrm{\Delta }=0.429\mu ^2`$. The contribution of the $`\pi `$N sigma-term is given by diagrams 1-4 of fig.3 and can be calculated by noting that it enters only in the isospin symmetric $`\pi `$N amplitude $`A^+`$. The corresponding part of this amplitude is denoted by $`A_\sigma ^+`$ and can be parametrized as
$$A_\sigma ^+(t;k^2,k^2)=\frac{\sigma _N}{\mu ^2f_\pi ^2}\left[k^2+k^2\mu ^2+\beta \left(tk^2k^2\right)\right]$$
(43)
and the value of $`\beta `$ can be extracted from scattering data. The evaluation of the diagrams of fig.3 yields, for the coefficients $`\alpha `$,
$$\alpha _{\sigma 1}+\alpha _{\sigma 2}+\alpha _{\sigma 3}+\alpha _{\sigma 4}=\frac{4}{m}\frac{\sigma _N}{f_\pi ^2\mu ^2}\left[\omega ^2\left(𝒒^2+\mu ^2\right)\right].$$
(44)
The term proportional to $`\left(𝒒^2+\mu ^2\right)`$ cancels the pion propagator, giving rise to a contact interaction, which does not contribute when the OPEP is regularized. The overall MEC contribution to the scattering length then becomes
$`a_{mec}(\omega )`$ $`=`$ $`{\displaystyle \frac{1}{4\pi (1+\mu /M)}}{\displaystyle \frac{1}{3f_\pi ^2}}\{[(2+{\displaystyle \frac{g_A^2}{4}}){\displaystyle \frac{\omega ^2}{m^2}}+f_\pi ^2\alpha _\mathrm{\Delta }{\displaystyle \frac{\omega ^2}{\mu ^2}}+{\displaystyle \frac{4\sigma _N}{m}}{\displaystyle \frac{\omega ^2}{\mu ^2}}]V_\pi `$ (45)
$``$ $`(32{\displaystyle \frac{\omega ^2}{\mu ^2}})\mu ^2{\displaystyle \frac{dV_\pi }{d\mu ^2}}_{\frac{g_A}{f_\pi }}\}`$ (46)
In order to establish the relationship between $`a_{mec}(\omega )`$ and $`\sigma _ϵ`$, we use eq.(14) and write
$`a_{mec}(\omega )`$ $`=`$ $`{\displaystyle \frac{1}{4\pi (1+\mu /M)}}{\displaystyle \frac{1}{3f_\pi ^2}}\{[(2+{\displaystyle \frac{g_A^2}{4}}){\displaystyle \frac{\omega ^2}{m^2}}+f_\pi ^2\alpha _\mathrm{\Delta }{\displaystyle \frac{\omega ^2}{\mu ^2}}+{\displaystyle \frac{4\sigma _N}{m}}{\displaystyle \frac{\omega ^2}{\mu ^2}}`$ (47)
$`+`$ $`c(32{\displaystyle \frac{\omega ^2}{\mu ^2}})]V_\pi (32{\displaystyle \frac{\omega ^2}{\mu ^2}})\sigma _ϵ\}.`$ (48)
In the soft pion limit $`(\omega 0)`$ this result becomes
$$a_{mec}(0)=\frac{1}{4\pi (1+\mu /M)}\frac{\sigma _ϵcV_\pi }{f_\pi ^2}.$$
(49)
For physical pions, on the other hand, one has
$$a_{mec}(\mu )=\frac{1}{4\pi (1+\mu /M)}\frac{1}{3f_\pi ^2}\left\{\sigma _ϵ\left[\left(2+\frac{g_A^2}{4}\right)\frac{\mu ^2}{m^2}+f_\pi ^2\alpha _\mathrm{\Delta }+\frac{4\sigma _N}{m}+c\right]V_\pi \right\}.$$
(50)
The first observation from eq.(49) is that $`a_{mec}(0)`$ is not just proportional to $`\sigma _ϵ`$, as in the PCAC result, $`a_{mec}^{PCAC}(0)`$, but the term $`cV_\pi `$ which appears in the epression (14) of $`\sigma _ϵ`$ is cancelled in $`a_{mec}(0)`$. The reason for this difference is that the usual meson exchange amplitude, $`a_{mec}`$, does not incorporate terms where the two pions are attached to the $`\pi NN`$ vertex through loop diagrams. These terms are instead present in the PCAC expression. The fact that the term in $`cV_\pi `$ may give a large contribution to $`\sigma _ϵ`$ indicates a possible importance also as an exchange correction. Moreover, inspecting eqs.(49) and (50), one notes that the contribution proportional to $`dV_\pi /d\mu ^2`$ is three times larger for soft pions than for physical pions, due to the strong energy dependence of the intermediate $`\pi \pi `$ amplitude of diagram 1. This feature is consistent with the results found by Chanfray, Ericson and Wambach , who studied the self energy $`\mathrm{\Pi }(\omega ,𝒌)`$ of a pion propagating in a gas of of pions. Using PCAC and the Hartree approximation, they found that
$$\mathrm{\Pi }(\omega ,𝒌)=\frac{\rho _s}{f_\pi ^2}\left[\mu ^2\frac{2}{3}\left(\omega ^2𝒌^2\right)\right],$$
(51)
where $`\rho _s`$ is the scalar density of the pions. Thus, for soft and physical pions one has, respectively, $`\mathrm{\Pi }(0,0)=\rho _s\mu ^2/f_\pi ^2`$ and $`\mathrm{\Pi }(\mu ,0)=\rho _s\mu ^2/3f_\pi ^2`$. As this self energy is related to the MEC amplitude, both must change in the same proportion when one goes from physical to soft pions.
In summary, the measurable meson exchange contribution written in eq.(50) has little relation to the quark condensate. Therefore, the pion-deuteron scattering length provides no exploitable information about this condensate. In the next section we discuss numerically the results produced here.
## V results and conclusions
We estimate the numerical implications of the results produced in the previous sections and adopt the following values for the various constants: M=1875.61 MeV, m=938.28 MeV , $`\mu `$=139.57 MeV , g<sub>A</sub>=1.26 , f<sub>π</sub>=93.3 MeV , $`\sigma _N`$=45 MeV , $`\alpha _\mathrm{\Delta }`$=-0.43 $`\mu ^2`$ , $`\lambda =\mu `$, $`\mathrm{}_4^r(\mu )=4.3/16\pi ^2`$, and $`a_3=m\sigma _N/4\mu ^2`$. As very little is known about the constant $`b_{17}^r(\mu )`$, we neglect it in eq.(21). With these inputs, we find a negative value for $`c`$: -0.30, which is strongly dominated by the derivative of the $`\pi `$N coupling constant and has opposite sign to the kinetic energy term. Thus one has
$`\left[\left(2+g_A^2/4\right)\mu ^2/m^2+f_\pi ^2\alpha _\mathrm{\Delta }+4\sigma _N/m+c\right]=\left[0.050.19+0.190.30\right]=0.25`$ .
Expressions (23) and (25) are based on the assumption that the short range components of the interaction are not important since the OPEP strongly dominates the deuteron. In order to test this hypothesis, we consider the case of a toy potential containing an OPEP tail and regularized by means of monopole form factors . It has the same form as eq.(9), with $`U_C`$, $`G`$ and $`U_T`$ given by
$`U_C`$ $`=`$ $`{\displaystyle \frac{e^{\mu r}}{\mu r}}{\displaystyle \frac{\mathrm{\Lambda }_C}{\mu }}{\displaystyle \frac{e^{\mathrm{\Lambda }_Cr}}{\mathrm{\Lambda }_Cr}}{\displaystyle \frac{1}{2}}{\displaystyle \frac{\mu }{\mathrm{\Lambda }_C}}\left({\displaystyle \frac{\mathrm{\Lambda }_C^2}{\mu ^2}}1\right)e^{\mathrm{\Lambda }_Cr},`$ (52)
$`G`$ $`=`$ $`\delta {\displaystyle \frac{1}{2}}{\displaystyle \frac{\mu }{\mathrm{\Lambda }_C}}\left({\displaystyle \frac{\mathrm{\Lambda }_C^2}{\mu ^2}}1\right)^2e^{\mathrm{\Lambda }_Cr},`$ (53)
$`U_T`$ $`=`$ $`\left(1+{\displaystyle \frac{3}{\mu r}}+{\displaystyle \frac{3}{\mu ^2r^2}}\right){\displaystyle \frac{e^{\mu r}}{\mu r}}{\displaystyle \frac{\mathrm{\Lambda }_T^3}{\mu ^3}}\left(1+{\displaystyle \frac{3}{\mathrm{\Lambda }_Tr}}+{\displaystyle \frac{3}{\mathrm{\Lambda }_T^2r^2}}\right){\displaystyle \frac{e^{\mathrm{\Lambda }_Tr}}{\mathrm{\Lambda }_Tr}}`$ (54)
$``$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{\mathrm{\Lambda }_T}{\mu }}\left({\displaystyle \frac{\mathrm{\Lambda }_T^2}{\mu ^2}}1\right)\left(1+\mathrm{\Lambda }_Tr\right){\displaystyle \frac{e^{\mathrm{\Lambda }_Tr}}{\mathrm{\Lambda }_Tr}},`$ (55)
where $`\mathrm{\Lambda }_C`$ and $`\mathrm{\Lambda }_T`$ are cut-offs for the central and tensor components and the parameter $`\delta `$ regulates the strength of the short range function $`G`$. The pure OPEP results are recoverd in the limit $`\mathrm{\Lambda }_C=\mathrm{\Lambda }_T\mathrm{}`$ and $`\delta =1`$. It yields a regularized version of eqs.(23) and (25), namely
$`V_\pi =\left({\displaystyle \frac{g_A}{f_\pi }}\right)^2{\displaystyle \frac{\mu ^3}{16\pi }}{\displaystyle 𝑑r\left[\left(U_CG\right)u^2+2\sqrt{8}U_Tuw+\left(U_CG2U_T\right)w^2\right]},`$ (56)
$`\mu ^2{\displaystyle \frac{dV_\pi }{d\mu ^2}}_{\frac{g_A}{f_\pi }}={\displaystyle \frac{3}{2}}V_\pi `$ (57)
$`\left({\displaystyle \frac{g_A}{f_\pi }}\right)^2{\displaystyle \frac{\mu ^3}{16\pi }}{\displaystyle 𝑑r\mu ^2\left[\frac{d\left(U_CG\right)}{d\mu ^2}u^2+2\sqrt{8}\frac{dU_T}{d\mu ^2}uw+\frac{d\left(U_CG2U_T\right)}{d\mu ^2}w^2\right]},`$ (58)
with
$`\mu ^2{\displaystyle \frac{dU_C}{d\mu ^2}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[\left(1+{\displaystyle \frac{1}{\mu r}}\right)e^{\mu r}{\displaystyle \frac{e^{\mathrm{\Lambda }_Cr}}{\mu r}}{\displaystyle \frac{1}{2}}{\displaystyle \frac{\mu }{\mathrm{\Lambda }_C}}\left({\displaystyle \frac{\mathrm{\Lambda }_C^2}{\mu ^2}}+1\right)e^{\mathrm{\Lambda }_Cr}\right],`$ (59)
$`\mu ^2{\displaystyle \frac{dG}{d\mu ^2}}`$ $`=`$ $`\delta {\displaystyle \frac{1}{4}}{\displaystyle \frac{\mu }{\mathrm{\Lambda }_C}}\left(3{\displaystyle \frac{\mathrm{\Lambda }_C^4}{\mu ^4}}2{\displaystyle \frac{\mathrm{\Lambda }_C^2}{\mu ^2}}1\right)e^{\mathrm{\Lambda }_Cr},`$ (60)
$`\mu ^2{\displaystyle \frac{dU_T}{d\mu ^2}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}[(1+{\displaystyle \frac{4}{\mu r}}+{\displaystyle \frac{9}{\mu ^2r^2}}+{\displaystyle \frac{9}{\mu ^3r^3}})e^{\mu r}3{\displaystyle \frac{\mathrm{\Lambda }_T^3}{\mu ^3}}(1+{\displaystyle \frac{3}{\mathrm{\Lambda }_Tr}}+{\displaystyle \frac{3}{\mathrm{\Lambda }_T^2r^2}}){\displaystyle \frac{e^{\mathrm{\Lambda }_Tr}}{\mathrm{\Lambda }_Tr}}`$ (61)
$``$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{\mathrm{\Lambda }_T}{\mu }}(3{\displaystyle \frac{\mathrm{\Lambda }_T^2}{\mu ^2}}1)(1+\mathrm{\Lambda }_Tr){\displaystyle \frac{e^{\mathrm{\Lambda }_Tr}}{\mathrm{\Lambda }_Tr}}],`$ (62)
In general, the deuteron binding energy is a function of the form $`ϵ(g_A,f_\pi ,\mu ,\mathrm{\Lambda }_C,\delta ,\mathrm{\Lambda }_T)`$. As $`g_A`$, $`f_\pi `$ and $`\mu `$ are kept fixed, the binding energy depends on the the short range parameters $`\mathrm{\Lambda }_C`$, $`\delta `$ and $`\mathrm{\Lambda }_T`$ . When constructing the deuteron, we fix two of them and look for the third one so as to have $`ϵ=2.2250`$ MeV.
In table 1 we display our results for $`V_\pi `$ and $`\mu ^2dV_\pi /d\mu ^2`$ as given by the the perturbative OPEP (pert) eqs. (56) and (58) and by the regularized OPEP (toy), eqs.(23) mand (25). The first feature to be noted is that the sensitivity to the regularization of the potential is much greater for $`V_\pi `$ than for $`\mu ^2dV_\pi /d\mu ^2`$, due to the fact that the latter is less influenced by the short distance components of the wave function. In the case of the calculation based on the regularized potential, the large variations of the inner parameters considerd change results only by a few percent. This suggests that the self consistency between the potential and the wave function is important. In table 2 we present our results for the case of the Argonne $`v_{14}`$ and super soft core C potentials and the values quoted also follow the pattern found in the case of the toy potential.
Inspection of these tables shows that the expectation values of the potential are about ten times larger than those of its derivative. Taking this information into eq.(35), one finds that this corresponds to an average pion momentum $`q=3\mu `$, which is relatively high. The disturbance of the QCD vacuum due to the NN interaction, represented by $`\sigma _ϵ`$, has a central value of about 10 MeV, which is about five times the binding energy and corresponds to about 10% of the total deuteron $`\sigma `$ term. Our results have the same magnitude but an opposite sign to that produced by Gammal and Frederico in the framework of the Skyrme model. The values of $`\sigma _ϵ`$ quoted in the tables are dominated by the component involving the constant $`c`$ in eq.(14). This in turn depends strongly on $`dg_A/d\mu ^2`$ which was calculated using chiral perturbation theory and contains an unknown constant. Hence our result has to be taken as an estimate of the magnitude of $`\sigma _ϵ`$.
The columns $`a_{mec}(0)`$, eq.(49) and $`a_{mec}(\mu )`$, eq.(50), correspond respectively to the quantities that have a relation to the condensate $`\sigma _ϵ`$. The difference between $`a_{mec}(0)`$ and $`a_{mec}(\mu )`$ stems in part from the factor 3, related to the off-shell behaviour of the intermediate pion-pion scattering amplitude, as discussed at the end of section IV. In the case of soft pions, it is worth noting that $`\frac{1}{3}a_{mec}(0)0.0007\mu ^1`$, in agreement with the value found by Robilotta and Wilkin for physical pions . The value for $`F_A^{exch}(0)`$, the many body electromagnetic term of the commutator amplitude, is also displayed.
In summary, we have studied the many body effects of the quark condensate in the deuteron through the Feynman-Hellmann theorem and found out that the part of the deuteron sigma commutator associated with the NN interaction is smaller than the pion-nucleon sigma-term, but five times larger than the binding energy. With the restricions mentioned previously ($`b_{17}^r`$ is not known), we find that $`\sigma _ϵ`$ could be dominated by the derivative of the $`\pi `$N coupling constant. We have also linked the changes in the condensate with meson exchange effects for probes that can couple to the pion field, namely Compton and pion scatterings. As far as the possibility of extracting $`\sigma _ϵ`$ from the pion-deuteron scattering length, our study has shown that meson exchange effects are comparable to the present experimental error . However the extrapolation to the soft limit produces important changes which tend to blur the contribution of $`\sigma _ϵ`$. The reason why the pion-deuteron scattering length is unexploitable is that the part of the exchange correction which is linked to the sigma commutator is reduced by a factor 3 when one goes from soft to physical pions, which makes it small. Moreover, in the last case, non static corrections appear, in such a way that the extraction of the interesting term becomes unfeasible. In the case of the Compton amplitude, instead, no such problem arises, since soft photons are directly accessible to experiment, opening the the possibility of empirical determination. The photons are by far a superior tool as a source of information on the quark condensate, not only in the deuteron, but also in nuclei.
Acknowledgments
We would like to thank G. Chanfray, J. Delorme, C.A. Dominguez and H. Leutwyler for useful discussions, J. Gasser, J. Goity and M. Mojžiš for exchanges of messages, and U-G. Meissner for help in dealing with aspects chiral perturbation theory. It is also our pleasure to acknowledge the hospitality of the Institute of Nuclear Theory and the Nuclear Theory Group of the University of Washington, USA, where this work was initiated. M.R.R. would also like to thank the hospitality of the Division de Physique Theorique de l’Institut de Physique Nucleaire, Orsay, France and FAPESP, for financial support.
Figure Captions
Fig.1 Seagull meson exchange diagram contributing to the Compton amplitude.
Fig.2. Diagrams contributing to the pion-deuteron scattering length in the pure pion-nucleon sector; the crosses in the propagators of figs. 9 and 10 indicate that they refer to antinucleons.
Fig.3. Diagrams contributing to the pion-deuteron scattering length due to the isospin-symmetric amplitude represented by the black square.
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# Mode Solutions of the Klein-Gordon equation in warped spacetimes
## 1 Introduction
A warped spacetime , , is a product manifold $`V_1\times V_2`$ endowed with (omitting the canonical projections) a metric <sup>1</sup><sup>1</sup>1The minus sign is dictated by our signature $`+\mathrm{}`$. See next Section.
$$(g)=\alpha (S)\gamma $$
(1)
where $`\alpha `$ and $`\gamma `$ respectively are metric tensors on $`V_1`$ and $`V_2`$. Either $`\alpha `$ or $`\gamma `$ is Lorentzian, so that $`(g)`$ is normal hyperbolic. $`S`$ is a positive function on $`V_1`$, it is convenient to set $`S=\mathrm{e}^{2i\mathrm{\Theta }}`$.
The most simple example of warped spacetime is probably given by the Friedman-Robinson-Walker (FRW) universe. But the warped structure accomodates a large number of metrics of physical interest in General Relativity.
In FRW, the wave equation has been studied in details for many decades, in relation with the old prblem of defining its positive-frequency solutions. To this end, the general solution is usually developed over the so-called mode solutions, so the Klein-Gordon equation gets reduced to an ordinary differential equation involving the time variable only.
In fact the property of being spatially homogeneous, exhibited by FRW universes, plays no role in this reduction, which permits to follow the same line in generalized FRW spacetimes where the spatial sections have not necessarily a constant curvature. The geometrical properties of these spacetimes have been systematically investigated by M. Sánchez .
Warped structures in general have been widely studied in the literature, the main issue being of course to establish relations between symmetry and curvature properties of the total spacetime and that of its factor manifolds.
In this article we are concerned with interesting features exhibited by the Klein Gordon equation in a warped spacetime of arbitrary dimension and type (with some preference however for the case where $`V_1`$ is Lorentzian). Assuming the minimal coupling we write
$$(^2+m^2)\mathrm{\Psi }=0$$
(2)
for a complex $`c`$-number-valued wave function.
In spite of well-known limitations, a particle interpretation of (2) remains of interest. In the operator approach to quantum field theory in curved spacetime, one-particle concepts play at least the role of useful tools. For instance, the kernel which projects any solution onto a positive-frequency subspace and determines a definition of the vacuum, is a solution of the KG equation (often referred to as the two-point Wightman function).
Our goal is to show that most results obtained up to now in the context of generalized FRW spacetimes can be systematically extended to any kind of warped spacetime.
We shall take advantage of the warped-product structure in order to carry out a systematic reduction of (2), and shall end up with an equation to be solved for a reduced wave function which depends only on the coordinates running in the first factor manifold $`V_1`$.
The principle of this procedure consists in developing the general solution of equation (2) over special ones that generalize the mode solutions arising in the customary FRW framework.
The possibility of reducing the KG equation with help of generalized modes can be tracked back to a remarkable feature of classical motion in warped spacetimes. That point is briefly exposed in Section 2, in parallel to a study of the differential operator which describes the quantum motion. Generalized mode solutions are introduced in Section 3, where the separation of variables is carried out.
Together with the KG equation, we analyse the sesquilinear form defined on its solutions. Our motivation is the fact that, for complex (resp. real) solutions sesquilinear (resp. bilinear) forms are fundamental with respect to quantum mechanics; they provide a framework for the construction of a ”complex structure positive operator” ensuring the splitting of any solution into positive and negative-frequency parts , .
We shall demonstrate in Sections 4-5 that the sesquilinear form defined for solutions to equation (2) (by conservation of the Gordon current) undergoes some kind of reduction. Indeed, under very general assumptions, the space of complex solutions to the reduced equation can be in turn endowed with a sesquilinear form of its own. Actually, a conserved vector-density defined on the first factor manifold is associated with the reduced wave equation; we display the relationship of this object with the usual Gordon current.
The possibility of extending this study to nonminimal couplings is discussed in Section 6.
Although we have in mind possible applications to quantum mechanics, this article is written from the viewpoint of differential geometry; all functions and tensors are supposed to be smooth, that is $`C^{\mathrm{}}`$.
Moreover all manifolds considered here are implicitly assumed to be connected.
### 1.1 Notation
In the product $`(V)=V_1\times V_2`$, the factors $`V_1,V_2`$ respectively have dimensions $`p,q`$. We define Type I (resp. Type II) by this property that $`(V_1,\alpha )`$ (resp. $`(V_2,\gamma )`$ ) is Lorentzian.
Throughout this paper we use coordinate charts adapted to the warped structure. The metric takes on the orthogonal form
$$ds^2=\alpha _{AB}dx^Adx^BS(x^C)\gamma (x^k)_{ij}dx^idx^j$$
(3)
where $`A,B,CI_1,i,j,kI_2I_1I_2=`$ and $`I_1I_2`$ covers the whole set of integers $`0,1,2,\mathrm{}.,p+q1`$.
\[For instance in 4 dimensions:
$`I_10,I_21,2,3`$ for FRW spacetimes, and we can take $`I_10,3I_21,2`$ for spherically symmetric universes\]
The minus sign in equation (3) is dictated by our signature $`+\mathrm{}`$. Type I necessarily corresponds to $`\gamma `$ positive definite. In contrast, Type II implies having the quadratic form $`\alpha `$ elliptic but negative definite.
$`g_{AB}=\alpha _{AB}`$ and $`g_{ij}=S\gamma _{ij}`$.
\[Example: In FRW spacetimes, , $`p=1,q=3`$ and we have $`A,B,C=0`$ and $`i,j,k=1,2,3`$\].
It is clear that
$$g_{\mu \nu }=\left(\begin{array}{cc}\alpha _{AB}& 0\\ 0& S\gamma _{ij}\end{array}\right),g^{\mu \nu }=\left(\begin{array}{cc}\alpha ^{AB}& 0\\ 0& S^1\gamma ^{ij}\end{array}\right)$$
where $`g^{AB}g_{BC}=\delta _C^A,g^{ij}g_{jk}=\delta _k^i`$ Notice that
$$g^{AB}=\alpha ^{AB}$$
if we define $`\alpha ^{AB}`$ (resp. $`\gamma ^{ij}`$) as the contravariant tensor inverting $`\alpha _{AB}`$ (resp. $`\gamma _{ij}`$), that is $`\alpha ^{AB}\alpha _{BC}=\delta _C^A,\gamma ^{AB}\gamma _{BC}=\delta _C^A`$. We obviously can write $`g_{AB}=\alpha _{AB}`$.
Notice that
$$g^{ij}=(S)^1\gamma ^{ij}$$
Caution that the Types I, II defined above must not be confused with the Classes $`A,B`$ introduced by Carot and da Costa for four-dimensional warped spacetimes . For $`p+q=4`$, intersection of Classes $`A_1,A_2`$ and $`B`$ with both Types give rise to six possibilities.
In references each class is combined with the Types by chosing a sign $`\pm `$ in the generic form of the metric. The contact with the notation of Ref. can be made as follows:
$$\alpha =h_1,\gamma =h_2$$
where $`h_1`$ and $`h_2`$ are the metrics assigned to the factor manifolds in Ref. . Our quadratic form $`ds^2=g_{\alpha \beta }dx^\alpha dx^\beta `$ and their have opposite signs.
Geometric objects corresponding to $`(V_1,\alpha ),(V_2,\gamma )`$ are affected by the index $`1,2`$ respectively. This label will be put on the left for connexions, curvature tensors and their contractions.
Remark
It is a trivial observation that the Minkowski metric is (globally) decomposable in several ways. Similarly, it may happen that a given spacetime can be considered as warped in several ways, so that, for instance, it may be of Class $`A`$ for one structure whilst it is of Class $`B`$ for another one . But here, we consider only one structure at one time, so the six possibilities mentioned above are mutually exclusive.
$`\eta =\sqrt{|g|}\epsilon `$ where $`\epsilon `$ is the Levi-Civitta tensor. Useful determinants are as follows: Setting
$$g=detg_{\alpha \beta }\gamma =det\gamma _{ij}\alpha =det\alpha _{AB}detg_{AB}$$
we have
$$detg_{ij}=(S)^qdet\gamma _{ij}$$
and therefore $`g=\alpha (S)^q\gamma `$ $`g=detg_{AB}(S)^qdet\gamma _{ij}`$ thus $`g=\alpha (S)^q\gamma `$. We shall rather use
$$\sqrt{|g|}=\sqrt{|\alpha |}S^{q/2}\sqrt{|\gamma |}$$
(4)
Volume elements: Setting
$$d_1^px=dx^A,d_2^qx=dx^j,AI_1,jI_2$$
we have $`d^{p+q}x=d_1^pxd_2^qx`$.
The volume elements of $`(V_1,\alpha )`$ and $`(V_2,\gamma )`$ respectively are $`\sqrt{|\alpha |}d_1^px`$ and $`\sqrt{|\gamma |}d_2^qx`$, whereas naturally the volume form of $`(V)`$ is $`\sqrt{|g|}d^{p+q}x`$.
For Type I, it is convenient to take $`x^A`$ running from $`0`$ to $`p1`$ and to factorize out the time coordinate by setting
$$\omega =dx^1\mathrm{}dx^{p1}$$
(5)
so that $`d_1^px=dx^0\omega `$.
With this convention, we have
$$\omega =\frac{1}{(p1)!}\epsilon _{0B_1\mathrm{}B_{p1}}dx^{B_1}\mathrm{}dx^{B_{p1}}$$
(6)
## 2 Classical and Quantum Motion in Warped Spacetime
### 2.1 Geodesic motion
An equation like (2) can be thought of as describing the quantum motion of a test particle, in the approximation where possible particle creation is neglected. It is in order to point out that, in any warped spacetime, the classical motion of a free particle (geodesic motion) already enjoys an interesting property which directly stems from the warping.
Indeed the equations of motion of a test particle in $`V`$ are canonically generated by a ”Hamiltonian function”
$$G(x,p)=\frac{1}{2}g^{\alpha \beta }p_\alpha p_\beta =\frac{1}{2}(\alpha ^{AB}p_Ap_BS^1\gamma ^{ij}p_ip_j)$$
which is a scalar in the cotangent bundle $`T_{}(V)`$. According to the canonical symplectic form of this bundle, we have the usual Poisson brackets $`\{x^\alpha ,p_\beta \}=\delta _\beta ^\alpha `$ , etc. Constants of the motion are characterized by a vanishing Poisson bracket with $`G`$. It is easy to verify that
###### Proposition 1
In any warped spacetime, with the metric written like in (1), geodesic motion admits the first integral
$$2K=\gamma ^{ij}p_ip_j$$
(7)
where $`p_\alpha `$ are the momenta.
\[ Proof: We see that $`\{K,\alpha ^{AB}p_Ap_B\}`$ vanishes. This is obvious since $`K`$ only depends on $`x^C,p_D`$. Then we observe that $`\{p_j,x^A\}=0`$. It follows that $`\{K,S\}=0`$, so finally $`\{g^{\alpha \beta }p_\alpha p_\beta ,K\}=0`$ and $`K`$ is a constant of the motion.\] Indeed we derive with help of the standard Poisson brackets that $`\{G,K\}=0`$.
For Type I, the quantity $`K/m`$ somehow generalizes the kinetic energy, although its conservation is ensured even if $`(V_2,\gamma )`$ fails to admit a group of translations. For instance, when $`V`$ is simply $`𝐑\times 𝐑^3`$ warped with some time-depending scale factor, the conservation of $`K`$ can also be derived from the existence of a translation group in $`𝐑^3`$. And in this case, $`K/m`$ is just the kinetic energy in the usual sense. But the point is that this property survives when $`𝐑^3`$ is replaced by any other three-dimensional manifold.
### 2.2 Quantum motion
We assume minimal coupling, so we write the KG equation as (2) where $`^2\mathrm{\Psi }=g^{\alpha \beta }_\alpha _\beta \mathrm{\Psi }`$. But we shall use the well-known formula
$$^2\mathrm{\Psi }=\frac{1}{\sqrt{|}g|}\mu (\sqrt{|g|}g^{\mu \nu }_\nu \mathrm{\Psi })$$
(8)
By formal analogy with (7) it is natural to consider that the quantum mechanical analog of $`K`$ is $`K_{\mathrm{quant}}=\frac{1}{2}\mathrm{\Delta }_2`$ where $`\mathrm{\Delta }_2`$ is the Laplace-Beltrami operator in $`(V_2,\gamma )`$. Indeed we have in obvious notations $`\mathrm{\Delta }_2=\gamma ^{ij}(_2)_i(_2)_j`$ acting on scalars. But we rather use the formula
$$\mathrm{\Delta }_2\mathrm{\Psi }=\frac{1}{\sqrt{|\gamma |}}_i(\sqrt{|\gamma |}\gamma ^{ij}_j\mathrm{\Psi })$$
Developing formula (8) we get
$$\sqrt{|g|}^2\mathrm{\Psi }=_A(\sqrt{|g|}g^{AB}_B\mathrm{\Psi })+_i(\sqrt{|g|}g^{ij}_j\mathrm{\Psi })$$
but $`g^{AB}=\alpha ^{AB}`$ and $`g^{ij}=S^1\gamma ^{ij}`$ thus
$$\sqrt{|g|}^2\mathrm{\Psi }=_A(\sqrt{|g|}g^{AB}_B\mathrm{\Psi })_i(\sqrt{|g|}S^1\gamma ^{ij}_j\mathrm{\Psi })$$
where $`_kS=0`$. But in view of (4) we have
$$\sqrt{|g|}^2\mathrm{\Psi }=_A(\sqrt{|g|}g^{AB}_B\mathrm{\Psi })S^1S^{q/2}_i(\sqrt{|\alpha |}\sqrt{|\gamma |}\gamma ^{ij}_j\mathrm{\Psi })$$
Since $`\alpha `$ only depends on $`x^A`$ we obtain
$$\sqrt{|g|}^2\mathrm{\Psi }=_A(\sqrt{|g|}g^{AB}_B\mathrm{\Psi })S^1S^{q/2}\sqrt{|\alpha |}_i(\sqrt{|\gamma |}\gamma ^{ij}_j\mathrm{\Psi })$$
Again develop $`\sqrt{|g|}`$ and remember that $`_A\gamma =0`$. We get
$$\sqrt{|g|}^2\mathrm{\Psi }=\sqrt{|\gamma |}_A(\sqrt{|\alpha |}S^{q/2}g^{AB}_B\mathrm{\Psi })S^1S^{q/2}\sqrt{|\alpha |}_i(\sqrt{|\gamma |}\gamma ^{ij}_j\mathrm{\Psi }$$
Again develop $`\sqrt{|g|}`$ hence
$$^2\mathrm{\Psi }=\frac{1}{\sqrt{|\alpha |}S^{q/2}}_A(\sqrt{|\alpha |}S^{q/2}g^{AB}_B\mathrm{\Psi })\frac{S^1}{\sqrt{|\gamma |}}_i(\sqrt{|\gamma |}\gamma ^{ij}_j\mathrm{\Psi })$$
It is convenient to define
$$D\mathrm{\Psi }=\frac{1}{\sqrt{|\alpha |}}S^{q/2}_A(\sqrt{|\alpha |}S^{q/2}g^{AB}_B\mathrm{\Psi })$$
irrespective of whether $`\mathrm{\Psi }`$ is a solution to (2) or not.
Indeed the second order differential operator $`D`$ only affects quantities depending on the $`x^A`$ variables. So we can write
$$^2\mathrm{\Psi }=D\mathrm{\Psi }S^1\frac{1}{\sqrt{|\gamma |}}_i(\sqrt{|\gamma |}\gamma ^{ij}_j\mathrm{\Psi })$$
In other words we have the identity
$$^2\mathrm{\Psi }=D\mathrm{\Psi }S^1\mathrm{\Delta }_2\mathrm{\Psi }$$
(9)
where $`\mathrm{\Delta }_2`$ is the $`q`$-dimensional Laplace-Beltrami operator, associated with the manifold $`(V_2,\gamma )`$. As an operator extended to functions on $`V`$, it does not affect the quantities depending on $`x^A`$ only.
It is clear that $`\mathrm{\Delta }_2`$ commutes with $`D`$, because these operators act on separate sets of variables. As $`S`$ does not depend on the $`x^j`$ coordinates, it is clear that $`\mathrm{\Delta }_2`$ (or equivalently $`K_{\mathrm{quant}}`$) commutes with $`^2`$. In the classical limit ($`\mathrm{}0`$) this property reduces to the conservation of $`K`$.
Let us re-arrange $`D`$ in order to simplify the expression of $`^2\mathrm{\Psi }`$.
Let us provisionally use coordinates where $`|\alpha |=1`$. We obtain
$$D\mathrm{\Psi }=S^{q/2}_A(S^{q/2}g^{AB}_B\mathrm{\Psi })$$
$$D\mathrm{\Psi }=_A(g^{AB}_B\mathrm{\Psi })+\frac{q}{2}(_A\mathrm{log}S)g^{AB}_B\mathrm{\Psi }$$
But $`g^{AB}=\alpha ^{AB}`$ and $`|\alpha |=1`$, thus
$$D\mathrm{\Psi }=\mathrm{\Delta }_1\mathrm{\Psi }+\frac{q}{2}\alpha ^{AB}(_A\mathrm{log}S)_B\mathrm{\Psi }$$
(10)
which is valid in all coordinates and for arbitrary $`\mathrm{\Psi }`$. This expression, where the coordinates $`x^j`$ are ignorable, is to be inserted into equation (9).
## 3 Mode Solutions, Product Solutions
Since $`[\mathrm{\Delta }_2,^2]`$ vanishes, it is clear that some solutions of the KG equation are also eigenstates of $`\mathrm{\Delta }_2`$. If $`\mathrm{\Phi }`$ is such a solution we have
$$\mathrm{\Delta }_2\mathrm{\Phi }=\lambda \mathrm{\Phi }$$
(11)
for some constant number $`\lambda \mathrm{Spec}.(V_2)`$. We shall generalize the terminology which is commonly used when $`V`$ is FRW universe and shall call $`\mathrm{\Phi }`$ a Mode Solution to the KG equation. For a solution in mode $`\lambda `$, the KG equation reduces to
$$(D+\lambda S^1+m^2)\mathrm{\Phi }=0$$
(12)
where the differential operator $`D`$ acting on $`\mathrm{\Phi }`$ affects the $`x^A`$’s only.
Finally, $`x^A`$ and $`x^j`$ are respectively ignorable in equations (12) and (11), which realizes separation of the variables.
Leaving aside the solving of (11), the original KG equation in $`(V)`$ has been reduced to a (linear) partial differential equation in $`p`$ dimensions.
From now on, we look for solutions to (2) in the form of a superposition of various mode solutions corresponding to all possible values taken by $`\lambda `$ in the spectrum of $`V_2`$.
In the special case where $`p=1`$, equation (12) is an ordinary 2nd order equation and its solutions form a two-dimensional vector space.
Otherwize, the space of solutions still has infinitely many dimensions.
Some special solutions of the wave equation (2) have the form of a product of functions compatible with the product structure of spacetime, namely
$$\mathrm{\Phi }=f(x^A)F(x^k)$$
(13)
We observe that
###### Proposition 2
Any product solution $`fF`$ to the KG equation is a mode solution, and $`F`$ is eigenfunction of $`\mathrm{\Delta }_2`$ .
Proof.
$$^2\mathrm{\Phi }=(Df)FS^1f\mathrm{\Delta }_2F$$
$$(^2+m^2)\mathrm{\Phi }=F(D+m^2)fS^1f\mathrm{\Delta }_2F$$
According to KG equation this quantity vanishes.
$$F(D+m^2)f=S^1f\mathrm{\Delta }_2F$$
Discarding a trivial case, neither $`f`$ nor $`F`$ can identically vanish. When $`f`$ and $`F`$ are not zero, we divide by $`fF`$ and multiply by $`S`$. We get
$$S\frac{(D+m^2)f}{f}=\frac{\mathrm{\Delta }_2F}{F}$$
The l.h.s. of this equation depends on $`x^A`$ only while the r.h.s. only depends on $`x^k`$. Both are thus necessarily constants, so there exists some $`\lambda `$ such that
$$\mathrm{\Delta }_2F=\lambda F$$
(14)
Let $`[\lambda ]`$ be the space of smooth functions on $`V_2`$ satisfying (14) for a given value of $`\lambda `$. For $`V_2`$ compact, $`\mathrm{\Delta }_2`$ has a discrete spectrum which is the infinite sequence
$$\mathrm{Spec}(V_2,\gamma )=\{\lambda _0=0,<\lambda _1,\mathrm{}<\lambda _n\mathrm{}\}$$
In this case $`_n`$ denotes $`[\lambda _n]`$ and we know that its dimension $`d(n)`$ is finite \[berg\].
The converse of the previous Proposition is not true, but
###### Proposition 3
In a warped spacetime of Type I with compact $`(V_2)`$, any mode solution corresponding to a given $`\lambda `$ in $`\mathrm{Spec}(V_2,\gamma )`$ is a finite sum of product solutions.
Let $`f_u(x^A)`$ be the coefficients of $`\mathrm{\Phi }`$ in a development over a basis $`F_1(x^j),\mathrm{}\mathrm{}`$.
$$\mathrm{\Phi }=\underset{u=1}{\overset{d(n)}{}}f_u(x^A)F_u(x^j)$$
(15)
$$(^2+m^2)\mathrm{\Phi }=F_u(x^j)[D+S^1\lambda _n+m^2]f_u(x^A)$$
(16)
This expression must vanish. Since $`F_1\mathrm{}F_d`$ form a basis, it is clear that each $`f_u`$ must be a solution to the equation
$$(D+\lambda _nS^1+m^2)f=0$$
(17)
It follows that each $`f_uF_u`$ is a product solution.
Let $`𝒮[\lambda ]`$ be the space of smooth functions on $`V_1`$ satisfying to (17) for a given value of $`\lambda `$.
Except in the very special case where $`p=1`$, (17) is a partial differential equation and has infinitely many linearly independent solutions.
Developping $`D`$ with help of (10) we obtain after simplification
$$\mathrm{\Delta }_1f+\frac{q}{2}\alpha ^{AB}(_A\mathrm{log}S)_Bf+(\lambda S^1+m^2)f=0$$
(18)
which is a $`p`$-dimensional problem only, formulated in terms of the metric $`\alpha `$. The $`x^j`$ do not arise in this equation.
To summarize, equation (2) has been reduced to a pair of equations involving separate sets of variables. These equations are (14) and (18). For Type I (resp. Type II) the Laplacian in the former is elliptic (resp. hyperbolic) whereas the partial differential operator in the latter is hyperbolic (resp. elliptic). Notice that (14) involves only the geometry of $`V_2`$ , whereas (18) involves not only the geometry of $`V_1`$ but also the shape of the warping function $`S`$.
For Type I with compact $`V_2`$, we write $`𝒮_n`$ for $`𝒮[\lambda _n]`$. The $`n`$th mode will be noted
$$_n=𝒮_n_n$$
Now, a comparison of (18) with (2) is in order: Equation (2) simply involves the $`p+q`$-dimensional Laplacian, which implies that the Gordon current in $`p+q`$ dimensions is conservative. In contrast, (18) not only involves the $`p`$-dimensional Laplacian $`\mathrm{\Delta }_1`$ and an innocent multiplicative operator $`\lambda S^1+m^2`$, but also first order partial differentiation. As a result, the $`p`$-dimensional sesquilinear field constructed with a couple of solutions $`f,h`$ to (18), that is to say
$$I^A=i(f^{}_1^Ahh_1^Af^{})$$
(19)
fails to be divergence-free.
For the same reason, the second order linear differential operator involved in (18) is not symmetric with respect to the scalar product
$$<f,h>_1=_{V_1}f^{}h\sqrt{|\alpha |}d_1^px$$
defined with help of the $`p`$-dimensional volume element $`\sqrt{|\alpha |}d_1^px`$ determined by the metric $`\alpha `$ in $`V_1`$.
This situation is more specially unpleasant in Type I spacetimes, where (18) includes the dynamical aspects of the original equation KG equation. In this case, a scalar product for any couple of solutions to (18) would be of interest for quantum mechanics. So, it is desirable to construct some conserved $`p`$-dimensional current, sesquilinear with respect to the couple $`f,h`$. But the presence of $`f`$ in (18) is an obstacle.
Before we focus our attention on Type I, let us first replace the vector field of formula (19) by a better candidate in order to make up a conservation law.
\[N.B. Statements about the divergence of a vector necessarily refer to a metric. In contradistinction, the divergence of a vector-density is intrinsic.\]
The difficulty associated with the presence of first derivatives in (18) can be circumvented by two manners.
* Use another metric (conformal to $`\alpha _{AB}`$) on the manifold $`V_1`$ and re-write (18) in terms of it. This procedure amounts to consider another differential operator that is symmetric in the sense of this new metric.
* Keep the metric of $`V_1`$ unaltered, but make a conformal change of function, say $`f=S^r\widehat{f}`$.
## 4 Conserved Currents
### 4.1 First Method
Let us consider in $`V_1`$ a new metric $`\stackrel{~}{\alpha }_{AB}`$ such that
$$\alpha _{AB}=U(x^C)\stackrel{~}{\alpha }_{AB}$$
for some conformal factor $`U(x^C)`$ which must be suitably chosen. It is clear that $`\alpha ^{AB}=U^1\stackrel{~}{\alpha }^{AB}`$, if we call $`\stackrel{~}{\alpha }^{AB}`$ the contravariant tensor inverting $`\alpha _{AB}`$. We set $`\stackrel{~}{}_1^A=\stackrel{~}{\alpha }^{AB}_B`$.
The determinants are related by
$$det\alpha _{AB}=U^pdet(\stackrel{~}{\alpha }_{AB})$$
(20)
In view of this formula, we find
$$\mathrm{\Delta }_1f=U^1\stackrel{~}{\mathrm{\Delta }}_1f+(_A\mathrm{log}U^{p/21})\alpha ^{AB}_Bf$$
This is to be inserted into $`Df`$ which is given by (10). We get
$$Df=U^1\stackrel{~}{\mathrm{\Delta }}_1f+(_A\mathrm{log}U^{p/21})\alpha ^{AB}_Bf+\frac{q}{2}(_A\mathrm{log}S)\alpha ^{AB}_Bf$$
(21)
The first derivatives of $`f`$ are eliminated from $`Df`$ provided that
$$U^{p/21}=\mathrm{const}.S^{q/2}$$
which is possible (for nontrivial $`S`$) under the condition that $`p2`$. In this case it is sufficient to take $`U^{p2}=S^q`$, which leads to
$$U=S^{\frac{q}{2p}}$$
(22)
Thus (20) entails
$$\sqrt{|\alpha |}=S^{\frac{pq}{42p}}\sqrt{|\stackrel{~}{\alpha }|}$$
(23)
With the choice (22) we simply have
$$Df=U^1\stackrel{~}{\mathrm{\Delta }}_1f$$
to be inserted into equation (17) for mode $`\lambda `$.
We end up with an equation of the form
$$\stackrel{~}{\mathrm{\Delta }}_1f+Qf=0$$
(24)
where $`Q=U(\lambda S^1+m^2)`$ but equation (22) implies that
$$Q=S^{\frac{q}{2p}}(\lambda S^1+m^2)$$
Notice that $`\stackrel{~}{\mathrm{\Delta }}_1+Q`$ in equation (24) is symmetric with respect to the scalar product $`<\stackrel{~}{f,h}>_1`$, defined with help of the $`p`$-dimensional volume $`\sqrt{|\stackrel{~}{\alpha }|}d_1^px`$. Moreover, the $`p`$-dimensional current
$$\stackrel{~}{I}^A(f,h)=i(f^{}\stackrel{~}{}_1^Ahh\stackrel{~}{}_1^Af^{})$$
(25)
is divergence-free in $`(V_1,\stackrel{~}{\alpha })`$ for any couple of solutions $`f,h`$ to the same equation (24). We mean the same $`\lambda `$ in (24) for $`f`$ and $`h`$ and of course $`\stackrel{~}{}_1^A=\stackrel{~}{\alpha }^{AB}_B`$ . In other words we have
$$_A(\sqrt{|\stackrel{~}{\alpha }|}\stackrel{~}{I}^A)=0$$
Remark: When $`p=1`$ then $`\stackrel{~}{I}^A`$ has a single component only, so the above conservation law reduces to the well-known constancy of the Wronskian.
So, provided $`p2`$, it is possible to eliminate $`\alpha ^{AB}_B\mathrm{log}S`$ from (21) by chosing $`U`$ as in (22).
It is noteworthy that in four-dimensional warped spacetime, we precisely have $`p=2`$ in the case of Class $`B`$. Cases with $`p=2`$ will be handled by the second method (next section).
### 4.2 Second Method
Let us try a conformal change of function involving a suitable power of $`S`$. So we introduce $`\widehat{f}`$ by setting $`f=S^r\widehat{f}`$, for some $`r`$ to be determined below. Our goal is to eliminate $`\widehat{f}`$.
Define in any dimension the Laplace-Beltrami operator on scalars $`\mathrm{\Delta }u=u`$. Well known that
$$\mathrm{\Delta }(uv)=(\mathrm{\Delta }u)v+u\mathrm{\Delta }v+2uv$$
We apply this formula in manifold $`V_1`$, where the Laplacian operator is $`\mathrm{\Delta }_1`$, to the product $`S^r\widehat{f}`$. In the present case
$$\mathrm{\Delta }_1f=(\mathrm{\Delta }_1S^r)\widehat{f}+S^r\mathrm{\Delta }_1\widehat{f}+2\alpha ^{AB}(_AS^r)_B\widehat{f}$$
But (10) tells that
$$Df=\mathrm{\Delta }_1f+\frac{q}{2}\alpha ^{AB}(_A\mathrm{log}S)_Bf$$
where $`_Bf=(_BS^r)\widehat{f}+S^r_B\widehat{f}`$. Thus
$$Df=\mathrm{\Delta }_1f+\frac{q}{2}\alpha ^{AB}(_A\mathrm{log}S)(_BS^r)\widehat{f}+\frac{q}{2}\alpha ^{AB}(_A\mathrm{log}S)S^r_B\widehat{f}$$
First order derivatives of $`\widehat{f}`$ get cancelled from (18) provided that $`2_AS^r+{\displaystyle \frac{q}{2}}S^r_A\mathrm{log}S=0`$. Thus
$$r=\frac{q}{4}$$
(26)
We are left with
$$Df=(\mathrm{\Delta }_1S^r)\widehat{f}+S^r\mathrm{\Delta }_1\widehat{f}+\frac{q}{2}\alpha ^{AB}(_A\mathrm{log}S)(_BS^r)\widehat{f}$$
$$S^rDf=S^r(\mathrm{\Delta }_1S^r)\widehat{f}+\mathrm{\Delta }_1\widehat{f}+\frac{qr}{2}\alpha ^{AB}(_A\mathrm{log}S)(_B\mathrm{log}S)\widehat{f}$$
In view of (26) we conclude that
$$S^{q/4}Df=\mathrm{\Delta }_1\widehat{f}+(S^{q/4}\mathrm{\Delta }_1S^{q/4})\widehat{f}\frac{q^2}{8}\alpha ^{AB}(_A\mathrm{log}S)(_B\mathrm{log}S)\widehat{f}$$
(27)
Equation (12) yields $`(D+\lambda S^1+m^2)f=0`$. Multiply by $`S^{q/4}`$, and use $`f=S^{q/4}\widehat{f}`$. We get
$$S^{q/4}Df+(\lambda S^1+m^2)\widehat{f}=0$$
where the first term is to be developed in terms of $`\widehat{f}`$ as in (27)
We end up with a reduced equation of the form $`\mathrm{\Delta }_1\widehat{f}+\mathrm{\Xi }\widehat{f}=0`$.
If $`\widehat{f}`$ and $`\widehat{h}`$ are two solutions to this equation (for the same $`\lambda `$) they correspond to $`f`$ and $`h`$ through the formulas
$$f=S^r\widehat{f},h=S^r\widehat{h}$$
(28)
and the $`p`$-dimensional current
$$J^A(f,h)=I^A(\widehat{f},\widehat{h})=i(\widehat{f}^{}_1^A\widehat{h}\widehat{h}_1^A\widehat{f}^{})$$
(29)
is conserved
$$_A\sqrt{|}\alpha |J^A=0$$
In order to compute $`J^A`$ we use (28) and write
$$_1^A\widehat{f}^{}=\alpha ^{AB}S^r_Bf^{}+\alpha ^{AB}(_BS^r)f^{}$$
(30)
$$_1^A\widehat{h}=\alpha ^{AB}S^r_Bh+\alpha ^{AB}(_BS^r)h$$
(31)
Inserting into (29) yields
$$iJ^A=S^rf^{}[\alpha ^{AB}S^r_Bh+\alpha ^{AB}(_BS^r)h]S^rh[\alpha ^{AB}S^r_Bf^{}+\alpha ^{AB}(_BS^r)f^{}]$$
After simplification we are left with
$$iJ^A=\widehat{f}^{}S^r_1^Ah\widehat{h}S^r_1^Af^{}$$
(32)
After cancellation of two terms we obtain
$$iJ^A=S^{2r}(f^{}\alpha ^{AB}_Bhh\alpha ^{AB}_Bf^{})$$
Remember that $`2r=\frac{1}{2}q`$. So finally
$$iJ^A=S^{q/2}(f^{}\alpha ^{AB}_Bhh\alpha ^{AB}_Bf^{})$$
(33)
$$J^A=S^{q/2}I^A$$
(34)
Now, under the assumption that (22) and (26) hold true we can assert
###### Proposition 4
If $`p2`$ then the vector fields $`\stackrel{~}{I}`$ on $`(V_1,\stackrel{~}{\alpha })`$ and $`J`$ on $`(V_1,\alpha )`$ correspond to the same (conserved) vector-density, in other words
$$\sqrt{|\alpha |}J^\alpha =\sqrt{|\stackrel{~}{\alpha }|}\stackrel{~}{I}^\alpha $$
Proof
Start from (29), that is
$$iJ^A=\widehat{f}^{}_1^A\widehat{h}\widehat{h}_1^A\widehat{f}^{}$$
use
$$\widehat{f}^{}=S^rf^{},\widehat{h}=S^rh$$
Multiply equation (32) by $`\sqrt{|\alpha |}`$ which is given by (23). Remember (26), hence
$$\frac{pq}{42p}2r=\frac{q}{2p}$$
Thus we find
$$i\sqrt{|\alpha |}J^A=S^{\frac{q}{2p}}(f^{}_1^Ahh_1^Af^{})\sqrt{|\stackrel{~}{\alpha }|}$$
At this stage it is convenient to notice that for all scalar function $`u(x^A)`$ we have
$$_1^Au=S^{\frac{q}{p2}}\stackrel{~}{}_{1}^{}{}_{}{}^{A}u$$
(35)
Taking this formula into account, we can write
$$i\sqrt{|\alpha |}J^A=(f^{}\stackrel{~}{}_1^Ahh\stackrel{~}{}_1^Af^{})\sqrt{|\stackrel{~}{\alpha }|}$$
which proves our assertion, according to (25).
#### 4.2.1 Comparison of both methods
The first method is widely employed in the literature concerning FRW universes, and we also used it in the framework of generalized FRW spacetimes . As we have just checked, the first method and the second one are equivalent when both can be carried out. But the former cannot be carried out when $`p=2`$. So we are forced to consider the latter as more fundamental, since it is not affected by any dimensional exception. In particular it remains the only one available in the important case of a Class $`B`$ four-dimensional spacetime. For a unified approach that encompasses all cases, we are led to a systematic use of the second method.
## 5 Sesquilinear Forms for Type I
In this section we specialize to warped spacetimes of Type I, with $`V_2`$ compact.
The special case $`p=1,q=3`$ concerns generalized FRW spacetimes; it has been treated in and yields results similar to those of the case $`p>1`$ considered below, where the quantity defined by formula (37) is replaced by $`iW(f^{},h)`$ in terms of the Wronskian. Extension to arbitrary $`q`$ is straightforward.
When $`p>1`$ we assume in addition that $`V_1𝐑\times \mathrm{compact}`$.
If $`L`$ is imbedded in $`(V_1,\alpha )`$ as a $`(p1)`$-dimensional spacelike surface, let $`dL_A`$ be the $`(p1)`$-dimensional surface element, that is
$$dL_A=\frac{1}{(p1)!}\sqrt{|\alpha |}\epsilon _{AB_1\mathrm{}B_{p1}}dx^{B_1}\mathrm{}dx^{B_{p1}}$$
(36)
Conservation of $`J^A`$ with respect to $`(V_1,\alpha )`$ implies that
$$(f;h)_1=_LJ^A(f,h)𝑑L_A$$
(37)
doesnot depend on the choice of $`L`$. This expression defines a sesquilinear form on $`𝒮_n`$.
### 5.1 The Gordon Current
Remember usual formulas for $`N`$-dimensional spacetime, with $`\alpha ,\beta =0,1\mathrm{}N1`$.
The space of arbitrary solutions to the KG equation is endowed with a sesquilinear form
$$(\mathrm{\Phi };\mathrm{\Omega })=j^\nu 𝑑\mathrm{\Sigma }_\nu $$
(38)
$$j^\nu (\mathrm{\Phi },\mathrm{\Omega })=i(\mathrm{\Phi }^{}^\nu \mathrm{\Omega }\mathrm{\Omega }^\nu \mathrm{\Phi }^{})$$
Integration is performed over an $`N1`$ dimensional spacelike surface $`\mathrm{\Sigma }`$. Notation $`d^{N1}x=dx^1dx^2\mathrm{}dx^{N1}`$. Be cautious that $`d^{N1}x`$ should not be confused with $`dx^{N1}`$.
Provided $`(\mathrm{\Sigma })`$ is defined by $`x^0=\mathrm{const}.`$ we can write
$$d\mathrm{\Sigma }_0=\frac{1}{(N1)!}\eta _{01\mathrm{}N1}dx^1\mathrm{}dx^{N1}=\sqrt{|g|}d^{N1}x$$
For the space components we first have
$$d\mathrm{\Sigma }_1=\frac{1}{(N1)!}\eta _{1\alpha _1\mathrm{}\alpha _{N1}}dx_1^\alpha \mathrm{}dx_{N1}^\alpha $$
Here the indices $`\alpha _1,\mathrm{}\alpha _{N1}1`$, thus one of them all must be $`0`$ (otherwize they would not be all different, so $`\eta _{1\alpha _1\mathrm{}\alpha _{N1}}`$ would vanish). Thus $`d\mathrm{\Sigma }_1`$ has $`dx^0`$ as a factor, say $`d\mathrm{\Sigma }_1=(\mathrm{}..)dx^0`$. But $`dx^0`$ is zero on $`(\mathrm{\Sigma })`$. Thus finally $`d\mathrm{\Sigma }_1`$ vanishes on $`(\mathrm{\Sigma })`$. By a similar argument we check that $`d\mathrm{\Sigma }_2`$, etc $`\mathrm{}d\mathrm{\Sigma }_{N1}`$ also vanish on $`(\mathrm{\Sigma })`$.
In order to integrate over $`(\mathrm{\Sigma })`$ we evaluate the differential form $`jd\mathrm{\Sigma }`$ on this surface; we can write
$$jd\mathrm{\Sigma }=j^0d\mathrm{\Sigma }_0$$
$$j^0=i(\mathrm{\Phi }^{}^0\mathrm{\Omega }\mathrm{\Omega }^0\mathrm{\Phi }^{})$$
Naturally $`^0\mathrm{\Phi }=g^{0\alpha }_\alpha \mathrm{\Phi }`$.
So far the formulas of this subsection are general. When we specify to warped spacetimes of Type I, then $`N=p+q`$ and $`V_1`$ is Lorentzian; $`V_2`$ is Riemannian so the label $`k`$ cannot be $`0`$. In coordinates adapted to the warped product structure, we know that $`g^{Ak}`$ all vanish; in particular $`g^{0k}`$ vanishes. We are left with
$$^0\mathrm{\Phi }=g^{0A}_A\mathrm{\Phi }$$
Similarly
$$^0\mathrm{\Omega }=g^{0A}_A\mathrm{\Omega }$$
We thus have
$$jd\mathrm{\Sigma }=i\sqrt{|g|}d^{n1}x(\mathrm{\Phi }^{}g^{0A}_A\mathrm{\Omega }\mathrm{\Omega }g^{0A}_A\mathrm{\Phi }^{})$$
(39)
Cf. (2.2) in ref. .
Let us now consider product solutions.
So we assume that $`\mathrm{\Phi }`$ and $`\mathrm{\Omega }`$ are solutions to (2) in the form
$$\mathrm{\Phi }=f(x^A)F(x^j)\mathrm{\Omega }=h(x^B)H(x^k)$$
but not necessarily on the same mode, say $`\mathrm{\Phi }_n,\mathrm{\Omega }_l`$ including the possibility that $`nl`$. We have
$$g^{0A}_A\mathrm{\Omega }=Hg^{0A}_Ah$$
$$g^{0A}_A\mathrm{\Phi }^{}=F^{}g^{0A}_Af^{}$$
Inserting into (39) yields
$$jd\mathrm{\Sigma }=i(\mathrm{\Phi }^{}Hg^{0A}_Ah\mathrm{\Omega }F^{}g^{0A}_Af^{})\sqrt{|g|}d^{n1}x$$
Develop $`\sqrt{|g|}`$ according to (4) where $`S`$ and $`\gamma `$ are positive, so
$$\sqrt{|g|}=\sqrt{|\alpha |}S^{q/2}\sqrt{\gamma }$$
According to (5) we can write $`d^{N1}x=\omega d_2^qx`$, therefore
$$jd\mathrm{\Sigma }=i(F^{}f^{}Hg^{0A}_AhhHF^{}g^{0A}_Af^{})\sqrt{|\alpha |}S^{q/2}\sqrt{\gamma }\omega d_2^qx$$
$$jd\mathrm{\Sigma }=i(f^{}g^{0A}_Ahhg^{0A}_Af^{})F^{}H\sqrt{|\alpha |}S^{q/2}\sqrt{\gamma }\omega d_2^qx$$
We now turn to integration.
If $`p=1`$ we can take $`\mathrm{\Sigma }=\{t_0\}\times V_2`$ where $`t_0`$ is a fixed value of the time coordinatr $`x^0`$. This case has been investigated in details in .
If $`p>1`$, let us take
$$\mathrm{\Sigma }=L\times V_2$$
where $`LV_1`$ is the submanifold defined by $`x^0=\mathrm{const}`$. Indeed this choice is compatible with the assumption made above that the equation of $`\mathrm{\Sigma }`$ in $`V_1\times V_2`$ is just $`x^0=\mathrm{const}`$.
Integrate the above formula; we obtain
$$(\mathrm{\Phi };\mathrm{\Omega })=i\left(_{V_2}F^{}H\sqrt{\gamma }d_2^qx\right)_L(f^{}g^{0A}_Ahhg^{0A}_Af^{})S^{q/2}\sqrt{|\alpha |}\omega $$
(40)
$$(\mathrm{\Phi };\mathrm{\Omega })=i<F,H>_2_L(f^{}g^{0A}_Ahhg^{0A}_Af^{})S^{q/2}\sqrt{|\alpha |}\omega $$
(41)
where we have factorized out the scalar product in $`(V_2,\gamma )`$
$$<F,H>_2=_{V_2}\sqrt{\gamma }F^{}Hd_2^qx$$
(42)
well-defined and positive for arbitrary couple of functions on $`V_2`$.
When $`F`$ and $`H`$ belong to $`_n`$ and $`_l`$ with $`nl`$ then $`<F,H>_2`$ vanishes. It follows that if two product solutions $`\mathrm{\Phi },\mathrm{\Omega }`$ belong to distinct modes, $`(\mathrm{\Phi };\mathrm{\Omega })`$ vanishes. Now using (15) and a similar development
$$\mathrm{\Omega }=\underset{s=1}{\overset{d(l)}{}}h_sH_s$$
over a basis $`H_1,\mathrm{}`$ of $`_l`$, we easily check that this property holds for any couple of mode solutions, in other words
###### Proposition 5
For Type I under our assumptions, two different modes are orthogonal in the sense of the sesquilinear form defined by the Gordon current.
\[This result extends Proposition 1 of Ref. \]
When $`\mathrm{\Phi }`$ and $`\mathrm{\Omega }`$ belong to the same mode we give to (41 a more compact formulation. A look at (33) (34) enables one to write
$$jd\mathrm{\Sigma }=j^0d\mathrm{\Sigma }_0=J^0\sqrt{|\alpha |}\omega (F^{}H\sqrt{\gamma }d_2^qx)$$
(43)
We obtain by integration
$$_\mathrm{\Sigma }jd\mathrm{\Sigma }=<F,H>_2_LJ^0\sqrt{|\alpha |}\omega $$
with At this stage it is convenient to remember that $`\omega `$ is given by (5) and to observe that
$$dL_0=\sqrt{|\alpha |}\omega $$
When $`A0`$, if $`B_1\mathrm{}B_{p1}`$ are all $`A`$, one of them must be $`0`$, thus $`dx^0`$ is a factor in $`dL_A`$.
On $`(L)`$, we can write $`dx^0=0`$, thus $`dL_A`$ vanishes for $`A0`$. Thus we have on this manifold $`J^AdL_A=J^0dL_0`$. Hence
$$_\mathrm{\Sigma }jd\mathrm{\Sigma }=<F,H>_2_LJ^AdL_A$$
A glance at (37) enables us to state
###### Theorem 1
In Type I warped spacetime, for product solutions, the sesquilinear map defined through the usual Gordon current gets factorized according to the formula
$$(\mathrm{\Phi };\mathrm{\Omega })=(f;h)_1<F,H>_2$$
where $`(f;h)_1`$ is defined by (37).
## 6 Extension to further couplings.
Up to now we have focussed on the minimal coupling of a free particle to gravity. Extending our results to an equation of the form
$$(^2+m^2+a(x))\mathrm{\Psi }=0$$
(44)
is straightforward provided that the additional term $`a`$ does not depend on the $`x^j`$’s. Otherwize, $`\mathrm{\Delta }_2`$ would not commute any more with $`^2+a(x)`$, and the mode decomposition would be impossible.
Assuming that $`a=a(x^A)`$, the differential operator in (44) still commutes with $`\mathrm{\Delta }_2`$. Equation (12) is replaced by
$$(D+\lambda S^1+m^2+a(x^A))\mathrm{\Phi }=0$$
Any product solution remains a mode. In (17)(18) $`m^2`$ must be replaced by $`m^2+a(x^A)`$ but the variables $`x^j`$ remain ignorable. Since $`a`$ acts on $`f`$ as a multiplicative operator, the status of the first derivatives $`f`$ in equation (18) is not modified; therefore condition (26) still avoid the occurence of them. Moreover formula (27) which determines $`Df`$ in terms of $`\widehat{f}`$ is an identity valid for arbitrary $`f`$. Again we end up with an equation of the form $`\mathrm{\Delta }_1\widehat{f}+\mathrm{\Xi }\widehat{f}=0`$ where $`\mathrm{\Xi }`$ now involves the additional term $`a`$. As well-known, this still ensures conservation of $`J^A`$, given without modification by (29). Finally the results of Section 4 remain valid.
These remarks permit to treat the case where the KG equation includes an external potential of the special (and rather artificial) form $`a(x^A)`$. It is more interesting to notice that our study could remain valid for curvature coupling, characterized by the addition to equation (2) of a curvature term $`\xi R_\alpha ^\alpha `$ where $`\xi `$ is a constant <sup>2</sup><sup>2</sup>2For dimensional reasons, $`\xi `$ necessarily vanishes in the limit where $`\mathrm{}0`$, thus curvature coupling cannot affect the geodesic motion, and $`R_\alpha ^\alpha `$ is the scalar curvature of $`(V,g)`$. But the condition which legitimates this extension is that the scalar curvature of spacetime depends only on the $`x^A`$’s.
We are thus led to investigate what kind of warped spacetime supports a scalar curvature of the form $`R(x^A)`$.
The following lemma will be useful.
###### Lemma 1
Let $`u(x^A)`$ be a function on $`V`$ satisfying $`_ju=0`$. Then the following quantities
$$^A_Bu,^Au,^A_Au$$
are independent from the coordinates $`x^k`$.
Proof:
$`_A_Bu=_A_Bu\mathrm{\Gamma }_{A}^{\mu }{}_{B}{}^{}_\mu u`$. The first term is obviously independent from $`x^k`$’s. Since $`_ju`$ is zero, only the coefficients $`\mathrm{\Gamma }_{A}^{C}{}_{B}{}^{}`$ give a contribution to the second term. But it is known that $`\mathrm{\Gamma }_{A}^{C}{}_{B}{}^{}=^1\mathrm{\Gamma }_{A}^{C}{}_{B}{}^{}`$, where $`{}_{}{}^{1}\mathrm{\Gamma }`$ is the connexion for $`(V_1,\alpha )`$ (see ; in 4 dimensions see formula (26) in . Thus the second term has the same property.
The statement about $`^Au`$ is obvious, if we remember that $`g^{AB}=\alpha ^{AB},g^{A0}=0`$.
Finally, we have $`^A_Au=\alpha ^{AB}_B_Au`$ which is independent of $`x^k`$ because of the first statement.
The Ricci tensor of a warped product can be expressed in terms of the warping function and of the geometry of the factor manifolds (See O’Neill , especially corollary 43, p. 211. When $`p+q=4`$, see Carot and da Costa , eqs (25), for a more transparent notation).
We are concerned with
$$R=\alpha ^{AB}R_{AB}S^1\gamma ^{ij}R_{ij}$$
(45)
In the first term, we know that $`R_{AB}`$ differs from $`{}_{}{}^{1}R_{AB}^{}`$ only by a function of $`e^\mathrm{\Theta }`$ and $`_B_A(e^\mathrm{\Theta })`$, which cannot depend on $`x^C`$ according to the previous result. Therefore $`R_{AB}`$ doesnot depend on the $`x^k`$’s
$$_jR_{AB}=0$$
This result, in turn, entails that $`R`$ might depend on these variables only through the quantity $`\gamma ^{ij}R_{ij}`$.
According to the literature cited, $`R_{ij}`$ takes on the form
$$R_{ij}=^2R_{ij}+N\gamma _{ij}$$
where $`N`$ is a function of $`\mathrm{\Theta },^A\mathrm{\Theta },_A\mathrm{\Theta }`$and $`^A_A\mathrm{\Theta }`$. By the above lemma, we see that
$$_jN=0$$
By contraction we obtain
$$\gamma ^{ij}R_{ij}=\gamma ^{ij}{}_{}{}^{2}R_{ij}^{}+qN$$
It is now clear that $`\gamma ^{ij}R_{ij}`$, and therefore equally $`R`$, depends on the $`x^k`$’s except when $`{}_{}{}^{2}R`$ is a constant; thus
###### Theorem 2
(Only) if $`(V_2,\gamma )`$ has constant scalar curvature $`\gamma ^{ij}{}_{}{}^{2}R_{ij}^{}`$, the scalar curvature $`R`$ of the warped product can be regarded as a function on $`V_1`$.
An well-known instance of this situation is the case of conventional FRW spacetime, where the space sections are homogeneous.
## 7 Concluding remarks
We have separated the variables in the KG equation with help of a generalized mode decomposition which is possible essentially because the motion of a test particle in a warped product spacetime admits a remarkable first integral.
At least for Type I, these modes are actually ”normal” , that is orthogonal, in the sense of the usual sesquilinear form associated with the Gordon current. This form itself has been analysed in terms of the vector field $`J`$ defined on the Lorentzian factor manifold and, under very large assumptions, a sesquilinear form has been defined on the solutions to the reduced wave equation.
Type I is perhaps physically the most interesting; in 4 dimensions it encompasses not only generalized FRW spacetimes, but also all kind of spherically symmetric spacetime, including the nonstationary ones. However further investigations of Type II might be of interest.
It is noteworthy that the whole picture remains valid in the presence of a curvature coupling term $`\xi R_\alpha ^\alpha `$ only under the condition that $`(V_2,\gamma )`$ has a constant scalar curvature. For FRW, abandoning spatial homogenity would not permit to carry out our mode decomposition when a curvature coupling term is introduced into the KG equation. In contrast, the mode decomposition remains possible with such a term for all kind of spherically symmetric spacetime.
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# 1 Introduction
## 1 Introduction
Long ago, the question has been addressed of constructing geometric field theories of gravity from the non-symmetric metric
$$G_{\mu \nu }g_{\mu \nu }+B_{\mu \nu },$$
(1.1)
where $`g_{\mu \nu }`$ is the symmetric part of $`G_{\mu \nu }`$ and $`B_{\mu \nu }`$ is an antisymmetric tensor field. Such a model was first discussed in 1925 by Einstein in an attempt to unify gravity and electromagnetism, and developed by himself and others . In some sense, modern string theory has revived this attempt at geometric unification of forces
But the couplings of a gauge field like a massless 2-form $`B_{\mu \nu }`$ and a (symmetric) metric $`G_{\mu \nu }`$ are severely restricted by consistency requirements, thereby disqualifying the old theory of Einstein and its modern versions (like the “nonsymmetric gravity theory” of Moffat ). The first proof of the physical inconsistency of these models has been given by T. Damour, S. Deser and J. McCarthy . They have shown that all “geometric” actions (see for a precise explanation of the sense of this term) homogeneous of order two in the number of derivatives violate standard physical requirements. In this paper, we confirm their results while relaxing some of the assumptions made by these authors.
To analyze the generic geometric models, they had to expand the action in powers of $`B_{\mu \nu }`$ about a classical symmetric background, and consider the merits of the resulting theory, that means study the consistency of the deformations of a theory of gravity coupled to a differential 2-form in terms of standard physical criteria: absence of negative-energy excitations and coherence of the degree-of-freedom content. The first point of their argument is that the absence of negative-energy excitations requires (whatever the gravitational background) the $`B`$ expansion to begin with the (quadratic) kinetic term $`H_{\mu \nu \lambda }H^{\mu \nu \lambda }`$ where $`H_{\lambda \mu \nu }_\lambda B_{\mu \nu }+_\mu B_{\nu \lambda }+_\nu B_{\lambda \mu }`$ is the field strength. This action has the usual gauge invariance $`\delta B_{\mu \nu }=_\mu ϵ_\nu _\nu ϵ_\mu `$. Secondly, geometric actions homogeneous in two derivatives generate powers of the undifferentiated $`B_{\mu \nu }`$. These higher-power terms generically violate the gauge invariance of the leading kinetic term. Furthermore, it is even impossible to overcome this problem by a deformation of the abelian gauge invariance .
A modern useful tool to study consistent deformations is the BRST field-antifield formalism. It is known that if $`s`$ is the BRST operator of the free theory, the first order vertices of consistent deformations are constrained to belong to the cohomological group $`H_0^n(s|d)`$ of $`s`$ modulo $`d`$ in ghost number $`0`$ and form degree $`n`$ (the dimension of spacetime) . By using the powerful tool of homological algebra, we completely determine all non trivial consistent deformations at first order in the coupling constant for a system of abelian 2-forms coupled to gravity for $`n4`$.<sup>1</sup><sup>1</sup>1For $`n<4`$, the 2-forms do not carry any degree of freedom, hence this case is not considered.
Our result is that all non trivial first-order consistent deformations are of the four following types:
* Strictly invariant polynomials in the curvatures $`R_{\mu \nu \rho \sigma }`$ and $`H_{\mu \nu \rho }^A`$.
* Couplings invariant under the abelian transformation up to a divergence. They generalize the purely gravitational deformations previously known in the literature (denoted as $`𝒜_{chiral}`$ in ) and the Chern-Simons type self-coupling of the 2-forms (which exist only for specific values of $`n`$).
* Exotic couplings of the 2-form gauge fields to gravitational Chern-Simons forms (see ). These interactions deform the gauge transformations of the free theory in a non trivial way.
* The Freedman-Townsend coupling (only allowed for $`n=4`$), where the abelian gauge transformations for the 2-forms are non trivially deformed.
We notice that non trivial deformations of “gravity + 2-form” are at most linear in the undifferentiated $`B_{\mu \nu }`$ which confirms the results of Damour et al. that geometric actions including higher powers of the undifferentiated $`B_{\mu \nu }`$ are physically inconsistent or trivial. No assumption is made here on the degree of homogeneity in the derivatives nor on the form of the gravity action (except that it must define a normal theory, according to the terminology of ).
In the next section, we determine all the cohomological class of $`H(s|d)`$ in ghost number $`0`$ and form degree $`n`$ and we relate those classes to the first-order consistent deformations of the theory.
## 2 The model
Throughout the paper we shall work in the vielbein formulation of gravity. Denoting the vielbein fields by $`e_\mu ^a`$, we define the inverse vielbeins, the metric, the Christoffel connection and the spin connection through,
$`E_a^\mu e_\mu ^b=\delta _b^a,e_\mu ^aE_a^\nu =\delta _\mu ^\nu ,g_{\mu \nu }=e_{\mu a}e_\nu ^a,g^{\mu \nu }=E^{a\mu }E_a^\nu ,`$ (2.1)
$$_\mu e_\nu ^a\omega _{\mu b}^ae_\nu ^b\mathrm{\Gamma }_{\mu \nu }^\rho e_\rho ^a=0,\omega _\mu ^{ab}=\omega _\mu ^{ba},\mathrm{\Gamma }_{\mu \nu }^\rho =\mathrm{\Gamma }_{\nu \mu }^\rho ,$$
(2.2)
where Lorentz indices ($`a,b,\mathrm{}`$) are raised and lowered with the Minkowski metric.
As we pointed out in the introduction, our calculations are valid for a wide range of gravity action. For simplicity, the model we consider here is the standard Einstein-2-form system. Its lagrangian reads,
$$_0/e=\frac{1}{2}R\frac{1}{12}g^{\mu \lambda }g^{\nu \sigma }g^{\rho \theta }H_{\mu \nu \rho }^AH_{A\lambda \sigma \theta },$$
(2.3)
where $`e`$ is the determinant of the vielbeins, $`R=R_{\mu \nu }^{ab}E_a^\mu E_b^\nu `$ is the Riemann curvature and $`H_{\mu \nu \rho }^A`$ are the $`2`$-form field strengths.
Because of the gauge invariance of the classical action under the diffeomorphisms, local Lorentz transformations and gauge transformations of the $`2`$-form potentials ($`B_{\mu \nu }^AB_{\mu \nu }^A+_\mu ϵ_\nu ^A_\nu ϵ_\mu ^A`$), we are led to introduce within the field-antifield formalism the following set of ghosts: $`\xi ^\mu `$ (diffeomorphisms ghosts), $`C^{ab}`$ (Lorentz ghosts), $`B_\mu ^A,B^A`$ ($`2`$-form ghosts). The extra ghosts $`B^A`$ (“ghosts of ghosts”) required in the 2-form sector arise because the gauge transformations of the $`2`$-form potentials are reducible. Indeed, they vanish when the gauge parameters are set to $`ϵ_\mu ^A=_\mu \lambda ^A`$ where $`\lambda ^A`$ are arbitrary functions. In the sequel we will denote by $`\mathrm{\Phi }^\mathrm{\Gamma }`$ all the fields and ghosts: $`\{\mathrm{\Phi }^\mathrm{\Gamma }\}=\{e_\mu ^a,B_{\mu \nu }^A,\xi ^\mu ,C^{ab},B_\mu ^A,B^A\}`$. The corresponding antifields will be denoted by $`\mathrm{\Phi }_\mathrm{\Gamma }^{}`$: $`\{\mathrm{\Phi }_\mathrm{\Gamma }^{}\}=\{e_a^\mu ,B_A^{\mu \nu },\xi _\mu ^{},C_{ab}^{},B_A^\mu ,B_A^{}\}`$. In terms of these variables the minimal solution of the BRST master equation reads:
$`S`$ $`=`$ $`{\displaystyle }d^nx(_0\xi ^\mu \mathrm{\Phi }_\mathrm{\Gamma }^{}_\mu \mathrm{\Phi }^\mathrm{\Gamma }+e_a^\mu (e_\nu ^a_\mu \xi ^\nu +e_\mu ^bC_b^a)+{\displaystyle \frac{1}{2}}C_K^{}C^IC^Jf_{IJ}^K`$ (2.4)
$`+B_A^{\mu \nu }(_\mu B_\nu ^A+_\mu \xi ^\rho B_{\rho \nu }^A)B_A^\mu (_\mu B^A+_\mu \xi ^\rho B_\rho ^A)).`$
In the above expression the $`f_{IJ}^K`$ denote the structure constants of $`𝒢_L`$. We have also used the notation, $`C^I=C^{ab}.`$ The action of the BRST differential on the variables is then defined as: $`s\mathrm{\Phi }^A=\frac{\delta ^RS}{\delta \mathrm{\Phi }_A^{}},s\mathrm{\Phi }_A^{}=\frac{\delta ^RS}{\delta \mathrm{\Phi }^A}`$.
In order to obtain the possible deformations of (2.3) we need to compute the local BRST cohomology $`H(s|d)`$ in the algebra of local forms. These are by definition linear combinations of spacetime forms $`\omega ^\tau (x,dx)`$ with coefficients that are local functions. By local functions it is meant functions which depend polynomially on the variables $`_{\mu _1\mathrm{}\mu _j}\mathrm{\Phi }^\mathrm{\Gamma }`$ and $`_{\mu _1\mathrm{}\mu _j}\mathrm{\Phi }_\mathrm{\Gamma }^{}`$ except the undifferentiated vielbeins $`e_\mu ^a`$ for which a smooth and regular dependence in an open neighborhood of some regular background configuration ($`\text{det}(e_\mu ^a)0`$) is allowed. We shall denote the algebra of local forms by $``$ and the algebra of local functions by $`𝒜`$. Thus: $`=\mathrm{\Omega }(M)𝒜;aa=_\tau \omega ^\tau \alpha _\tau ,\omega ^\tau \mathrm{\Omega }(M),\alpha _\tau 𝒜.`$
For gravitational theories it turns out that the local BRST cohomology in form degree $`n`$ and ghost number $`g`$ can easily be obtained from the cohomology $`H(s)`$ of the BRST differential $`s`$ . Indeed, one has the following isomorphism<sup>2</sup><sup>2</sup>2This isomorphism holds because we assume the spacetime manifold to be homeomorphic to $`𝐑^n`$.:
$$H_g^n(s|d,)\frac{H^{g+n}(s,𝒜)}{𝐑}$$
(2.5)
In particular, this isomorphism holds in ghost number $`0`$ and therefore for the consistent deformations. It is implemented as follows: if $`\alpha ^n`$ is a representative of $`H^n(s,𝒜)`$, then the corresponding representative of $`H_0^n(s|d,)`$ is given by $`a_0^n=\frac{1}{n!}b^n\alpha ^n`$ where $`b`$ is the operator defined by $`b=dx^\mu \frac{}{\xi ^\mu }`$ .
The first part of the analysis of $`H(s)`$ consists in finding new generators of the algebra $`𝒜`$ which isolate a contractible part of the algebra with respect to the differential $`s`$. Taking into account the results of we define the new basis of generators of $`𝒜`$ as:
$`\{𝒯^r\}`$ $`=`$ $`\{D_{a_1}\mathrm{}D_{a_k}R_{ab}^I,D_{a_1}\mathrm{}D_{a_k}H_{bcd}^A:k=0,1,\mathrm{}\},`$ (2.6)
$`\{𝒯_{\overline{r}}^{}\}`$ $`=`$ $`\{D_{a_1}\mathrm{}D_{a_k}\widehat{\mathrm{\Phi }}_\mathrm{\Gamma }^{}:k=0,1,\mathrm{}\},`$ (2.7)
$`\widehat{\xi }^a`$ $`=`$ $`\xi ^\mu e_\mu ^a,\widehat{C}^I=C^I+\xi ^\mu \omega _\mu ^I,`$ (2.8)
$`\widehat{B}^A`$ $`=`$ $`B^A+\xi ^\mu B_\mu ^A+{\displaystyle \frac{1}{2}}\xi ^\nu \xi ^\mu B_{\mu \nu }^A,`$ (2.9)
$`\{U_l\}`$ $`=`$ $`\{_{(\mu _1\mathrm{}\mu _k}e_{\mu )}^a,_{(\mu _1\mathrm{}\mu _k}\omega _{\mu )}^I,_{(\mu _1\mathrm{}\mu _k}B_{\nu _1)}^A,_{(\mu _1\mathrm{}\mu _k}B_{[\nu _1)\nu _2]}^A:`$ (2.10)
$`k=0,1,\mathrm{}\},`$
$`\{V_l\}`$ $`=`$ $`\{_{(\mu _1\mathrm{}\mu _k}se_{\mu )}^a,_{(\mu _1\mathrm{}\mu _k}s\omega _{\mu )}^I,_{(\mu _1\mathrm{}\mu _k}sB_{\nu _1)}^A,_{(\mu _1\mathrm{}\mu _k}sB_{[\nu _1)\nu _2]}^A:`$ (2.11)
$`k=0,1,\mathrm{}\},`$
where,
$`H_{abc}^A`$ $`=`$ $`E_a^\mu E_b^\nu E_c^\rho H_{\mu \nu \rho }^A`$ (2.12)
$`R_{ab}^I`$ $`=`$ $`R_{ab}^{cd}=E_a^\mu E_b^\nu R_{\mu \nu }^{cd},`$ (2.13)
$`\omega _\mu ^I`$ $`=`$ $`\omega _\mu ^{ab},`$ (2.14)
$`\{\widehat{\mathrm{\Phi }}_\mathrm{\Gamma }^{}\}`$ $`=`$ $`\{\widehat{e}_a^b,\widehat{B}_A^{ab},\widehat{\xi }_a^{},\widehat{C}_I^{},\widehat{B}_A^a,\widehat{B}_A^{}\},`$ (2.15)
$`\widehat{e}_a^b`$ $`=`$ $`e_\mu ^be_a^\mu /e,\widehat{C}_I^{}=C_I^{}/e,`$ (2.16)
$`\widehat{B}_A^{ab}`$ $`=`$ $`e_\mu ^ae_\nu ^bB_A^{\mu \nu }/e,\widehat{B}_A^a=e_\mu ^aB_A^\mu /e,\widehat{B}_A^{}=B_A^{}/e`$ (2.17)
$`\widehat{\xi }_a^{}`$ $`=`$ $`E_a^\mu (\xi _\mu ^{}A_\mu ^IC_I^{}B_{\mu \nu }^AB_A^\nu B_\mu ^AB_A^{})/e`$ (2.18)
It is easily shown that any local function can be expressed in terms of the new generators.
The first advantage of the new basis is that on all but the variables $`U_l`$ and $`V_l`$, the action of $`s`$ takes the familiar form:
$$s=\delta +\gamma ,$$
(2.19)
with,
$`\delta 𝒯^r=\delta \widehat{\xi }^a=\delta \widehat{C}^I=\delta \widehat{B}^A=0,`$ (2.20)
$`\delta \widehat{e}_a^b=(R_a^b{\displaystyle \frac{1}{2}}\delta _a^bR){\displaystyle \frac{1}{2}}H_{acd}^AH_A^{bcd}{\displaystyle \frac{1}{12}}\delta _a^bH_{cde}^AH_A^{cde},`$ (2.21)
$`\delta \widehat{C}_{ab}^{}=2\widehat{e}_{[ab]}^{},\delta \widehat{\xi }_a^{}=D_b\widehat{e}_a^b+{\displaystyle \frac{1}{p_A!}}\widehat{B}_A^{b_1b_2}\widehat{H}_{ab_1b_2}^A`$ (2.22)
$`\delta \widehat{B}_A^{a_1a_2}=D_b\widehat{H}_A^{ba_1a_2},\delta \widehat{B}_A^{a_1}=D_b\widehat{B}_A^{ba_1},\delta \widehat{B}_A^{}=D_b\widehat{B}_A^b,`$ (2.23)
$`\delta D_{a_1}\mathrm{}D_{a_k}\widehat{\mathrm{\Phi }}_\mathrm{\Gamma }^{}=D_{a_1}\mathrm{}D_{a_k}\delta \widehat{\mathrm{\Phi }}_\mathrm{\Gamma }^{},`$ (2.24)
$`\gamma 𝒯^r=(\widehat{\xi }^aD_a+\widehat{C}^I\delta _I)𝒯,\gamma 𝒯_{\overline{r}}^{}=(\widehat{\xi }^aD_a+\widehat{C}^I\delta _I)𝒯_{\overline{r}}^{}`$ (2.25)
$`\gamma \widehat{\xi }^a=\widehat{C}_b^a\widehat{\xi }^b,\gamma \widehat{C}^I={\displaystyle \frac{1}{2}}\widehat{C}^J\widehat{C}^Kf_{KJ}^I+\widehat{F}^I,\gamma \widehat{B}^A=\widehat{H}^A,`$ (2.26)
where
$$\widehat{R}^I=\frac{1}{2}\widehat{\xi }^c\widehat{\xi }^dR_{cd}^I,\widehat{H}^A=\frac{1}{6}\widehat{\xi }^c\widehat{\xi }^b\widehat{\xi }^aH_{abc}^A.$$
(2.27)
In the above equations, the notation $`\widehat{C}^I\delta _I𝒯^r`$ stands for the Lorentz infinitesimal transformation of $`𝒯^r`$ with the infinitesimal parameter replaced by the ghost $`\widehat{C}^I`$. For example, if $`𝒯^r=𝒯^a`$ is a contravariant vector then $`\widehat{C}^I\delta _I𝒯^a=\widehat{C}_b^a𝒯^b`$.
The gradings associated to $`\delta `$ and $`\gamma `$ are respectively the antighost number (denoted antigh) and the pureghost number (denoted puregh); their sum is equal to the ghost number (denoted gh). The table summarizes the various gradings associated to the operators, the fields, the ghosts and the antifields.
The second advantage of the new basis is that it exhibits a manifestly contractible part of the algebra $`𝒜`$. Indeed, by construction, the variables $`U_l`$ and $`V_l`$ are mapped on each other by the BRST differential,
$$sU_l=V_l,sV_l=0,$$
(2.28)
and since the BRST differential does not mix the $`U_l`$ and $`V_l`$ with the rest of the variables, each couple $`\{U_l,V_l\}`$ drops out from the cohomology and we have $`H(s,𝒜)H(s,𝒜_2)`$ where $`𝒜_2`$ is the algebra generated by the set $`\{𝒯^r,𝒯_{\overline{r}}^{},\widehat{\xi }^a,\widehat{C}^I,\widehat{B}^A\}`$.
To summarize, the above discussion indicates that to obtain $`H_0^n(s|d,)`$ we only need to calculate $`H^n(s,𝒜_2)`$. This is the subject of the next section.
## 3 BRST cohomology in $`𝒜_2`$ and consistent vertices
In order to get the cohomology $`H^n(s,𝒜_2)`$ we need to solve the equation,
$$s\alpha ^n=0,\alpha ^n𝒜_2,$$
(3.1)
where two solutions of (3.1) are identified if they differ by an $`s`$-exact contribution, i.e, $`\alpha ^n\alpha ^n+s\beta ^{n1}`$ with $`\beta ^{n1}𝒜_2`$.
The approach we follow is identical to the one developed in so we only emphasize here the main ideas and the new results. A more in-depth presentation will be given in .
First, the cocycle $`\alpha ^n`$ is decomposed according to a degree called the $`\widehat{\xi }^a`$-degree which counts the polynomial degree in the variables $`\widehat{\xi }^a`$,
$$\alpha ^n=\underset{k=l}{\overset{n}{}}\alpha _k.$$
(3.2)
According to the $`\widehat{\xi }^a`$-degree, the BRST differential decomposes into four parts, $`s=s_0+s_1+s_2+s_3`$, which can be read off from (2.20)-(2.26). The first term is given by,
$$s_0=\delta +\gamma _L$$
(3.3)
where $`\delta `$ is the Koszul-Tate differential and $`\gamma _L`$ is the longitudinal exterior derivative along the gauge orbits of $`𝒢_L`$,
$$\gamma _L=\frac{1}{2}\widehat{C}^J\widehat{C}^Kf_{JK}^I\frac{}{\widehat{C}^I}+\widehat{C}^I\delta _I,$$
(3.4)
with,
$$\delta _{ab}\widehat{\xi }_c=\eta _{bc}\widehat{\xi }_a\eta _{ac}\widehat{\xi }_b,\delta _I\widehat{C}^J=f_{IK}^J\widehat{C}^K.$$
(3.5)
$`s_1`$ plays the rôle of an exterior covariant derivative whose differentials are the $`\widehat{\xi }^a`$. Its action is given by,
$$s_1𝒯^r=\widehat{\xi }^aD_a𝒯^r,s_1𝒯_{\overline{r}}^{}=\widehat{\xi }^aD_a𝒯_{\overline{r}}^{},s_1\widehat{\xi }^a=s_1\widehat{C}^I=s_1\widehat{B}^A=0.$$
(3.6)
Finally, the operators $`s_2`$ and $`s_3`$ are given by,
$$s_2=\widehat{R}^I\frac{}{\widehat{C}^I},s_3=\widehat{H}^A\frac{}{\widehat{B}^A}.$$
(3.7)
According to the $`\widehat{\xi }^a`$-degree, eq. (3.1) decomposes into the following tower of equations:
$`0`$ $`=`$ $`s_0\alpha _l,`$ (3.8)
$`0`$ $`=`$ $`s_0\alpha _{l+1}+s_1\alpha _l,`$ (3.9)
$`0`$ $`=`$ $`s_0\alpha _{l+2}+s_1\alpha _{l+1}+s_2\alpha _l,`$
$`\mathrm{}`$
Up to trivial terms which only modify components of higher $`\widehat{\xi }^a`$-degree, eq. (3.8) indicates that $`\alpha _l`$ is an element of the cohomology $`H(s_0,𝒜_2)`$. Given the definition of $`s_0`$, this cohomology is analyzed in a very similar fashion as the standard BRST cohomology for non-gravitational theories. Using the acyclicity of $`\delta `$ in antighost number $`k>0`$, one can show that each class of $`H(s_0,𝒜_2)`$ admits an antifield independent representative so that the cocycle condition becomes, $`\gamma _L\alpha _l=0`$. This equation is the well known coboundary condition for the Lie algebra cohomology of $`𝒢_L`$ in a $`𝒢_L`$-module so the most general form for the non-trivial $`\alpha _l`$ is,
$$\alpha _l=\alpha _l^i(\widehat{\xi }^a,𝒯^r)\omega _i(\theta _K(\widehat{C}^I),\widehat{B}^A),\delta _I\alpha _l^i(\widehat{\xi }^a,𝒯^r)=0.$$
(3.11)
In (3.11), the $`\omega _i`$ are polynomials in the $`\widehat{B}^A`$ and the $`\theta _K(\widehat{C}^I)`$ which are the primitive elements of the Lorentz Lie algebra cohomology. In $`n=2r`$ and $`n=2r+1`$ dimensions, they are given by,
$`\theta _K(C)`$ $`=`$ $`C_{a_1}^{a_2}D_{a_2}^{a_3}\mathrm{}D_{a_{2K}}^{a_1},K=1,2,\mathrm{},r1,`$ (3.12)
$`\theta _r(C)`$ $`=`$ $`\{\begin{array}{cc}C_{a_1}^{a_2}D_{a_2}^{a_3}\mathrm{}D_{a_{2r}}^{a_1}\hfill & \text{for}n=2r+1\hfill \\ ϵ_{a_1b_1\mathrm{}a_rb_r}C^{a_1b_1}D^{a_2b_2}\mathrm{}D^{a_rb_r}\hfill & \text{for}n=2r\hfill \end{array}`$ (3.15)
where $`D_a{}_{}{}^{b}=C_a{}_{}{}^{c}C_{c}^{}^b`$.
Let us first consider the case $`l=n`$. Since we are interested in solutions of (3.1) in ghost number $`n`$, in (3.11) we necessarily have $`\omega _i(\theta _K(\widehat{C}^I),\widehat{B}^A)=k_i`$ where the $`k_i`$ are constants (the $`\widehat{\xi }^a`$ are of ghost number 1). In that case, no further condition is imposed on $`\alpha _n`$ and we have,
$$\alpha =\alpha _n=L(𝒯)\widehat{\mathrm{\Theta }},\delta _IL=0,$$
(3.16)
where $`\widehat{\mathrm{\Theta }}=\widehat{\xi }^0\mathrm{}\widehat{\xi }^{n1}`$. The above elements of $`H(s)`$ give rise in $`H_0^n(s|d,)`$ to the cocycles,
$$\alpha =L(𝒯)d^nx,\delta _IL=0,$$
(3.17)
which constitute the first set of consistent interactions announced in the introduction. Since these vertices are strictly gauge invariant they do not require any modification in the gauge transformations of the theory.
We now turn our attention to the case $`l<n`$. To that end, we substitute the general form (3.11) into eq. (3.9). By doing so, one easily proves that $`\alpha _l^i(\widehat{\xi }^a,𝒯^r)`$ has to obey the following equation,
$$s_1\alpha _l^i+\delta \alpha _{l+1,1}^i=0,$$
(3.18)
where $`\alpha _{l+1,1}^i=\alpha _{l+1,1}^i(\widehat{\xi }^a,𝒯^r,𝒯_{\overline{r}}^{})`$ is of antighost number $`1`$ and satisfies $`\delta _I\alpha _{l+1,1}^i(\widehat{\xi }^a,𝒯^r,𝒯_{\overline{r}}^{})=0`$. Trivial solutions of (3.18) ($`\alpha _l^i=s_1\beta _{l1}^i+\delta \beta _l^i`$) are irrelevant since they amount to trivial contributions in $`\alpha `$.
Eq. (3.18) along with its coboundary condition define the so-called “invariant characteristic cohomology” $`H_{char}^{inv}(d)`$ (with $`d`$ formally substituted by $`s_1`$) which plays a central rôle in the analysis of any local field theory. Theorems concerning $`H_{char}^{inv}(d)`$ in the cases of pure gravity and $`p`$-form gauge theory can be found respectively in and . The extension of those results in the case of gravity coupled to a system of $`2`$-forms is fully treated in . Here, we may restrict our attention to solutions of eq. (3.18) in ghost number $`n2`$. Indeed, taking into account (3.12) and (3.15) we see that the $`\theta _K(\widehat{C}^I)`$ and $`\widehat{B}^A`$ are at least of ghost number $`2`$; therefore, to construct solutions of (3.1) in ghost number $`n`$, we necessarily have $`ln2`$. In that case, the most general solution of (3.18) is up to trivial terms given by ,
$$\alpha _l^i=P^i(f_K,\widehat{H}^A)+\delta _l^{n3}k_A^i\overline{\widehat{H}}^A+\delta _l^{n2}\delta _4^nk_{AB}^i\overline{\widehat{H}}^A\overline{\widehat{H}}^B.$$
(3.19)
In (3.19), the $`k_A^i`$ and $`k_{AB}^i=k_{BA}^i`$ are constants, $`\overline{\widehat{H}}^A`$ is the Hodge dual of $`\widehat{H}^A`$ and the $`f_K`$ are generators for the the $`𝒢_L`$-invariant polynomials in the $`\widehat{R}^I`$. The $`f_K`$ are given in the cases $`n=2r`$ and $`n=2r+1`$ by:
$`f_K`$ $`=`$ $`\widehat{R}_{a_1}^{a_2}\widehat{R}_{a_2}^{a_3}\mathrm{}\widehat{R}_{a_{2K}}^{a_1},K=1,2,\mathrm{},r1,`$ (3.20)
$`f_r`$ $`=`$ $`\{\begin{array}{cc}\widehat{R}_{a_1}^{a_2}\widehat{R}_{a_2}^{a_3}\mathrm{}\widehat{R}_{a_{2r}}^{a_1}\hfill & \text{for}n=2r+1\hfill \\ ϵ_{a_1b_1\mathrm{}a_rb_r}\widehat{R}^{a_1b_1}\mathrm{}\widehat{R}^{a_rb_r}\hfill & \text{for}n=2r.\hfill \end{array}`$ (3.23)
From (3.19) we see that $`\alpha _l^i`$ involves two kinds of contributions. The first consists of the polynomials $`P^i(f_K,\widehat{H}^a)`$ which obey (3.18) without the need for a term $`\alpha _{l+1,1}^i`$. One says that the corresponding $`\alpha _l^i`$ are strongly $`s_1`$-closed. On the other hand, the last two terms of (3.19) are only weakly $`s_1`$-closed since they require a term $`\alpha _{l+1,1}^i`$ in order to satisfy eq. (3.18).
Let us first consider the BRST cocycles generated by the strongly $`s_1`$-closed $`\alpha _l^i`$. Their part of lowest $`\widehat{\xi }^a`$-degree is of the form,
$$\alpha _l=P^i(f_K,\widehat{H}^A)\omega _i(\theta _K(\widehat{C}^I),\widehat{B}^A).$$
(3.24)
We need to complete these $`\alpha _l`$ by terms of higher $`\widehat{\xi }^a`$-degree in order to obtain BRST cocycles. To do this, we associate to the $`\theta _K`$ the following quantities :
$$q_K=\underset{k=0}{\overset{2K1}{}}()^k\frac{(2K)!(2K1)!}{(2K+k)!(2Kk1)!}Str(\widehat{𝒞}\widehat{𝒟}^k\widehat{}^{2Kk1}),$$
(3.25)
with,
$$\widehat{𝒞}=\widehat{C}^IT_I,\widehat{𝒟}=\widehat{𝒞}^2,\widehat{}=\widehat{R}^IT_I.$$
(3.26)
In (3.26), the $`T_I`$ are the matrices of the adjoint representation of $`so(n1,1)`$ except for $`n=2r`$ in which the spinor representation is used. At lowest $`\widehat{\xi }^a`$-degree the $`q_K`$ begin with $`\theta _K`$ and are such that $`sq_K=f_K`$.
The usefulness of the $`q_K`$ stems from the fact that if $`\alpha `$ is a cocycle with $`\alpha _l`$ as in (3.24), then one may assume that $`\alpha =\alpha (f_K,q_K,\widehat{H}^A,\widehat{B}^A)=P^i(f_K,\widehat{H}^A)\omega _i(q_K,\widehat{B}^A)`$, i.e., $`\alpha `$ is obtained from $`\alpha _l`$ by replacing the $`\theta _K`$ by the $`q_K`$ . Furthermore, one can show that if $`\alpha =\alpha (f_K,q_K,\widehat{H}^A,\widehat{B}^A)=s\beta `$ then $`\beta =\beta (f_K,q_K,\widehat{H}^A,\widehat{B}^A)`$. The analysis of the BRST cocycles arising from (3.24) is therefore reduced to the investigation of the BRST cohomology in the so-called small algebra generated by the variable $`q_K,f_K,\widehat{H}^A,\widehat{B}^A`$.
This problem has been studied in detail for the pure gravitational case in . It is shown that in ghost number $`n`$ the BRST cocycles in the small algebra are of the “Chern-Simons” form, $`\alpha =q_KP(f_K)`$. Here, the only difference comes from the presence of $`\widehat{H}^A`$ and $`\widehat{B}^A`$ among the generators of the small algebra; they are related by $`s\widehat{B}^A=\widehat{H}^A`$. The procedure is nearly identical to and one can show that in ghost number $`n`$ the BRST cocycles are again of the Chern-Simons form, i.e.,
$$\alpha =q_KP_K(f_L,\widehat{H}^A)+\widehat{B}^AQ_A(f_K,\widehat{H}^B).$$
(3.27)
The corresponding consistent deformations are obtained from (3.27) by the substitution $`\widehat{\xi }^adx^\mu `$. They constitute the second set of vertices announced in the introduction. Note that they are invariant up to a boundary term and therefore do not require a modification of the gauge transformations.
We now consider the BRST cocycles generated by the weakly $`s_1`$-closed $`\alpha _l^i`$. Taking into account (3.12), (3.15) and (3.19), we see that in spacetime dimension $`n>4`$, the only possibility we have in ghost number $`n`$ for $`\alpha _l`$ is,
$$\alpha _{n3}=k_A^1\overline{\widehat{H}}^A\theta _1,$$
(3.28)
with $`\theta _1=C_a^bC_b^cC_c^a`$. $`\alpha _{n3}`$ has to be completed with terms of higher $`\widehat{\xi }^a`$-degree in order to obtain a BRST cocycle. To that end we introduce the following notation :
$$q^A=\overline{\widehat{H}}^A+\overline{\widehat{B}}_1^A+\overline{\widehat{B}}_2^A+\overline{\widehat{B}}_3^A,$$
(3.29)
where the $`\overline{\widehat{B}}_j^A`$ are defined through, $`\delta \overline{\widehat{B}}_1^A+d\overline{\widehat{H}}^A=0`$, $`\delta \overline{\widehat{B}}_2^A+d\overline{\widehat{B}}_1^A=0`$, $`\delta \overline{\widehat{B}}_3^A+d\overline{\widehat{B}}_2^A=0`$ ($`\overline{\widehat{B}}_j^A`$ is proportional to the dual of the antifield of antighost $`j`$, with the $`\widehat{\xi }^a`$ as differentials).
Using the above definition, it is straightforward to complete (3.28) to a BRST cocycle. Indeed, as can be seen by direct substitution,
$`\alpha `$ $`=`$ $`k_A^1q^Aq_1`$ (3.30)
$`=`$ $`k_A^1\overline{\widehat{H}}^A\theta _1+\alpha _{n2}+\alpha _{n1}+\alpha _n.`$ (3.31)
is a solution of (3.28). The corresponding elements of $`H(s|d,)`$ are again obtained from (3.30) by the substitution $`\widehat{\xi }^adx^\mu `$. Their components of antighost number zero give rise to the following class of consistent vertices:
$$V=k_A^1\overline{H}^ATr(\omega R\frac{1}{3}\omega ^3)$$
(3.32)
where $`\overline{H}^A`$ is the dual of the field strength $`H^A`$, $`\omega =\omega _\mu ^Idx^\mu T_I`$ and $`R=\frac{1}{2}dx^\mu dx^\nu R_{\mu \nu }^IT_I`$ (the $`T_I`$ are the matrices of the adjoint representation of $`so(n1,1)`$). These vertices require a modification in the gauge transformations since they are not invariant under the original ones. They belong to the third type of interactions described in the introduction.
In the particular case of spacetime dimension $`n=4`$ further couplings are possible. Indeed, in ghost number $`n=4`$ we have the following candidates for $`\alpha _l`$,
$$\alpha _l=k_A^1\overline{\widehat{H}}^A\theta _1+k_A^2\overline{\widehat{H}}^A\theta _2+k_{ABC}\overline{\widehat{H}}^A\overline{\widehat{H}}^B\widehat{B}^C,$$
(3.33)
where $`k_{ABC}=k_{BAC}`$ and $`\theta _2=ϵ_{abcd}C^{ab}C^{ce}C_e^d`$. The above $`\alpha _l`$ are easily completed to the following BRST-cocycles,
$$\alpha =k_A^1q^Aq_1+k_A^2q^Aq_2+k_{ABC}q^Aq^B\widehat{B}^A.$$
(3.34)
The first two terms yield consistent vertices identical to (3.32) (with the trace taken once in the adjoint representation and once in the spinorial representation of $`so(n1,1)`$):
$$V=k_A^1\overline{H}^ATr_{adj}(\omega R\frac{1}{3}\omega ^3)+k_A^2\overline{H}^ATr_{sp}(\omega R\frac{1}{3}\omega ^3)$$
(3.35)
The last term in (3.34) produces the Freedman-Townsend coupling,
$$V=k_{ABC}\overline{H}^A\overline{H}^BB^C.$$
(3.36)
This vertex is again not invariant under the original gauge transformations and a modification of these transformations is imposed.
## 4 Comments
In this article we have computed the local BRST cohomology $`H(s|d)`$ in ghost number $`0`$ and form degree $`n`$ in order to obtain all the first-order vertices of gravity coupled to a system of $`2`$-form gauge fields.
The first two types of couplings we obtain (strictly invariant polynomials in the curvatures and Chern-Simons forms) are consistent to all orders in the coupling constant since they are gauge invariant (up to a total derivative for the Chern-Simons forms) under the original gauge transformations. These vertices correspond to antifield independent representative of the BRST cohomology.
The last two types of vertices we describe are not invariant under the original gauge transformations since they depend non-trivially on the antifields. In that case a modification of the gauge transformations is required as well as the addition of higher-order vertices in order to obtain a theory consistent to all orders in the coupling constant. For the exotic couplings of the 2-form gauge fields to gravitational Chern-Simons forms, the full theory is of the Chapline-Manton type in which the 2-form curvature present in (2.3) is replaced by $`H_A^{}=H_A+k_ATr(\omega F\frac{1}{3}\omega ^3)`$. The Freedman-Townsend coupling gives rise to a non-polynomial full theory. Polynomiality can be restored by describing the theory with a “first-order formulation” and the full lagrangian is simply a covariantized version of the original Freedman-Townsend theory.
Our analysis can be extended to cover the couplings of gravity to $`p`$-forms .
## Acknowledgements
We are grateful to Marc Henneaux for suggesting the problem. We also thank him along with Glenn Barnich and Friedemann Brandt for useful discussions. This work is supported in part by the “Actions de Recherche Concertées” of the “Direction de la Recherche Scientifique - Communauté Française de Belgique” and by IISN - Belgium (convention 4.4505.86).
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# The Equilibrium Structure of Cosmological Halos: From Dwarf Galaxies to X-ray Clusters
(13 March 2000)
Abstract
A new model for the postcollapse equilibrium structure of virialized objects which condense out of the cosmological background universe is described and compared with observations and simulations of cosmological halos from dwarf galaxies to X-ray clusters. The model is based upon the assumption that virialized halos are isothermal, which leads to a prediction of a unique nonsingular isothermal sphere for the equilibrium structure, with a core radius which is approximately 1/30 times the size and a core density which is proportional to the mean background density at the epoch of collapse. These predicted nonsingular isothermal spheres are in good agreement with the observations of the internal structure of dark-matter-dominated halos from dwarf galaxies to X-ray clusters.
Poster paper presented at the Institute for Theoretical Physics Conference on Galaxy Formation and Evolution, March 14-17, 2000, at the University of California at Santa Barbara.
Equilibrium structure of virialized halos
\[Shapiro, Iliev, and Raga 1999, MNRAS, 307, 203 (Einstein-de Sitter case); Iliev and Shapiro 2000 (low-density, open and flat with cosmological constant)\]
* Problem and Motivation
+ Question: What equilibrium structure forms when a density perturbation collapses out of the expanding background universe and virializes?
+ An analytical model for the structure (e.g. mass profile, temperature, velocity dispersion, radius) of virialized cosmological halos would be a valuable tool for the semi-analytical modeling of galaxy and cluster formation in a hierarchical clustering model like CDM.
- Earlier work adopted crude approximations involving either uniform spheres or singular isothermal spheres which resulted from top-hat perturbation collapse
- What is a more realistic outcome, even for the simple top-hat problem?
+ N-body simulations of CDM predict dark matter halo profiles with singular density profiles, but a finite density core is required to explain:
- Dwarf galaxy rotation curves
- Cluster mass profiles inferred from gravitational lensing
+ As a result, the cold, collisionless nature of CDM has recently been re-examined to allow for variations which affect the post-collapse equilibrium structure of halos.
+ Suppose we ignore the details of this relaxation process and adopt the assumption that the final equilibrium is isothermal.
+ Solve this problem and compare the result with dwarf galaxy rotation curves and X-ray cluster data.
* Model:
+ Top-hat density perturbation collapses and virializes
+ Virialization leads to a truncated isothermal sphere in hydrostatic equilibrium (TIS) $``$ solution of the Lane-Emden equation (modified for $`\mathrm{\Lambda }0`$)
+ Total energy of top-hat is conserved thru collapse and virialization
+ Postcollapse temperature set by virial theorem (including effect of finite boundary pressure)
* Is the solution uniquely determined? – No, some additional information is required:
+ Minimum-Energy Solution: Boundary pressure is that for which the conserved top-hat energy is the minimum possible for an isothermal sphere of fixed mass within a finite truncation radius.
+ The Self-Similar Spherical Cosmological Infall Solution (Bertschinger 1985) confirms this choice if we identify the virialized object with the spherical region of post-shock gas and shell-crossing dark matter $``$ explains dynamical origin of boundary pressure adopted above as the result of thermalizing the energy of infall.
Summary of the TIS Solution
* Top-hat perturbation $``$ unique, nonsingular TIS (minimum-energy configuration)
$``$ universal, self-similar density profile for the postcollapse equilibrium of cosmic structure
+ Unique scale and amplitude set by top-hat mass and collapse epoch
+ Same density profile for gas and dark matter (no bias)
* Matter-Dominated Case (see Table 1 and Fig. 1)
+ Finite core size: $`r_0=0.034\times `$ radius $`r_t`$
+ Central density: $`\rho _0=514\times `$ surface density $`\rho _t`$
+ $`T=2.16T_{\mathrm{uniform}\mathrm{sphere}}=0.72T_{\mathrm{singular}\mathrm{isothermal}\mathrm{sphere}}`$
+ At intermediate radii, $`\rho `$ drops faster than $`r^2`$
* Flat, $`\mathrm{\Lambda }0`$ Case
+ Profile varies with epoch of collapse, approaching case I above for early collapse.
For example: for $`\mathrm{\Omega }_0=1\lambda _0=0.3`$ for $`z_{coll}=(0;0.5;1)`$:
- $`r_t/r_0=(30.04;29.68;29.54)`$
- $`\rho _0/\rho _t=(529.9;520.8;517.2)`$
- $`T/T_{\mathrm{uniform}\mathrm{sphere}}=(2.188;2.170;2.163)`$
Table 1: The Postcollapse Virial Equilibrium Resulting
from the Collapse of Top-Hat Density Perturbations:
Einstein-de Sitter Universe
| | Uniform | Singular | Our |
| --- | --- | --- | --- |
| | Sphere | Isothermal | Solution |
| | | Sphere | |
| $`\eta =\frac{r_t}{r_m}`$….. | 0.5 | 0.417 | 0.554 |
| $`{\displaystyle \frac{k_BT_{\mathrm{vir}}}{\left(\frac{2}{5}\frac{GMm}{r_{\mathrm{vir}}}\right)}}`$ | 1 | 3 | 2.16 |
| $`{\displaystyle \frac{\rho _0}{\rho _t}}`$……….. | 1 | $`\mathrm{}`$ | 514 |
| $`{\displaystyle \frac{\rho }{\rho _t}}`$……….. | 1 | 3 | 3.73 |
| $`{\displaystyle \frac{r_t}{r_0}}\mathrm{}\mathrm{}\mathrm{}\mathrm{}..`$ | – NA – | $`\mathrm{}`$ | 29.4 |
| $`{\displaystyle \frac{\rho }{\rho _b\left(t_{\mathrm{coll}}\right)}}`$….. | $`18\pi ^2`$ | $`18\pi ^2\left({\displaystyle \frac{6}{5}}\right)^3\pi ^5`$ | |
| | $`178`$ | $`307`$ | 130.5 |
Our solution = minimum-energy, truncated, nonsingular, isothermal sphere
Note: $`\rho _b`$ cosmic mean matter density
Density Profile of Halo which Forms from Top-Hat
Perturbation Collapse and Virialization
Fig.1: Density profile of truncated isothermal sphere which forms from the virialization of a top-hat density perturbation in a matter-dominated universe. Radius $`r`$ is in units of $`r_m`$ \- the top-hat radius at maximum expansion, while density $`\rho `$ is in terms of the density $`\rho _{SUS}`$ of the standard uniform sphere approximation for the virialized, post-collapse top-hat. (TIS = our solution, SUS = uniform sphere, SIS = singular isothermal sphere). Bottom panel shows logarithmic slope of density profile.
Direct Comparison with NFW Profile
Fig. 2: Continuous line = TIS profile; Dashed lines = “NFW” = Navarro, Frenk, and White (1996, 1997) profile:
$$\rho \left(r\right)=\frac{\delta _c\rho _{b0}}{cx\left(cx+1\right)^2},x=\frac{r}{r_{200}}$$
Range of $`c`$ appropriate for X-ray clusters to early forming dwarf galaxies.
Dwarf Galaxy Rotation Curves
Q: How well does our TIS profile match the observed mass profiles of dark-matter-dominated dwarf galaxies? The observed rotation curves of dwarf galaxies can be fit according to the following density profile with a finite density core (Burkert 1995):
$$\rho \left(r\right)=\frac{\rho _{0,Burkert}}{\left(r/r_c+1\right)\left(r^2/r_c^2+1\right)}$$
A: The TIS profile gives a nearly perfect fit to the Burkert profile. (see Fig. 3)
Fig. 3: Rotation Curve Fit. Best fit parameters:
$$\frac{\rho _{0,Burkert}}{\rho _{0,TIS}}=1.216,\frac{r_c}{r_{0,TIS}}=3.134$$
Solid line = Best fit TIS; Dashed line = Burkert profile, where $`\sigma _{TIS}^2=v^2/3=k_BT/m.`$
Q: How well does this best fit TIS profile predict the $`r_{max}`$ and $`v_{max}`$?
A: $`{\displaystyle \frac{r_{max,Burkert}}{r_{max,TIS}}}=1.13,{\displaystyle \frac{v_{max,Burkert}}{v_{max,TIS}}}=1.01`$
(i.e. excellent agreement)
The $`v_{max}r_{max}`$ relation for dwarf and LSB galaxies.
Q: Can the TIS halo model explain the observed correlation of $`v_{max}`$ and $`r_{max}`$ for dwarf spiral and LSB galaxies?
A: Yes, when the TIS halo model is combined with the Press-Schechter model which predicts the typical collapse epoch for objects of a given mass (i.e. the mass of the $`1\sigma `$-fluctuations vs. $`z_{coll}`$). (See Fig. 4) For the three untilted CDM models plotted, a cluster normalized Einstein-de Sitter model, and COBE-normalized low-density models ($`\mathrm{\Omega }_0=0.3`$ and $`\lambda _0=0`$ or 0.7), only the flat models yield a reasonable agreement with the observed $`v_{max}r_{max}`$ relation.
Fig. 4: Dwarf galaxies (triangles) and LSB galaxies (squares) from Kravtsov et al. (1998); Burkert: fit to data (Mori & Burkert 2000); SCDM: $`\mathrm{\Omega }_0=1`$, $`\lambda _0=0`$, $`\sigma _{8h^1}=0.5`$ (cluster normalized); OCDM: $`\mathrm{\Omega }_0=0.3`$, $`\lambda _0=0`$ (COBE normalized); $`\mathrm{\Lambda }`$CDM: $`\mathrm{\Omega }_0=0.3`$, $`\lambda _0=0.7`$, (COBE normalized); $`h=0.65`$ for all.
Galaxy Halo $`M\sigma _v`$ Relation
Q: How well does our TIS halo model scaling relation predict the velocity dispersion of galactic halos which form in the CDM model according to N-body simulations?
A: Antonuccio-Delogu, Becciani, & Pagliaro (1999) used an N-body treecode at high-res ($`256^3`$ particles) to simulate galactic halos in the region of a single and a double cluster. The agreement with the TIS model is good.
Fig. 6: Velocity dispersion vs. mass for galactic haloes in cluster regions: (upper panel) double cluster, (lower panel) single cluster.
X-Ray Cluster Scaling Relations
Q: How well does the TIS halo model predict the internal structure of X-ray clusters found by gas-dynamical/N-body simulations of X-ray cluster formation in the CDM model?
A: As shown below and in Fig. 5, our TIS model predictions agree astonishingly well with the mass-temperature and the radius-temperature virial relations and integrated mass profiles derived from numerical simulations by Evrard, Metzler and Navarro (1996). Apparently, these simulation results are not sensitive to the discrepancy between our prediction of finite density core and the N-body predictions of a density cusp for clusters in CDM.
* Mass Profile – Temperature Relation
$$r_Xr_{10}\left(X\right)\left(\frac{T}{10\mathrm{keV}}\right)^{1/2};X\frac{\rho \left(r\right)}{\rho _b}$$
Fig.5 (triangles) fit to CDM simulation results by Evrard, Metzler and Navarro (1996); (continuous line) TIS prediction.
* Mass-Temperature and Radius-Temperature Virial Relations
Evrard, Metzler and Navarro (1996)
$`M_{500}`$ $`=`$ $`\left(1.11\pm 0.16\right)\times 10^{15}\left({\displaystyle \frac{T}{10\mathrm{keV}}}\right)^{3/2}h^1M_{},`$
$`r_{500}`$ $`=`$ $`\left(1.24\pm 0.09\right)\left({\displaystyle \frac{T}{10\mathrm{keV}}}\right)^{1/2}h^1\mathrm{Mpc}.`$
$`M_{200}`$ $`=`$ $`1.45\times 10^{15}\left({\displaystyle \frac{T}{10\mathrm{keV}}}\right)^{3/2}h^1M_{},`$
$`r_{200}`$ $`=`$ $`1.85\left({\displaystyle \frac{T}{10\mathrm{keV}}}\right)^{1/2}h^1\mathrm{Mpc}.`$
Our solution
$`M_{500}`$ $`=`$ $`1.11\times 10^{15}\left({\displaystyle \frac{T}{10\mathrm{keV}}}\right)^{3/2}h^1M_{},`$
$`r_{500}`$ $`=`$ $`1.24\left({\displaystyle \frac{T}{10\mathrm{keV}}}\right)^{1/2}h^1\mathrm{Mpc}.`$
$`M_{200}`$ $`=`$ $`1.55\times 10^{15}\left({\displaystyle \frac{T}{10\mathrm{keV}}}\right)^{3/2}h^1M_{},`$
$`r_{200}`$ $`=`$ $`1.88\left({\displaystyle \frac{T}{10\mathrm{keV}}}\right)^{1/2}h^1\mathrm{Mpc}.`$
$`\beta `$-fits to X-ray Cluster Brightness and
Density Profiles
$$\rho _{gas}=\frac{\rho _0}{\left(1+r^2/r_c^2\right)^{3\beta /2}},I=\frac{I_0}{\left(1+\theta ^2/\theta _c^2\right)^{3\beta 1/2}}$$
Q: How well does the TIS model for the internal structure of X-ray clusters predict the observed and simulated X-ray brightness profile of clusters?
A: It predicts gas density profiles and brightness profiles which are well-fit by a $`\beta `$profile, with $`\beta `$-values for the TIS $`\beta `$-fit which are close to those of simulated clusters in the CDM model, but somewhat larger than the conventional observational result that $`\beta 2/3`$. However, recent X-ray results suggest that the true $`\beta `$-values are larger than 2/3 when measurements at larger radii are used and when central cooling flows are excluded from the fit.
| Brightness profile observations | $`\beta `$ |
| --- | --- |
| Jones and Foreman (1999) | 0.4-0.8, ave. 0.6 |
| Jones and Foreman (1992) | $`2/3`$ |
| Balland and Blanchard (1997) | 0.57 (Perseus) |
| | 0.75 (Coma) |
| Durret et al. (2000) | 0.53 (incl. cooling flow) |
| | 0.82 (excl. cooling flow) |
| Vikhlinin, Forman, & Jones (1999) | 0.7-0.8 |
| (fit by Henry 2000) | |
| TIS $`\beta `$-fit ($`r_c/r_{0,TIS}=2.639`$) | 0.904 |
| Gas density profile simulations | $`\beta `$ |
| Metzler and Evrard (1997) | 0.826 (DM) |
| | 0.870 (gas) |
| Eke, Navarro, and Frenk (1998) | 0.82 |
| Lewis et al. (1999) (adiabatic) | $`1`$ |
| Takizawa and Mineshige (1998) | $`0.9`$ |
| Navarro, Frenk, and White (1995) | 0.8 |
| TIS $`\beta `$-fit ($`r_c/r_{0,TIS}=2.416`$) | 0.846 |
X-ray Cluster Gas Entropy
Q: Can the TIS model for the internal structure of X-ray clusters explain the observed correlation between the gas entropy near the cluster center and cluster virial temperature?
A: Yes, but only for high-T clusters (i.e. $`T>`$ few keV) for which energy release feedback effects were probably not big enough to alter the entropy of the equilibrium halo. (See Fig. 7)
$$S=T/n_e^{2/3},\text{at }r=0.1r_{vir}$$
Fig. 7: Cluster entropy vs. temperature. Data: Ponman, Cannon, and Navarro (1999) Nature, 397; Error bars: T – 90% confidence level, entropy – span of variation from $`r=0.05r_{vir}`$ to $`r=0.2r_{vir}`$. Our solution: thick line = S at $`r=0.1r_{vir}`$, dashed lines = S at $`r=0.05r_{vir}`$ (lower), and $`0.2r_{vir}`$ (upper).
Cluster Mass Profiles Deduced from Strong Gravitational Lensing
Q: Can the TIS halo model explain the mass profile with a finite density core measured by Tyson, Kochanski, and Dell’Antonio (1998) for cluster CL 0024+1654 at $`z=0.39`$ using the strong gravitational lensing of background galaxies by the cluster to infer the cluster mass distribution?
A: Yes, the TIS model not only provides a good fit to the shape of the projected surface mass density distribution of this cluster (see Fig. 8) within the arcs, but when we match the central value as well as the shape, our model predicts the overall mass, and a cluster velocity dispersion in close agreement with the value $`\sigma _v=1200`$ km/s measured by Dressler and Gunn (1992).
Fig. 8: Projected surface density of cluster CL 0024+1654 inferred from lensing measurements, together with the best fit TIS model.
Summary
* TIS profile fits dwarf galaxy rotation curves; combined with Press-Schechter formalism matches results for observed $`v_{max}r_{max}`$ relation for dwarf galaxies
* Predicted mass-velocity dispersion relation agrees with high resolution N-body simulations of galactic halo formation by Antonuccio-Delogu et al. (1999)
* Predicted mass-radius-temperature scaling relations match simulation results from X-ray clusters in CDM model
+ Our solution derives empirical fitting formulae of Evrard, Metzler and Navarro (1996)
+ Agrees well with X-ray cluster observations at $`z=0`$
* Fits high temperature X-ray cluster entropy data
* X-ray brightness profile is predicted to match $`\beta `$-fit with $`\beta 0.9`$, larger than typically observed, but similar to results of gas-dynamical/N-body simulations of X-ray clusters in CDM model
* Fits the cluster mass profile with finite core derived from strong gravitational lensing data of Tyson et al. (1999) on CL 0024+1654
* Predicted mass profile is close to NFW profile for low values of concentration parameter, outside the core
This work was supported in part by NASA Grants NAG5-2785, NAG5-7363, and NAG5-7821, NSF Grant ASC-9504046, and Texas Advanced Research Program Grant 3658-0624-1999, and benefitted from PRS’ participation at the Aspen Center for Physics in summer 1998 and 1999.
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# Acoustic coupling between two air bubbles in water
## 1. Introduction
Bubbles play an important role in the sound propagation in everyday life liquids. For example, the murmur of the brooks essentially originates, as first suggested by Bragg, in the oscillations of air bubbles captured and dragged along by the water. The so-called “hot chocolate effect”, namely the rising of sound pitch when one repeatedly taps the bottom of the mug in which some instant coffee or chocolate is being dissolved, is explained by the releasing into the water of the tiny bubbles trapped in the powder. Bubble dynamics and acoustic properties of liquids containing a large number of bubbles have been widely studied for a long time. A very good and complete review in this domain was achieved by Leighton. Inter alia, the problem of the interaction of two neighbouring bubbles has been discussed using fluid dynamics tools or the acoustic-electrostatic analogy. Moreover, the free oscillations of a system of two (and even three) air cavities formed in a metal plate lying on a water surface have been theoretically and experimentally investigated in detail (including cubic nonlinearities). The aim of the present article is to present a simple, readily reproducible, experimental study of the forced oscillation regime of a two-air bubble system in water. We begin with a short introductory theory in which we show that the two-bubble system is mostly equivalent to a set of two magnetically coupled electric circuits.
## 2. Theoretical model
An air bubble in water will be considered as a perfect sphere<sup>1</sup><sup>1</sup>1 The correction to the Minnaert angular frequency due to deviation from the spherical shape can be shown to be negligible. of radius $`R(t)=R_0+\xi (t)`$, with variation $`\xi `$ much smaller than equilibrium value $`R_0`$. It can be shown that $`\xi (t)`$ oscillates with Minnaert’s angular frequency $`\omega _0=\sqrt{3\gamma P_0/\rho _0R_0^2}`$, where $`\gamma `$ is the specific heat ratio $`C_p/C_v`$ of air, and $`P_0`$ and $`\rho _0`$ respectively stand for the equilibrium pressure<sup>2</sup><sup>2</sup>2 The pressure difference accross the bubble boundary due to air-water surface tension is about $`1\%P_0`$ for a typical radius of $`1mm`$ (see eq. (65b)) and will be neglected: $`P_0`$ is also the equilibrium pressure of enclosed air. and mass density of water. This oscillation is damped through several mechanisms: of course the acoustic radiation damping (thanks to which the bubble noise is audible), but also the viscous and thermal dampings. We will neglect, in the following simplified theory, the last two ones. Moreover, allowing for the typical $`1KHz`$ acoustic frequency and $`1mm`$ bubble size we deal with in our experiment, we will neglect any sound propagation in the enclosed air. We thus deliberately restrict the present study to the (radial) fundamental resonance of the air bubble-water system.
### 2.1 One-bubble free oscillation
Let us consider one bubble with radius $`R_0`$ immersed in an infinite volume of water at equilibrium pressure $`P_0`$. Let $`P(\stackrel{}{r},t)`$ be the actual pressure at site $`\stackrel{}{r}`$ and instant $`t`$. The extra pressure $`p(\stackrel{}{r},t)`$ is defined as $`P(\stackrel{}{r},t)P_0`$. According to Minnaert’s assumption, the enclosed air undergoes isentropic transformations and its (extra) pressure $`p(t)`$ is homogeneous inside the bubble. Then, neglecting air’s inertia as well as the air-water surface tension, $`p(t)`$ and the radius variation $`\xi (t)`$ are linked by:
$$\frac{p(t)}{P_0}+\frac{3\gamma \xi (t)}{R_0}=0$$
(1)
On the other hand, it can be easily shown that (extra) pressure $`p(r,t)`$ at distance $`r`$ from the center of the bubble follows a d’Alembert-like 1D equation, the solution of which exactly reads, for $`rR_0`$:
$$p(r,t)=\frac{1}{r}\rho _0R_0^2\left[\xi ^{^{\prime \prime }}\frac{R_0}{c}\xi ^{^{\prime \prime \prime }}+\mathrm{}+(\frac{R_0}{c})^k\xi ^{(2+k)}+\mathrm{}\right](t\frac{rR_0}{c})$$
(2)
where $`c`$ is the sound velocity in water. If the acoustic wavelength $`\lambda `$ is much larger than $`r`$ (i.e., under the circumstances, if condition $`r\omega _0/c1`$ is fulfilled), then $`p(r,t)`$ can be approximated by:
$`p(r,t)`$ $``$ $`\rho _0R_0^2\left[{\displaystyle \frac{\xi ^{^{\prime \prime }}(t)}{r}}{\displaystyle \frac{\xi ^{^{\prime \prime \prime }}(t)}{c}}\right]`$ (3)
$``$ $`\rho _0R_0^2\left[{\displaystyle \frac{\xi ^{^{\prime \prime }}(t)}{r}}+{\displaystyle \frac{\omega _0^2}{c}}\xi ^{^{}}(t)\right]`$ (4)
Then, equalling $`p(t)`$ in eq. (1) to $`p(R_0,t)`$ in eq. (4), one gets, all calculations carried out:
$$\xi ^{^{\prime \prime }}+\frac{\omega _0^2R_0}{c}\xi ^{^{}}+\frac{3\gamma P_0}{\rho _0R_0^2}\xi =\xi ^{^{\prime \prime }}+\mathrm{\Gamma }_{rad}\xi ^{^{}}+\omega _0^2\xi =0$$
(5)
which is the well-known differential equation of a weakly<sup>3</sup><sup>3</sup>3 Ratio $`\mathrm{\Gamma }_{rad}/\omega _0=\omega _0R_0/c`$ is actually assumed to be much smaller than unity, as a consequence of the the validity condition of eq. (4). damped 1D harmonic oscillator.
### 2.2 Two-bubble free oscillation
Let us now add a second bubble, with the same (equilibrium) radius $`R_0`$, at a distance $`d`$ apart from the first one. Let $`\stackrel{}{r_i}`$ ($`i=1,2`$) be the (equilibrium) position of the $`i^{th}`$ bubble center, $`\xi _i(t)`$ its radius variation, $`p_i(t)`$ the (inner) extra pressure of the enclosed air, and $`p_i(\stackrel{}{r},t)`$ (resp. $`\stackrel{}{u}_i(\stackrel{}{r},t)`$ the would-be (outer)extra pressure (resp. displacement with respect to equilibrium) at point $`\stackrel{}{r}`$ and instant $`t`$ in the water medium if bubble $`i`$ was alone. Then, allowing for the superposition principle for small displacements, we assume that overall water extra pressure and displacement respectively read:
$`p(\stackrel{}{r},t)`$ $`=`$ $`p_1(\stackrel{}{r},t)+p_2(\stackrel{}{r},t)`$ (6)
$`\stackrel{}{u}(\stackrel{}{r},t)`$ $`=`$ $`\stackrel{}{u}_1(\stackrel{}{r},t)+\stackrel{}{u}_2(\stackrel{}{r},t)`$ (7)
with, of course, $`p_i(t)`$ and $`\xi _i(t)`$ still linked by eq.(1). On the surface of the first bubble: $`r_1=\left|\stackrel{}{r}\stackrel{}{r}_1\right|=R_0`$, $`r_2=\left|\stackrel{}{r}\stackrel{}{r}_2\right|d`$, and $`p(\stackrel{}{r},t)=p_1(t)`$. A similar constraint is required on the surface of the second bubble, where $`r_1d`$ and $`r_2=R_0`$. If the bubble spacing $`d`$ is much smaller than $`\lambda `$ (i.e. $`\omega _0d/c1`$)<sup>4</sup><sup>4</sup>4 In our experiments $`\lambda `$ is of order $`1m`$, while $`d`$ ranges from $`1`$ to $`5cm`$., then eq. (4) is available and we finally get the following pair of coupled motion equations:
$`\xi _1^{^{\prime \prime }}+\alpha \xi _2^{^{\prime \prime }}+\mathrm{\Gamma }_{rad}(\xi _1^{^{}}+\xi _2^{^{}})+\omega _0^2\xi _1=0`$ (8)
$`\alpha \xi _1^{^{\prime \prime }}+\xi _2^{^{\prime \prime }}+\mathrm{\Gamma }_{rad}(\xi _1^{^{}}+\xi _2^{^{}})+\omega _0^2\xi _2=0`$ (9)
where $`\alpha =R_0/d(<0.5)`$ can be regarded as a dimensionless coupling constant. Observe, by the way, that if double condition: $`R_0d\lambda `$ is fulfilled, eqs. (8) and (9) are available (with $`\alpha 0`$), and dynamic variables $`\xi _i`$ are still coupled by radiation damping, since the dissipation terms do not involve $`\alpha `$.
Defining symmetrical and antisymmetrical normal variables $`\varphi _s(t)`$ and $`\varphi _a(t)`$ as respectively the sum and the difference of $`\xi _1(t)`$ and $`\xi _2(t)`$, we get the uncoupled equations system:
$`(1+\alpha )\varphi _s^{^{\prime \prime }}+2\mathrm{\Gamma }_{rad}\varphi _s^{^{}}+\omega _0^2\varphi _s`$ $`=`$ $`0`$ (10)
$`(1\alpha )\varphi _a^{^{\prime \prime }}+\omega _0^2\varphi _a`$ $`=`$ $`0`$ (11)
It is noteworthy that, as far as only radiation is concerned, the symmetrical mode’s damping rate is twice the single-bubble’s one, while the antisymmetrical mode is undamped.This feature is easily understood in terms of constructive (resp. destructive) interference between the acoustic waves radiated by each bubble, and parallels a well-known situation in the atomic physics domain (super- and sub- radiant quantum states of a couple of identical atoms interacting with each other through the E.M. field). From eqs. (10) and (11), it is clear that the symmetrical mode has the lower angular frequency $`\omega _s=\omega _0/\sqrt{1+\alpha }`$, and the antisymmetrical mode the higher one $`\omega _a=\omega _0/\sqrt{1\alpha }`$. Observe that, leaving apart the calculation of radiative damping, it is very easy to derive above expressions of $`\omega _{s,a}`$ using the following trick. Let us consider the water as an uncompressible fluid (i.e. $`c\mathrm{}`$). The water displacement due to bubble $`i`$’s motion simply reads:
$$\stackrel{}{u}_i(\stackrel{}{r},t)=\xi _i(t)\frac{R_0^2}{r_i^2}\stackrel{}{e}_{r_i}$$
(12)
with $`\stackrel{}{e}_{r_i}=(\stackrel{}{r}\stackrel{}{r}_i)/|\stackrel{}{r}\stackrel{}{r}_i|=(\stackrel{}{r}\stackrel{}{r}_i)/r_i`$. Then, allowing for eq. (7), the overall water kinetic energy $`T`$ is:
$`T`$ $`=`$ $`{\displaystyle \frac{1}{2}}\rho _0{\displaystyle d^3r}({\displaystyle \frac{\stackrel{}{u}}{t}})^2`$ (13)
$`=`$ $`{\displaystyle \frac{1}{2}}M_0\left(\xi _1^{}_{}{}^{}2+\xi _2^{}_{}{}^{}2+2\alpha \xi _1^{^{}}\xi _1^{^{}}\right)`$ (14)
where $`M_0=4\pi R_0^3\rho _0`$ is the effective mass of either bubble. On the other hand, the total potential energy $`V`$ associated with the isentropic compressibility of the enclosed air reads:
$$V=\frac{1}{2}K\left(\xi _1^2+\xi _2^2\right)$$
(15)
where $`K=12\pi \gamma R_0P_0`$ is the effective stiffness of either bubble. The Lagrange equations derived from $`L=TV`$ are:
$`\xi _1^{^{\prime \prime }}+\alpha \xi _2^{^{\prime \prime }}+\omega _0^2\xi _1=0`$ (16)
$`\alpha \xi _1^{^{\prime \prime }}+\xi _2^{^{\prime \prime }}+\omega _0^2\xi _2=0`$ (17)
which is exactly the $`c\mathrm{}`$ limit of eqs. (8) and (9). It is worth noticing that eqs. (14) through (17) are formally equivalent to those of a system of two ($`L,C`$) electric circuits coupled by mutual induction with coefficient $`\alpha L`$. In this analogy, $`M_0`$ and $`K`$ respectively correspond to $`L`$ and $`1/C`$, and the $`\xi _i`$’s to the electric charges $`q_i`$ of either capacitor.
### 2.3 Forced oscillation
Let us now suppose that the above studied two-bubble system is driven by an external acoustic source with an angular frequency $`\omega `$ near Minnaert’s one, $`\omega _0`$. The phase difference of the driving pressures on both bubbles can therefore be neglected, since $`\omega d/c1`$. Let $`p_{ei}(t)`$ be the external pressure undergone by bubble $`i`$. Motion eqs. (8) and (9) are then completed in:
$`\xi _1^{^{\prime \prime }}+\alpha \xi _2^{^{\prime \prime }}+\mathrm{\Gamma }_{rad}(\xi _1^{^{}}+\xi _2^{^{}})+\omega _0^2\xi _1={\displaystyle \frac{p_{e1}(t)}{\rho _0R_0}}`$ (18)
$`\alpha \xi _1^{^{\prime \prime }}+\xi _2^{^{\prime \prime }}+\mathrm{\Gamma }_{rad}(\xi _1^{^{}}+\xi _2^{^{}})+\omega _0^2\xi _2={\displaystyle \frac{p_{e2}(t)}{\rho _0R_0}}`$ (19)
or equivalently:
$`\varphi _s^{^{\prime \prime }}+{\displaystyle \frac{2\mathrm{\Gamma }_{rad}}{1+\alpha }}\varphi _s^{^{}}+\omega _s^2\varphi _s`$ $`=`$ $`F_{es}(t)`$ (20)
$`\varphi _a^{^{\prime \prime }}+\omega _a^2\varphi _a`$ $`=`$ $`F_{ea}(t)`$ (21)
with $`F_{es}(t)=(p_{e1}(t)+p_{e2}(t))/\rho _0R_0(1+\alpha )`$ and $`F_{ea}(t)=(p_{e1}(t)p_{e2}(t))/\rho _0R_0(1\alpha )`$. Solving for $`\varphi _s`$ and $`\varphi _a`$ in above eqs. (20) and (21), one gets $`\xi _1(t)`$ and $`\xi _2(t)`$, and consequently (using eq. (4)) quantities $`p_1(\stackrel{}{r},t)`$ and $`p_2(\stackrel{}{r},t)`$ at any point $`\stackrel{}{r}`$ of the medium. At last, comparing external (applied) pressure $`p_e(\stackrel{}{r},t)`$ with the actual overall extra pressure $`p(\stackrel{}{r},t)=p_e(\stackrel{}{r},t)+p_1(\stackrel{}{r},t)+p_2(\stackrel{}{r},t)`$, we can experimentally measure the two-bubble system’s response as a function of $`\omega `$. In this respect (and provided that the excitation-detection geometry allows it), resonances are expected for $`\omega =\omega _s`$ and $`\omega =\omega _a`$.
## 3. Experiments
Our aim is to demonstrate the existence of both above mentioned modes. From an experimental point of view, it turns out to be easier to implement a forced oscillation scheme than a free oscillation one. We therefore present the former hereafter.
### 3.1 Experimental setup
A small net (see fig.1), made up with a gauze maintained with a wire, is designed to catch up an air bubble in water and to fix it at any desired position without appreciably modifying acoustic impedance and spherical symmetry. Two such devices are used for studying the two-bubble system. The external driving source is a speaker and extrapressure $`p(\stackrel{}{r},t)`$ is measured with a small microphone. A function generator, to which the speaker is connected, produces a c.w. sinusoidal signal with a frequency slowly swept from $`f_{low}`$ to $`f_{high}`$. The signal delivered by the microphone is transmitted to a lock-in amplifier which compares it with the reference one (delivered by the function generator) and decomposes it into real and imaginary parts. Both parts can be seen on an oscilloscope and recorded with a computer (see fig.2).
In a preliminary set of experiments, without any bubble in the aquarium, the response of the microphone is calibrated for different speaker-microphone configurations. Two kinds of configurations are presented in figure 3. In figs. 3(a) and 3(b) the configuration is deliberately asymmetrical: the microphone is mainly susceptible to bubble 2’s motion, while the speaker selectively drives bubble 2 (fig. 3(a)) or bubble 1 (fig. 3(b)), so that $`F_{ea}(t)`$ is nonzero: both modes can be excited and the associated motions detected. In fig. 3(c), the speaker is placed far from the bubbles; then, not only the phases, but also the amplitudes of the external pressures $`p_{1e}(t)`$ and $`p_{2e}(t)`$ undergone on either bubble are appreciably equal. In such a symmetrical excitation configuration, $`F_{ea}(t)=0`$, so that the antisymmetrical mode remains unexcited. Observe, by the way, that since distances $`r_1`$ and $`r_2`$ between the bubbles and the microphone are equal, the latter would detect no contribution from the antisymmetrical mode even though it was excited (see eqs.(4) and (6): $`r_1=r_2`$ and $`\xi _1=\xi _2`$ yields $`p_1(r_1,t)+p(r_2,t)=0`$ ).
### 3.2 Results and discussion
In figs. 4(a) and 4(b), the imaginary part $`Imp`$ of the output signal from the lock-in amplifier is displayed versus the speaker frequency $`f`$ for various values of the bubbles spacing $`d`$. Figures 4(a) and 4(b) respectively correspond to configurations 3(a) and 3(b). Two resonances can be made out in fig 4(a) and (though at a lesser degree) in fig 4(b). Observe that the sign of the signal at the higher frequency resonance is changed from 4(a) to 4(b), while the lower frequency one remains unchanged. This is consistent with the latter signal being associated with the symmetrical mode’s resonance ($`\omega _s=\omega _0/\sqrt{1+\alpha }<\omega _0`$, and $`F_{es}`$ unchanged from configuration 3(a) to 3(b)), and the former one with the antisymmetrical mode’s resonance $`(\omega _a=\omega _0/\sqrt{1\alpha }>\omega _0`$, and $`F_{ea}`$ changed into $`F_{ea}`$ from configuration 3(a) to 3(b)).
It is noteworthy that both resonances have appreciably the same width. This is in contradiction with simplified eqs.(10) and (11) (or (20) and (21)), in which only the radiation damping was considered. In fact, as mentioned in introduction, other kinds of damping should be taken into account: if viscous damping is absolutely negligible for such large bubbles, thermal damping is not (see fig.8 in ), and may be at the origin of the linewidth. Further discussion of this point is out of the scope of the present paper. In figure 5, we have plotted, for both symmetrical and antisymmetrical modes, and for $`R_02mm`$, the inverse squared frequency $`f^2`$ (multiplied by a factor of $`10^7`$) versus the inverse bubble spacing $`d^1`$, in order to get a visual check of theoretical relations:
$`{\displaystyle \frac{1}{f_s^2}}={\displaystyle \frac{1}{f_0^2}}+{\displaystyle \frac{R_0}{f_0^2}}{\displaystyle \frac{1}{d}}`$ (22)
$`{\displaystyle \frac{1}{f_a^2}}={\displaystyle \frac{1}{f_0^2}}{\displaystyle \frac{R_0}{f_0^2}}{\displaystyle \frac{1}{d}}`$ (23)
Although experimental points are appreciably aligned, the measured slopes are about $`40\%`$ below theoretical prediction, suggesting that coupling constant $`\alpha `$ has been overestimated. In fact, theoretical value $`\alpha =R_0/d`$ was derived in eq.(14) when integrating the water kinetic energy density $`\frac{1}{2}\rho _0(\frac{\stackrel{}{u}}{t})^2`$ over the whole space<sup>5</sup><sup>5</sup>5 More precisely: over the whole space outside the two bubbles (the inner air’s inertia being negligible). Nevertheless, it can be shown that the coefficient of coupling term $`\xi _1^{^{}}\xi _2^{^{}}`$ in integral (14) does not depend on the bubbles radius $`R_0`$.. This inertial coupling is naturally lowered if some obstacle lies between the bubbles and consequently screens (part of) the water flow<sup>6</sup><sup>6</sup>6 The effective mass $`M_0`$ is modified too, but at a lesser degree.. Now, this is exactly what happens in configurations 3(a) and 3(b): to be able to excite the antisymmetrical mode, we are compelled to insert the speaker between the two bubbles, thus bringing about the above screening effect. In order to check this interpretation, we performed the same experiment with configuration 3(c), and recorded, for the symmetrical branch of the linear fitting of fig. 5, a slope of about $`90\%`$ of the theoretically predicted value.
At last, it should be noted that, when blocked from the top by the net, the bubble is, strictly speaking, no longer spherical. As mentioned in footnote 1, such a deviation from the spherical shape is (almost) of no consequence, and we have used the formulas derived above in section 2 with $`R_0`$ standing for the radius of the sphere of equivalent volume (i.e. the radius of the bubble before it is captured by the net). Nevertheless, this feature of our experimental setup raises the following difficulty: since, in course of motion, the fixed point of the bubble is no longer its centre (as implicitly assumed in the theoretical model) but its top, expression (12) of the water displacement is no longer correct; a dipolar term should be added to the spherical monopolar one. As a consequence, the kinetic energy $`T`$ derived in eq.(14) is modified too. More precisely, an exact calulation shows that $`T`$ should be multiplied by a factor of $`7/6`$ and the coupling constant $`\alpha `$ by a factor of $`1+R_0^2/(4d^2)`$. These corrections lie within our experimental accuracy. We have consequently neglected them. In this respect, it may be noted that the gauze in our device does not act like a rigid wall because the water can flow through the meshes of the net. The situation is therefore different from that discused in other studies considering the influence of the proximity of a rigid boundary on the Minnaert frequency.
## 4. Conclusion
As a conclusion, the acoustic inertial coupling between two air bubbles in water has been experimentally put in evidence. Theoretical analysis shows that the two-bubble system is formally equivalent to a set of two magnetically coupled ($`L`$,$`C`$) electric circuits, with two eigenmodes, respectively symmetrical and antisymmetrical. Experimental measurements and theoretical predictions are in $`10\%`$ accuracy agreement.
## 5. Acknowledgment
We gratefully thank Professor F.Massias for enriching discussions about the correction to the Minnaert frequency in the case of a bubble fixed from the top by a wide-mesh net.
## Figure Captions
Simple tool for capturing the bubble
Diagram of the experimental setup. In the experiment, we pump air into a tube immersed in water to produce the bubbles. The radii difference between these bubbles is small and will be neglected. (It can be shown that a small radii difference yields second order correction of $`\omega _s`$ and $`\omega _a`$).
Different geometrical configurations
(a) Spectra of symmetrical and antisymmetrical modes in configuration 3(a).
(b) Change of sign of $`Imp`$ for the antisymmetrical mode when configuration 3(b) is adopted.
Linear fitting of the plot $`\frac{1}{f_{s,a}^2}10^7(Hz^2)`$ vs. $`\frac{1}{d}(cm^1)`$ for the two modes. Average of the resonance frequency of the two bubbles: $`1499Hz`$; corresponding radius: $`0.217cm`$; slopes for the two fitting lines: $`0.578`$ and $`0.550`$ $`(cmsec^2)`$; slopes of theoretical prediction: $`\pm 0.966(cmsec^2)`$.
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# Running of the heavy quark production current and 1/|𝐤| potential in QCD
## I Introduction
For systems involving a heavy quark-antiquark pair near threshold it is useful to combine the QCD coupling constant expansion with an expansion in powers of the relative velocity $`v`$. This expansion is facilitated by using non-relativistic QCD formulated as an effective field theory with an explicit power counting in $`v`$ . For the potential between a heavy quark and antiquark this double expansion takes the form
$`V(𝐩,𝐩^{})={\displaystyle \underset{n=2}{\overset{\mathrm{}}{}}}V^{(n)},V^{(n)}={\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}V^{(n,j)},`$ (1)
$`\text{where}V^{(n)}v^n,V^{(n,j)}v^n\alpha _s^j.`$ (2)
The even terms $`V^{(2k)}`$ are first generated at tree-level (order $`\alpha _s`$), and the odd terms $`V^{(2k+1)}`$ are first generated at one loop (order $`\alpha _s^2`$).
Matrix elements for non-relativistic QCD systems typically depend on logarithms of the velocity $`v`$. For small $`v`$ it is convenient to sum large logarithms of the form $`\alpha _s(mv)\mathrm{ln}(v)`$ and $`\alpha _s(mv^2)\mathrm{ln}(v)`$ by using renormalization group equations in the effective theory. This reorganizes the series in $`j`$ in Eq. (1) so that:
$`V^{(n)}={\displaystyle \underset{j}{}}\stackrel{~}{V}^{(n,j)},\text{where}\stackrel{~}{V}^{(n,j)}v^n\alpha _s(m)^j{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}[\alpha _s(m)\mathrm{ln}(v)]^k.`$ (3)
The simplest example of such a summation is the use of a running coupling constant, $`\alpha _s(\mu )`$, instead of the coupling at the matching scale $`\alpha _s(m)`$. For $`\mu <m`$ the running coupling includes a series of $`\alpha _s(m)\mathrm{ln}(\mu /m)`$ terms. However, it should be emphasized that the complete set of renormalization group logarithms are not determined by the simple replacement $`\alpha _s(m)\alpha _s(\mu )`$.
A complication in summing the logarithms is the presence of two low energy scales: the typical momenta of the heavy quark $`mv`$ (soft scale), and typical energy $`mv^2`$ (ultrasoft scale). Two approaches have been proposed for dealing with the presence of two scales, one stage and two stage running. With two stage running, one matches QCD onto an effective theory at the scale $`\mu =m`$ and then runs to the scale $`\mu =mv`$. At $`\mu =mv`$ one matches onto an effective theory called pNRQCD which has composite $`q\overline{q}`$ fields, and then considers the running to $`\mu =mv^2`$. In vNRQCD , the running occurs in a single stage using a velocity renormalization group. The velocity renormalization group takes into account that the scales $`mv`$ and $`mv^2`$ are tied together by the heavy quark equation of motion for all $`\mu <m`$. We use dimensional regularization and the $`\overline{\mathrm{MS}}`$ scheme with both the ultrasoft scale $`\mu _Umv^2`$ and the soft scale $`\mu _Smv`$ given in terms of a single subtraction point velocity $`\nu `$: $`\mu _U=m\nu ^2`$ and $`\mu _S=m\nu `$. The renormalization group equations are written for the variable $`\nu `$. It is an interesting question as to whether both the one and two stage methods of running will sum the full set of $`\alpha _s\mathrm{ln}(v)`$ terms. In this paper only the single stage running will be considered.
The running of the Coulomb potential at one and two loops, $`\stackrel{~}{V}^{(2,1)}`$ and $`\stackrel{~}{V}^{(2,2)}`$, is determined by the running coupling constant $`\alpha _s(\mu )`$ At three loops ultrasoft gluons can contribute to the running of the static potential.. The summation of the leading logarithms for the the $`v^0`$ potential was carried out in Refs. . In this paper we extend this analysis to the $`1/|𝐤|`$ potential by calculating its two loop anomalous dimension.
If $`mv^2\mathrm{\Lambda }_{\mathrm{QCD}}`$, the $`v`$ expansion can be applied to non-relativistic QCD systems in a perturbative manner. This is the case for $`t\overline{t}`$ production near threshold where $`mv^24\mathrm{GeV}`$, and is the situation that will be analyzed in this paper. Of particular interest is the Coulombic regime where $`v\alpha _s`$. In this regime the expansion in Eq. (1) has the form
$`V`$ $`=`$ $`\left[\stackrel{~}{V}^{(2,1)}\right]+\left[\stackrel{~}{V}^{(2,2)}\right]+\left[\stackrel{~}{V}^{(2,3)}+\stackrel{~}{V}^{(1,2)}+\stackrel{~}{V}^{(0,1)}\right]+\mathrm{}`$ (4)
$``$ $`\left[{\displaystyle \frac{\alpha _s}{v^2}}\right]+\left[{\displaystyle \frac{\alpha _s^2}{v^2}}\right]+\left[{\displaystyle \frac{\alpha _s^3}{v^2}}+{\displaystyle \frac{\alpha _s^2}{v}}+\alpha _sv^0\right]+\mathrm{}.`$ (5)
If the $`\alpha _s\mathrm{ln}(v)`$ dependence is treated perturbatively these terms are referred to as the leading order (LO), next-to-leading order (NLO), and next-to-next-to-leading order (NNLO) potentials. Note that since the $`1/|𝐤|`$ potential first occurs at one loop, it only contributes at NNLO. When the series in $`\alpha _s\mathrm{ln}v`$ are summed, the terms in Eq. (4) will be referred to as leading-log (LL), next-to-leading log (NLL) and next-to-next-to-leading log (NNLL), respectively. In the Coulomb regime, the Coulomb potential must be kept to all orders. Each additional Coulomb insertion gives a $`\alpha _s/v^2`$ plus a factor of $`v`$ (from the potential loop), so each new Coulomb interaction costs a factor of $`\alpha _s/v1`$.
To study the threshold production of $`t\overline{t}`$, a non-relativistic expansion must also be made for the electromagnetic production current<sup>§</sup><sup>§</sup>§We will ignore effects associated with the top quark width.
$`\overline{t}\gamma ^it`$ $`=`$ $`{\displaystyle \underset{𝐩}{}}c_1\left(\psi _𝐩^{}𝝈^i\chi _𝐩^{}\right){\displaystyle \frac{c_2}{2m^2}}\left(\psi _𝐩^{}𝐩𝝈𝐩^𝐢\chi _𝐩^{}\right){\displaystyle \frac{c_3}{m^2}}\left(\psi _𝐩^{}𝐩^\mathrm{𝟐}𝝈^𝐢\chi _𝐩^{}\right)+\mathrm{}.`$ (6)
The fields $`\psi ^{}`$ and $`\chi ^{}`$ create non-relativistic top quarks and antiquarks respectively. The $`c_1`$ term contributes at order $`v^0`$, and the $`c_2`$ and $`c_3`$ terms contribute at order $`v^2`$. The current on the LHS of Eq. (6) is conserved, and has no anomalous dimension in QCD. However the non-relativistic current operators on the RHS are scale dependent in the effective theory, and the coefficients $`c_j`$ therefore depend on logarithms of $`\mu `$. The coefficients $`c_j`$ each have an expansion in $`\alpha _s`$. The matching at $`\mu m`$ is known to order $`\alpha _s^2`$ for $`c_1`$, and to order $`\alpha _s`$ for $`c_2`$ and $`c_3`$, so the production current is known to NNLO with partial N<sup>3</sup>LO results. At LL order, one needs the tree-level matching $`c_1=1`$ at $`\mu =m`$, and the $`v^2`$ coefficients $`c_2`$ and $`c_3`$ can be set to zero. There is no one loop anomalous dimension for the operator $`\psi _𝐩^{}\sigma ^i\chi _𝐩^{}`$, so the LL result is that $`c_1=1`$ at all $`\nu `$.
At NLL order, we need the one loop matching for $`c_1`$ at $`\mu =m`$, and the two loop running for $`c_1`$. The coefficients $`c_{2,3}`$ first enter at NNLL, at which order one would also need the three-loop anomalous dimension for $`c_1`$. At two loops, the anomalous dimension for $`c_1`$ was computed at the matching scale $`\mu =m`$
$`\mu {\displaystyle \frac{}{\mu }}c_1(\mu )|_{\mu =m}=C_F\left({\displaystyle \frac{1}{3}}C_F+{\displaystyle \frac{1}{2}}C_A\right)\alpha _s^2(m),`$ (7)
by studying the two loop matching condition for $`c_1`$. For $`\mu <m`$ the anomalous dimension no longer has the simple form Eq. (7), but depends on the running of the quark potential. This anomalous dimension was computed in Ref. , and depends on the running values $`[\stackrel{~}{V}^{(2,1)}]^2`$, $`\stackrel{~}{V}^{(2,1)}\times \stackrel{~}{V}^{(0,1)}`$, and $`\stackrel{~}{V}^{(1,2)}`$ (see Eq. (74) below). It is interesting that the structure of this result implies that determining the RHS of Eq. (7) for $`\mu <m`$ requires the LL values of $`\stackrel{~}{V}^{(2,1)}`$ and $`\stackrel{~}{V}^{(0,1)}`$, but the NLL value of $`\stackrel{~}{V}^{(1,2)}`$. The one loop running of $`\stackrel{~}{V}^{(2,1)}`$ is well known, and the running of $`\stackrel{~}{V}^{(2,2)}`$ was computed in Refs. . The two loop running of $`\stackrel{~}{V}^{(1,2)}`$ is computed in this paper. Using the running of these terms in the potential, we arrive at a complete NLL expression for $`c_1`$. The running of the non-relativistic scalar current is also briefly discussed.
In section II the effective theory is reviewed. We explain how reparameterization invariance fixes the value of the lowest order coupling of an ultrasoft gluon to the Coulomb potential. The details of the computation of the two loop anomalous dimension for the $`1/|𝐤|`$ potential are given in section III, and in the appendices. In section III we give a derivation of the anomalous dimension using on-shell potentials, while in appendix B we repeat the derivation in the presence of off-shell potentials. Readers not interested in the technical details can skip to section IV, where our results are discussed. In section IV we expand our renormalization group improved results in powers of $`\alpha _s`$ to compare to finite order calculations in the literature. For the color singlet $`1/|𝐤|`$ potential the first $`\alpha _s\mathrm{ln}(v)`$ term in the series was computed in Ref. , and our result for this term agrees with theirs. In Ref. , Kniehl and Penin computed the $`\alpha _s^3\mathrm{ln}^2(\alpha _s)`$ terms in the wavefunction at the origin which they refer to as “non-renormalization group logarithms”, since they do not involve factors of the $`\beta `$-function for $`\alpha _s`$. We show that the second term in the series generated by our NLL production current agrees with the result in Ref. . Thus, the solution of the renormalization group equations in the velocity renormalization group method does include these logarithms. Finally, we discuss our result for the NLL $`1/|𝐤|`$ potential and production current.
## II The vNRQCD Lagrangian
The vNRQCD effective Lagrangian has the form
$`=_u+_p+_s.`$ (8)
The ultrasoft Lagrangian $`_u`$ involves the fields $`\psi _𝐩`$ which annihilate a quark, $`\chi _𝐩`$ which annihilate an antiquark, and $`A^\mu `$ which annihilate and create ultrasoft gluons. The potential Lagrangian $`_p`$ contains operators with four or more quark fields including the quark-antiquark potential. Finally, the soft Lagrangian $`_s`$ contains all terms that involve soft particles which have energy and momenta of order $`mv`$. The terms we need in the ultrasoft Lagrangian include
$`_u`$ $`=`$ $`{\displaystyle \frac{1}{4}}F^{\mu \nu }F_{\mu \nu }+{\displaystyle \underset{𝐩}{}}\psi _𝐩^{}\left\{iD^0{\displaystyle \frac{\left(𝐩i𝐃\right)^2}{2m}}+{\displaystyle \frac{𝐩^4}{8m^3}}\right\}\psi _𝐩`$ (10)
$`+{\displaystyle \underset{𝐩}{}}\chi _𝐩^{}\left\{iD^0{\displaystyle \frac{\left(𝐩i𝐃\right)^2}{2m}}+{\displaystyle \frac{𝐩^4}{8m^3}}\right\}\chi _𝐩.`$
The covariant derivative is $`D^\mu =^\mu +ig\mu _U^ϵA^\mu =(D^0,𝐃)`$, so that $`D^0=^0+ig\mu _U^ϵA^0`$, $`𝐃=ig\mu _U^ϵ𝐀`$, and involves only the ultrasoft gluon fields. The ultrasoft scale parameter $`\mu _U=m\nu ^2`$, where $`\nu v`$ is the subtraction velocity. This $`v`$ scaling for $`\mu _U`$ is required for a consistent power counting in $`d`$ dimensions . The covariant derivative on $`\psi _𝐩`$ and $`\chi _𝐩`$ contain the color matrices $`T^A`$ and $`\overline{T}^A`$ for the $`\mathrm{𝟑}`$ and $`\overline{\mathrm{𝟑}}`$ representations, respectively.
The Lagrangian $`_p`$ includes both the traditional quark potential and ultrasoft corrections to this potential which we will denote by $`_{pu}`$:
$`_p={\displaystyle \underset{𝐩,p^{}}{}}V_{\alpha \beta \lambda \tau }(𝐩,𝐩^{})\mu _S^{2ϵ}\psi _{𝐩^{}}^{}{}_{\alpha }{}^{}\psi _{𝐩}^{}{}_{\beta }{}^{}\chi _{𝐩^{}}^{}{}_{\lambda }{}^{}\chi _𝐩{}_{\tau }{}^{}+_{pu}.`$ (11)
The terms we need in $`_{pu}`$ are fixed by reparameterization invariance and will be described below. The on-shell potential $`V(𝐩,𝐩^{})_{\alpha \beta \lambda \tau }`$ has an expansion in $`\alpha _s`$ and $`v`$, and $`\alpha ,\beta ,\lambda ,\tau `$ denote color and spin indices. We will use the color basis in which the potential $`V`$ is written as a linear combination of $`11`$ and $`T\overline{T}`$. The tree level diagrams in Fig. 1 generate terms of $`𝒪(v^{2k}\alpha _s)`$ in the QCD potential. The order $`v^2`$ Coulomb potential is
$`V^{(2)}`$ $`=`$ $`(T^A\overline{T}^A){\displaystyle \frac{𝒱_c^{(T)}}{𝐤^2}}+(11){\displaystyle \frac{𝒱_c^{(1)}}{𝐤^2}},`$ (12)
where the coefficients $`𝒱_c^{(T,1)}`$ have an expansion in $`\alpha _s`$. The order $`v^0`$ potential includes
$`V^{(0)}`$ $`=`$ $`(T^A\overline{T}^A)\left[{\displaystyle \frac{𝒱_2^{(T)}}{m^2}}+{\displaystyle \frac{𝒱_r^{(T)}(𝐩^2+p^2)}{2m^2𝐤^2}}+{\displaystyle \frac{𝒱_s^{(T)}}{m^2}}𝐒^2+{\displaystyle \frac{𝒱_\mathrm{\Lambda }^{(T)}}{m^2}}\mathrm{\Lambda }(𝐩^{},p)+{\displaystyle \frac{𝒱_t^{(T)}}{m^2}}T(𝐤)\right]`$ (14)
$`+(11)\left[{\displaystyle \frac{𝒱_2^{(1)}}{m^2}}+{\displaystyle \frac{𝒱_s^{(1)}}{m^2}}𝐒^2\right],`$
where $`𝐤=𝐩^{}𝐩`$ and
$`𝐒`$ $`=`$ $`{\displaystyle \frac{𝝈_1+𝝈_2}{2}},\mathrm{\Lambda }(𝐩^{},p)=i{\displaystyle \frac{𝐒(p^{}\times p)}{𝐤^2}},T(𝐤)=𝝈_1𝝈_2{\displaystyle \frac{3𝐤𝝈_1𝐤𝝈_2}{𝐤^2}}.`$ (15)
Matching the two diagrams in Fig. 1 to the $`v^2`$ and $`v^0`$ potentials at $`\mu =m`$ gives
$`𝒱_c^{(T)}`$ $`=`$ $`4\pi \alpha _s(m),𝒱_r^{(T)}=4\pi \alpha _s(m),𝒱_s^{(T)}={\displaystyle \frac{4\pi \alpha _s(m)}{3}}+{\displaystyle \frac{1}{N_c}}\pi \alpha _s(m),`$ (16)
$`𝒱_\mathrm{\Lambda }^{(T)}`$ $`=`$ $`6\pi \alpha _s(m),𝒱_t^{(T)}={\displaystyle \frac{\pi \alpha _s(m)}{3}},𝒱_s^{(1)}={\displaystyle \frac{(N_c^21)}{2N_c^2}}\pi \alpha _s(m),`$ (17)
$`𝒱_c^{(1)}`$ $`=`$ $`0,𝒱_2^{(T)}=0,𝒱_2^{(1)}=0.`$ (18)
The LL values for the coefficients of the $`V^{(2)}`$ and $`V^{(0)}`$ potentials will be needed below and are summarized in Appendix A. They are obtained by using the tree-level matching values in Eq. (16) and running using the one loop anomalous dimensions computed in Ref. .
The order $`1/v`$ potential includes
$`V^{(1)}`$ $`=`$ $`{\displaystyle \frac{\pi ^2}{m|𝐤|}}\left[𝒱_k^{(T)}(T^A\overline{T}^A)+𝒱_k^{(1)}(11)\right].`$ (19)
Tree level matching gives $`V^{(1)}=0`$, and $`V^{(1)}`$ has zero one loop anomalous dimension, so $`V^{(1)}=0`$ at LL order as well. At one loop in QCD, order $`1/v`$ potentials of the form in Eq. (19) are generated. The one loop matching of the on-shell potential at $`\mu =m`$ gives
$`𝒱_k^{(T)}`$ $`=`$ $`\alpha _s^2(m)\left({\displaystyle \frac{7C_A}{8}}{\displaystyle \frac{C_d}{8}}\right),𝒱_k^{(1)}=\alpha _s^2(m){\displaystyle \frac{C_1}{2}}.`$ (20)
The color group theory factors are $`C_F=(N_c^21)/(2N_c)`$, $`C_A=N_c`$, $`C_1=(N_c^21)/(4N_c^2)`$, and $`C_d=N_c4/N_c`$. For the color singlet channel Eq. (20) agrees with Refs. . In the language of the threshold expansion the value of the coefficients in Eqs. (16) and (20) are from integrating out off-shell potential gluons at the hard scale $`\mu m`$ where $`\nu =1`$ . In the next section we will compute the two loop anomalous dimension for $`V^{(1)}`$, and determine the NLL value of the coefficients in Eq. (19).
In addition we need ultrasoft corrections to the potential which are contained in $`_{pu}`$. Reparameterization invariance restricts the form of some of these terms by requiring that only the linear combination $`𝐩i𝐃`$ can appear. Here the covariant derivative acts on a quark or antiquark field with label $`𝐩`$. The reparameterization invariant form of the $`T\overline{T}`$ Coulomb potential operator is
$`{\displaystyle \frac{\left[\psi _𝐩^{}^{}T^A\psi _𝐩\right]\left[\chi _𝐩^{}^{}\overline{T}^A\chi _𝐩\right]}{(𝐩^{}𝐩)^2}}`$ $``$ $`{\displaystyle \frac{1}{2}}\left[\psi _𝐩^{}^{}{\displaystyle \frac{T^A}{(𝐩^{}𝐩+𝐢\stackrel{}{𝐃})^2}}\psi _𝐩\right]\left[\chi _𝐩^{}^{}\overline{T}^A\chi _𝐩\right]`$ (22)
$`+{\displaystyle \frac{1}{2}}\left[\psi _𝐩^{}^{}T^A\psi _𝐩\right]\left[\chi _𝐩^{}^{}{\displaystyle \frac{\overline{T}^A}{(𝐩^{}+𝐩+𝐢\stackrel{}{𝐃})^2}}\chi _𝐩\right],`$
where $`\stackrel{}{𝐃}=\stackrel{}{𝐃}+\stackrel{}{𝐃}`$. Terms in the $`v`$ expansion are then generated by expanding Eq. (22) with $`𝐃𝐩^{}𝐩`$. As written the ordering of color generators in Eq. (22) is ambiguous. The correct ordering in the expansion is to write factors of $`\stackrel{}{D}`$ to the right of the $`T^A`$, and factors of $`\stackrel{}{D}`$ to the left of the $`T^A`$. Expanding Eq. (22) and keeping only the terms that we will need gives
$`_{pu}`$ $`=`$ $`{\displaystyle \frac{2i𝒱_c^{(T)}f^{ABC}}{𝐤^4}}\mu _S^{2ϵ}\mu _U^ϵ𝐤(g𝐀^C)\psi _𝐩^{}^{}T^A\psi _𝐩\chi _𝐩^{}^{}\overline{T}^B\chi _𝐩`$ (25)
$`+𝒱_c^{(T)}\mu _S^{2ϵ}\psi _𝐩^{}^{}\left[{\displaystyle \frac{i𝐤\stackrel{}{}}{𝐤^4}}{\displaystyle \frac{\stackrel{}{}^2}{2𝐤^4}}+2{\displaystyle \frac{(𝐤\stackrel{}{})^2}{𝐤^6}}\right]T^A\psi _𝐩\chi _𝐩^{}^{}\overline{T}^A\chi _𝐩`$
$`+𝒱_c^{(T)}\mu _S^{2ϵ}\psi _𝐩^{}^{}T^A\psi _𝐩\chi _𝐩^{}^{}\left[{\displaystyle \frac{i𝐤\stackrel{}{}}{𝐤^4}}{\displaystyle \frac{\stackrel{}{}^2}{2𝐤^4}}+2{\displaystyle \frac{(𝐤\stackrel{}{})^2}{𝐤^6}}\right]\overline{T}^A\chi _𝐩,`$
where $`\stackrel{}{}=\stackrel{}{}+\stackrel{}{}`$. The first term couples an ultrasoft gluon to a four quark operator.
The terms in the second and third lines in Eq. (25) contain order $`1/v`$ and $`v^0`$ terms from the multipole expansion of the Coulomb potential. The first term could also have been determined from the on-shell matching calculation shown in Fig. 2. However, by determining the terms in Eq. (25) using reparameterization invariance rather than matching we know that the coefficients remain equal to $`𝒱_c`$ to all orders in perturbation theory. This saves us from the extra work that would be involved in computing the anomalous dimensions for these terms in $`_{pu}`$.
The terms in the soft Lagrangian include
$`_s`$ $`=`$ $`{\displaystyle \underset{q}{}}\left\{\left|q^\mu A_q^\nu q^\nu A_q^\mu \right|^2+\overline{\phi }_qq/\phi _q+\overline{c}_qq^2c_q\right\}`$ (28)
$`g^2\mu _S^{2ϵ}{\displaystyle \underset{𝐩,p^{},q,q^{}}{}}\{{\displaystyle \frac{1}{2}}\psi _{𝐩^{}}^{}{}_{}{}^{}[A_q^{}^\mu ,A_q^\nu ]U_{\mu \nu }^{(\sigma )}\psi _𝐩+{\displaystyle \frac{1}{2}}\psi _{𝐩^{}}^{}{}_{}{}^{}\{A_q^{}^\mu ,A_q^\nu \}W_{\mu \nu }^{(\sigma )}\psi _𝐩`$
$`+\psi _{𝐩^{}}^{}{}_{}{}^{}[\overline{c}_q^{},c_q]Y^{(\sigma )}\psi _𝐩+(\psi _{𝐩^{}}^{}{}_{}{}^{}T^BZ_\mu ^{(\sigma )}\psi _𝐩)(\overline{\phi }_q^{}\gamma ^\mu T^B\phi _q)\}+(\psi \chi ,T\overline{T}).`$
The fields $`A_q^\mu `$ and $`c_q`$ are the soft gluon and ghost fields, and $`\phi _q`$ is a massless soft quark field with $`n_f`$ flavor components. $`U`$, $`W`$, $`Y`$, and $`Z`$ are functions of $`(𝐩,𝐩^{},q,q^{})`$ and matrices in spin and the index $`\sigma `$ denotes the relative order in the $`v`$ expansion. For Feynman gauge and the case $`𝐩^2=𝐩_{}^{}{}_{}{}^{2}`$ these functions were derived in Ref. . Beyond one loop terms proportional to $`(𝐩^2𝐩_{}^{}{}_{}{}^{2})`$ will be needed in $`_s`$, since besides its soft energy $`q^0mv`$ the $`A_q`$ gluons can carry away a residual energy of order $`mv^2`$. The LL values for the functions $`U`$, $`W`$, $`Y`$, and $`Z`$ can be found in Appendix A of Ref. . In addition, some NLL contributions to $`_s`$ will be needed and will be discussed in section III A. These additional contributions are obtained from one loop matching with two-loop renormalization group improvement.
As an aside, note that it is not necessary to consider the ultrasoft renormalization of the vertices given in the soft Lagrangian in Eq. (28). One might think that diagrams such as
(29)
would effect the running of the coefficients in the soft Lagrangian. Diagrams analogous to the one in Eq. (29), but with only soft gluons and quarks generate the running of the coefficient functions in the soft LagrangianIn Ref. the running of $`_s`$ was calculated by examining loops with soft gluons that contribute to the Compton scattering process prior to integrating out the soft quarks. There it was noted that all ultraviolet divergent diagrams are in one-to-one correspondence with graphs in HQET, so that the running of this Lagrangian could be obtained from the known running in HQET . in Ref. . However, the graph in Eq. (29) has only one heavy quark, so a distinction between soft and ultrasoft gluons is unnecessary at this point. Noting that the soft vertices will always occur in pairs, it is in fact consistent to only dress pairs of the soft vertices by ultrasoft gluons:
(30)
Since these diagrams involve two heavy quarks, both types of gluons can occur. Since in the end it is only graphs such as Eq. (30) with two soft vertices that are relevant for constructing the theory, it is consistent to include the ultrasoft renormalization of soft vertices as a contribution to the four-quark operator in Eq. (30), rather than treating the subgraph as a contribution to the soft vertex, as in Eq. (29). Along with the diagrams in Eq. (30) there are graphs in which the ultrasoft gluon is exchanged between the two heavy quarks.
## III Two loop anomalous dimension for $`V^{(1)}`$
To calculate the NLL anomalous dimension for the $`1/|𝐤|`$ potentials we need to consider graphs in the effective theory of orderThis is the size (in $`v`$) of the amputated diagrams, so in contrast to the general power counting formula in Ref. we are not including the powers of $`v`$ generated by external lines. $`\alpha _s^3/v`$. These diagrams come in two classes, those with soft gluons, and those with a single ultrasoft gluon. The total anomalous dimension for the $`1/|𝐤|`$ potential is $`\gamma ^{(1,T)}=\nu d/d\nu 𝒱_k^{(1,T)}`$. Since $`\mu _S=m\nu `$ and $`\mu _U=m\nu ^2`$, $`\gamma ^{(1,T)}`$ can be written as the sum of a soft and ultrasoft anomalous dimension
$`\gamma ^{(1,T)}=\gamma _S^{(1,T)}+\mathrm{\hspace{0.17em}2}\gamma _U^{(1,T)},`$ (32)
where $`\gamma _S^{(1,T)}=\mu _S/\mu _S𝒱_k^{(1,T)}`$ and $`\gamma _U^{(1,T)}=\mu _U/\mu _U𝒱_k^{(1,T)}`$. In the remainder of this section we discuss the computation of these two loop anomalous dimensions in detail. The calculation is split into two parts, graphs with soft vertices and graphs with ultrasoft vertices. Graphs with soft gluons only contribute to $`\gamma _S`$, while those with an ultrasoft gluon contribute to both $`\gamma _S`$ and $`\gamma _U`$.
### A Soft contributions
The order $`\alpha _s^3/v`$ diagrams containing soft gluons that contribute to the anomalous dimension are shown in Fig. 3. The sum of diagrams forms a gauge invariant set.
The two loop graphs in Figs. 3a and 3b involve an iteration of a potential and a soft loop. The vertices in these graphs are of LL order (tree level matching with one loop running) and are given in Ref. . In Figs. 3a and 3b, counting powers of $`v`$ from the propagators and from the loop measures gives a $`v^1`$, so the sum of powers of $`v`$ for the three (amputated) vertices must give an overall $`1/v^2`$. The $`v`$ scaling for the two soft vertices is $`\sigma +\sigma ^{}2`$. In Fig. 3 we show two possibilities: a) has one $`V^{(2)}`$ insertion and two soft vertices from Eq. (28) such that $`\sigma ^{}+\sigma =2`$, and b) has one insertion of a $`V^{(0)}`$ potential from Eq. (14) and two soft vertices such that $`\sigma ^{}=\sigma =0`$. We could also have a $`V^{(1)}`$ potential plus two soft vertices where $`\sigma +\sigma ^{}=1`$; however this diagram is identically zero. The graphs where the potential and soft loop are exchanged simply give a factor of 2.
The one loop graphs in Fig. 3c, 3d, and 3e involve additional vertices in $`_s`$, denoted $``$, which are of NLL order (from one loop matching with two loop running). The one loop matching calculation for these vertices is sketched in Fig. 4.
There are a large number of diagrams in the full theory (graphs on the left hand side), so only a few representative examples have been shown. To obtain the values for the operators on the right hand side we subtract purely soft effective theory diagrams from those in the full theory. To see how these operators arise, it is useful to recall that in the threshold expansion soft heavy quarks have a propagator
$`{\displaystyle \frac{1}{q_0+iϵ}}=\mathrm{P}{\displaystyle \frac{1}{q_0}}i\pi \delta (q_0),`$ (33)
where $`\mathrm{P}`$ stands for the principal value. In our approach, off-shell potential gluons and soft quarks are integrated out at the scale $`m`$ when constructing the effective theory. When integrating out the soft heavy quarks the principal value term in Eq. (33) goes directly into a coefficient in the soft Lagrangian since this term is consistent with the scaling in the soft regime, $`q_0mv`$. For instance, in Eq. (28), $`U_{00}^{(0)}=1/q_0`$. The delta function contribution in Eq. (33) is associated with the potential regime since $`q_00`$. When the delta function appears in a loop in the full theory (or threshold expansion) it forces gluons in the loop to have zero energy, or in other words to become potential gluons. It is these contributions which do not appear in the soft effective theory diagrams and must be made up by the operators shown on the right hand side of Fig. 4.
The total contribution to the anomalous dimensions from the soft diagrams in Fig. 3 is
$`\gamma _S^{(T)}`$ $`=`$ $`{\displaystyle \frac{\beta _0(7C_AC_d)}{8\pi }}\alpha _s^3(m\nu ){\displaystyle \frac{8C_A(C_A+C_d)}{3\pi }}\alpha _s^3(m\nu ),`$ (34)
$`\gamma _S^{(1)}`$ $`=`$ $`{\displaystyle \frac{\beta _0C_1}{2\pi }}\alpha _s^3(m\nu )+{\displaystyle \frac{16C_AC_1}{\pi }}\alpha _s^3(m\nu ),`$ (35)
where $`\beta _0=11C_A/34T_Fn_f/3`$, $`n_f`$ is the number of massless soft quarks, and $`𝒱_c(\nu )=4\pi \alpha _s(m\nu )`$ was used. In Eq. (34) the terms proportional to $`\beta _0`$ can be inferred from Eq. (20). They simply turn the $`\alpha _s(m)`$’s in the matching result into running $`\alpha _s`$’s. At one and two loops terms proportional to the $`\beta `$-function for $`\alpha _s`$ completely determine the soft anomalous dimension for the Coulomb potential. However, the $`1/|𝐤|`$ potentials have additional contributions because the soft diagrams in Fig. 3c, 3d, and 3e have infrared divergences. The IR divergences from purely soft diagrams are not true IR divergences in the effective theory. For instance, in general they do not match up with IR divergences in QCD. The true IR divergences are from momenta $`<mv^2`$, and the desired ultraviolet divergences are from momenta $`m`$. Instead, the soft IR divergences are from momenta $`<mv`$, and match up with ultrasoft UV divergences that are from momenta $`mv`$ to carry these UV divergences up to the hard scale. Writing the soft IR divergence
$`{\displaystyle \frac{1}{ϵ_{IR}}}={\displaystyle \frac{1}{ϵ_{UV}}}\left({\displaystyle \frac{1}{ϵ_{UV}}}{\displaystyle \frac{1}{ϵ_{IR}}}\right),`$ (36)
the first term contributes to the anomalous dimension in Eq. (34). The $`1/ϵ_{\mathrm{UV}}1/ϵ_{\mathrm{IR}}`$ term can be ignored since it simply takes the corresponding ultrasoft UV divergence up to the hard scale $`m`$.<sup>\**</sup><sup>\**</sup>\**This result, that there are no true IR divergences in the soft regime has not been proven to all orders in perturbation theory. However, its is likely that all IR divergences can be attributed to ultrasoft and collinear gluons in the spirit of the Coleman-Norton theorem , plus IR divergences associated with the Coulomb regime that are reproduced by iterations of the potential. Thus, to compute $`\gamma _S`$, all divergences from soft loops should be treated as UV divergences. The terms not proportional to $`\beta _0`$ in Eq. (34) can be inferred from the result of the ultrasoft calculation in Eq. (49) of the next section.
Despite the fact that Eq. (34) can be inferred without a direct calculation it is worthwhile to examine the diagrams in Fig. 3. Consider the graphs in Figs. 3a and 3b in Feynman gauge. The sub-loop with soft gluons is divergent, while the remaining potential loop is convergent. There is also a set of one loop diagrams (not shown) where the soft sub-loop is replaced by the one loop counterterms for $`V`$ derived in Ref. . We find that these counterterm graphs exactly cancel against a set of divergences in Fig. 3a and 3b. For Fig. 3b there is an exact cancellation, and so there is no operator mixing between $`V^{(0)}`$ and $`V^{(1)}`$. However, there are divergences that appear in Fig. 3a that have no corresponding counterterm graphs. Consider the soft gluon case (the ghost and soft quark cases are similar). After performing the $`k^0`$ integration the loop integral for Fig. 3a is
$`{\displaystyle d^{d1}𝐪\frac{1}{(E𝐪^2/m)(𝐩q)^2}\left[d^dt\frac{U_{\mu \nu }^{(\sigma )}U^{\mu \nu (\sigma ^{})}}{t^2[(t^0)^2(𝐭+qp^{})^2]}\right]},`$ (37)
where $`𝐩`$ and $`𝐩^{}`$ are the momenta of the incoming and outgoing quarks, $`E=𝐩^\mathrm{𝟐}/m`$ and in Eq. (37) the $`U`$’s depend on $`𝐪`$, $`𝐭`$, $`𝐩^{}`$ and $`t^0`$. For the divergences of interest performing the $`t`$ integral (the soft loop) in $`d=42ϵ`$ dimensions gives a factor of
$`{\displaystyle \frac{(𝐩^{}{}_{}{}^{\mathrm{𝟐}}𝐪^\mathrm{𝟐})^2}{m^2(𝐩^{}𝐪)^{4+2ϵ}ϵ}}+\mathrm{},`$ (38)
where the ellipsis denote order $`ϵ^0`$ terms. The remaining loop integration is finite. For the one loop graph that corresponds to the soft sub-loop, the loop momenta $`𝐪`$ in Eq. (38) is replaced by $`𝐩^{}`$, and this diagram vanishes on-shell by energy conservation so no counterterm is generated. Performing the final integration and including the factor of $`2`$ from the left-right symmetric graph gives
$`\text{Fig.}\text{3}a=i𝒱_c^{(T)}(\mu _S)\alpha _s^2(\mu _S){\displaystyle \frac{\beta _0\mu _S^{2ϵ}}{32m|𝐤|}}(T^AT^B\overline{T}^A\overline{T}^B)\left[{\displaystyle \frac{1}{ϵ}}+2\mathrm{ln}\left({\displaystyle \frac{\mu _S^2}{|𝐤|^2}}\right)+\mathrm{}\right].`$ (39)
The divergence in Eq. (39) contributes to the anomalous dimension for $`𝒱_k`$. This seems slightly unusual because it was generated by a sub-divergence, whereas usually only the overall divergence in a diagram is relevant. In this diagram the loop integral with soft gluons generates a $`1/ϵ`$ pole, and the remaining potential loop integral generates the $`1/|𝐤|`$ factor. The corresponding diagrams in QCD include graphs such as the vacuum polarization of one of the gluons in the box diagram. In this full theory graph the subdivergence due to the vacuum polarization insertion would be canceled by a counterterm. However, in the effective theory this divergence is instead absorbed into the coefficients of terms in the quark potential because the potential gluon components have been integrated out. Since matching the full theory box diagram gives a contribution to the $`1/|𝐤|`$ potential, it is not surprising that gluon vacuum polarization in the box graph contributes to the anomalous dimension of this potential. An alternative to the approach used to derive Eq. (39) is to use off-shell matching and running. In this case the soft anomalous dimension analysis will be different (since, for instance, the off-shell potential is gauge dependent). This approach is discussed in Appendix B, where it is shown that the final answer for physical observables is unchanged.
The diagrams in Figs. 3c, 3d, and 3e are tedious to evaluate due to the large number of diagrams necessary for the matching calculation in Fig. 4.<sup>††</sup><sup>††</sup>††In fact once a contribution to $``$ has been identified it is simpler to directly evaluate the two-loop diagram obtained by combining the steps in Fig. 4 and Fig. 3. Since the graphs in Figs. 3c, 3d, and 3e are one-loop diagrams there is no cancellation from counterterms. In Feynman gauge
$`\text{Fig.}\text{3}c,d,e`$ $`=`$ $`i{\displaystyle \frac{\pi ^2\mu _S^{2ϵ}}{m|𝐤|}}{\displaystyle \frac{\alpha _s^3(\mu _S)}{4\pi ϵ}}\{{\displaystyle \frac{4}{3}}C_1T_Fn_f(11)+(C_AT_Fn_f{\displaystyle \frac{1}{3}}C_dT_Fn_f)(T^A\overline{T}^A)`$ (41)
$`+{\displaystyle \frac{37}{3}}C_AC_1(11)({\displaystyle \frac{7}{4}}C_AC_d+{\displaystyle \frac{65}{12}}C_A^2)(T^A\overline{T}^A)\}.`$
Only the terms proportional to $`n_f`$ in Eq. (41) have been checked by direct calculation. Also, in Eq. (41) we have assumed that the $``$ operators simply run with the strong coupling constant, $`\alpha _s(m\nu )`$. To motivate this, recall that a theory of propagating soft quarks and gluons has the same singularity structure as Heavy Quark Effective Theory (HQET). The operators that are generated in Fig. 4 correspond to diagrams in HQET with insertions of $`^2/(m)`$ or a $`𝐩𝐀/m`$ vertex. Since these operators do not run in HQET the operators that are constructed in Fig. 4 should also simply run with the strong coupling constant.
### B Ultrasoft contributions
Next consider the order $`\alpha _s^3/v`$ diagrams with ultrasoft gluons. At one loop we have a single insertion of the $`V^{(1)}`$ potential dressed by an ultrasoft gluon with $`gA^0`$ couplings. These diagrams are obviously zero in Coulomb gauge and in Feynman gauge it was shown in Ref. that, with any potential, the sum of this set of one loop diagrams is also identically zero. At two loops in an arbitrary gauge there are many possible contributions. There are diagrams with two $`V^{(2)}`$ vertices, and an ultrasoft tadpole generated by the seagull $`𝐀^\mathrm{𝟐}`$ operator attached to one of the quark lines. These graphs do not have logarithmic divergences and therefore do not contribute to the anomalous dimension. There are also diagrams that are zero because they are odd in an ultrasoft momenta which we can omit. Next consider the topologies shown in Fig. 5. Several classes of diagrams are generated depending on the vertices used:
1. vertex $`1`$ and $`2`$ are $`V^{(2)}`$, and the ultrasoft gluon couples with $`𝐩𝐀/m`$,
2. vertex $`1`$ and $`2`$ are $`V^{(2)}`$, the ultrasoft gluon couples with $`gA^0`$, and there are two insertions of $`𝐩/m`$ on quark lines,
3. vertex $`1`$ and $`2`$ are $`V^{(2)}`$, the ultrasoft gluon couples with $`gA^0`$, and there is one insertion of $`^\mathrm{𝟐}/m`$ on a quark line,
4. vertex $`1`$ and $`2`$ are $`V^{(2)}`$, the ultrasoft gluon couples with $`gA^0`$, and there is one insertion of $`𝐩^\mathrm{𝟒}/m^3`$ on a quark line (together with the expansion of factors of the energy that appear in lower order diagrams which should be viewed as corrections to the effective theory states, see Refs. ),
5. vertex $`1`$ is $`V^{(2)}`$, vertex $`2`$ is $`V^{(0)}`$ and the ultrasoft gluon couples with $`gA^0`$,
6. vertex $`1`$ is the order $`1/v`$ potential from the multipole expansion of the Coulomb potential given in Eq. (25), vertex $`2`$ is $`V^{(2)}`$, and the ultrasoft gluon couples with $`gA^0`$ vertices.
7. vertex $`1`$ is the order $`v^0`$ potential in Eq. (25), vertex $`2`$ is $`V^{(2)}`$, there is one insertion of $`𝐩/m`$ on a quark line, and the ultrasoft gluon couples with $`gA^0`$ vertices.
8. vertex $`1`$ and $`2`$ are $`v^0`$ potentials from Eq. (25), and the ultrasoft gluon couples with $`gA^0`$ vertices.
Insertions of operators with $``$’s only need to be considered on quark propagators where the multipole expansion was used; together with the graphs in cases 6, 7, and 8 they build up the sub-leading terms in the multipole expansion.
The graphs in cases 1, 2, and 5 do not give any contribution to the two loop anomalous dimension. The reason is that all ultraviolet divergences are exactly canceled by one loop counterterm graphs (these counterterms are generated by one loop graphs having one potential insertion dressed by an ultrasoft gluon, and in Feynman gauge were calculated in Refs. ). The graphs with the topology in Fig. 5c,d are ultraviolet finite. For the remaining diagrams we can identify the corresponding counterterm graph by simply shrinking the smallest loop containing the ultrasoft gluon to a point. The only subtlety occurs in Fig. 5h which has overlapping potential and ultrasoft loop integrals. For heavy scalars this graph was analyzed in Ref. with the threshold expansion, while in an effective theory for heavy fermions this diagram was analyzed in Ref. for the case where the ultrasoft gluon is massive and couples with derivatives at the vertices. Since the nature of the overlapping integrals does not depend on the structure of the numerator or on having a massive gluon, this analysis will not be repeated here. In the effective theory this diagram comes with the appropriate factor of $`2`$, so that it is canceled by the two one loop counterterm graphs that correspond to either making the loop integral on the right large or making the loop integral on the left large. Note that for a $`\varphi ^3`$ relativistic theory in $`d=6`$ , the sub-divergences for this diagram are also canceled in this way, but leave an overall divergence. For the non-relativistic effective theory diagram this overall divergence is not present.
For the graphs in case 4, the sum of insertions on quark lines inside the loop shared by the ultrasoft gluon are finite. The graphs with insertions on quark lines outside this loop are ultraviolet divergent, but are exactly canceled by the counterterm diagrams that correspond to shrinking the ultrasoft loop to a point. Therefore, the graphs in case 4 also do not contribute to the anomalous dimension.
In a general gauge the graphs in cases 3, 6, 7, and 8 will contribute to the anomalous dimension. There are also additional graphs with the vertex in Eq. (22) which involves the coupling of an ultrasoft $`𝐀^i`$ gluon to a potential vertex. Together this set of graphs gives a gauge invariant contribution to the anomalous dimension, so we can simplify their computation by choosing the most convenient gauge. We will choose Coulomb gauge since the graphs in cases 3, 6, 7, and 8 involve ultrasoft $`A^0`$ gluons and vanish in this gauge. The result for the infinite part of the the remaining diagrams in Coulomb gauge is:
$`=`$ $`0,`$ (43)
$`=`$ $`{\displaystyle \frac{i\alpha _s(\mu _U)\mu _S^{2ϵ}[𝒱_c^{(T)}]^2}{8\pi m|𝐤|ϵ}}\left[C_AC_1\mathrm{\hspace{0.25em}1}1{\displaystyle \frac{1}{4}}C_A(C_A+C_d)T^A\overline{T}^A\right],`$ (44)
$`=`$ $`{\displaystyle \frac{i\alpha _s(\mu _U)\mu _S^{2ϵ}[𝒱_c^{(T)}]^2}{8\pi m|𝐤|ϵ}}{\displaystyle \frac{C_AC_1}{3}}11,`$ (45)
$`=`$ $`{\displaystyle \frac{i\alpha _s(\mu _U)\mu _S^{2ϵ}[𝒱_c^{(T)}]^2}{8\pi m|𝐤|ϵ}}\left[{\displaystyle \frac{2}{3}}C_AC_1\mathrm{\hspace{0.25em}1}1{\displaystyle \frac{1}{12}}C_A(C_A+C_d)T^A\overline{T}^A\right].`$ (46)
Equations (5b,c) include the factor of four from their left-right and up-down mirror graphs. In the graphs in Eq. (5) the factor
$$\frac{\mu _U^{2ϵ}\mu _S^{4ϵ}}{ϵ\left|𝐤\right|^{1+2ϵ}E^{2ϵ}}=\frac{\mu _S^{2ϵ}}{|𝐤|}\left[\frac{1}{ϵ}+\mathrm{ln}\left(\frac{\mu _S^2}{\left|𝐤\right|^2}\right)+\mathrm{ln}\left(\frac{\mu _U^2}{E^2}\right)\right]$$
(48)
occurs, where $`E=𝐩^\mathrm{𝟐}/m`$. The first term is included in the result displayed in Eq. (5), and the remaining terms are not shown. For $`\nu v`$, we have $`\mu _S\left|k\right|mv`$, $`\mu _UEmv^2`$, so the soft and ultrasoft scale factors correctly minimize possible large logarithms in the effective theory.
From the sum of diagrams in Eq. (5) we have the following result for the soft and ultrasoft anomalous dimensions of the order $`1/|𝐤|`$ potentials:
$`\gamma _S^{(T)}=\gamma _U^{(T)}`$ $`=`$ $`C_A(C_A+C_d){\displaystyle \frac{\alpha _s(m\nu ^2)[𝒱_c^{(T)}(\nu )]^2}{12\pi ^3}},`$ (49)
$`\gamma _S^{(1)}=\gamma _U^{(1)}`$ $`=`$ $`C_AC_1{\displaystyle \frac{\alpha _s(m\nu ^2)[𝒱_c^{(T)}(\nu )]^2}{2\pi ^3}}.`$ (50)
Recall that the soft and ultrasoft anomalous dimensions are given by differentiating with respect to $`\mathrm{ln}\mu _S`$ and $`\mathrm{ln}\mu _U`$ respectively.
It is interesting to ask how the result in Eq. (49) would be reproduced if terms that vanish by the equations of motion were included in our effective Lagrangian. It is possible to include a term in our effective Lagrangian which makes the graph in Eq. (5b) give no contribution to the anomalous dimension, but the final result for observables remains invariant. This example is discussed in Appendix B.
## IV Results
In this section the NLL heavy quark $`1/|𝐤|`$ potential and production current are discussed. After presenting the renormalization group improved results, we re-expand to compare to the finite order results in Refs. and . We also discuss the behavior of the NLL $`1/|𝐤|`$ potential and the NLL production current as we run down from $`\nu =1`$ to $`\nu =v`$, where $`v`$ is the Coulombic velocity.
### A The NLL $`1/|𝐤|`$ potentials
The total soft and ultrasoft anomalous dimensions for the on-shell $`1/|𝐤|`$ potentials are obtained by adding Eqs. (34) and (49). Using the LL coefficient for the Coulomb potential, $`𝒱_c(\nu )=4\pi \alpha _s(m\nu )`$, we find
$`\gamma _S^{(T)}`$ $`=`$ $`{\displaystyle \frac{1}{8\pi }}\beta _0(7C_AC_d)[\alpha _s(m\nu )]^3{\displaystyle \frac{8}{3\pi }}C_A(C_A+C_d)[\alpha _s(m\nu )]^3`$ (52)
$`+{\displaystyle \frac{4}{3\pi }}C_A(C_A+C_d)\alpha _s(m\nu ^2)[\alpha (m\nu )]^2,`$
$`\gamma _S^{(1)}`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}\beta _0C_1[\alpha _s(m\nu )]^3+{\displaystyle \frac{16}{\pi }}C_AC_1[\alpha _s(m\nu )]^3{\displaystyle \frac{8}{\pi }}C_AC_1\alpha _s(m\nu ^2)[\alpha _s(m\nu )]^2,`$ (53)
$`\gamma _U^{(T)}`$ $`=`$ $`{\displaystyle \frac{4}{3\pi }}C_A(C_A+C_d)\alpha _s(m\nu ^2)[\alpha _s(m\nu )]^2,`$ (54)
$`\gamma _U^{(1)}`$ $`=`$ $`{\displaystyle \frac{8}{\pi }}C_AC_1\alpha _s(m\nu ^2)[\alpha _s(m\nu )]^2.`$ (55)
In full QCD (taking the QCD scale parameter $`\mu =m`$), the first logarithm generated by the ultrasoft anomalous dimension corresponds to a $`\mathrm{ln}(E/m)`$, while the first logarithm generated by the soft anomalous dimension corresponds to a $`\mathrm{ln}(|𝐤|/m)`$. For the velocity renormalization group, the total anomalous dimension is simply $`\gamma _S+2\gamma _U`$:
$`\nu {\displaystyle \frac{}{\nu }}𝒱_k^{(T)}(\nu )`$ $`=`$ $`{\displaystyle \frac{1}{8\pi }}\beta _0(7C_AC_d)[\alpha _s(m\nu )]^3{\displaystyle \frac{8}{3\pi }}C_A(C_A+C_d)[\alpha _s(m\nu )]^3`$ (57)
$`+{\displaystyle \frac{4}{\pi }}C_A(C_A+C_d)\alpha _s(m\nu ^2)[\alpha (m\nu )]^2,`$
$`\nu {\displaystyle \frac{}{\nu }}𝒱_k^{(1)}(\nu )`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}\beta _0C_1[\alpha _s(m\nu )]^3+{\displaystyle \frac{16}{\pi }}C_AC_1[\alpha _s(m\nu )]^3{\displaystyle \frac{24}{\pi }}C_AC_1\alpha _s(m\nu ^2)[\alpha _s(m\nu )]^2.`$ (58)
Note that no other $`1/v`$ potentials are generated for $`\nu <1`$ by operator mixing. Integrating Eq. (57) using the one loop $`\beta `$-function for $`\alpha _s`$ and the one loop boundary condition in Eq. (20) gives
$`𝒱_k^{(T)}(\nu )`$ $`=`$ $`{\displaystyle \frac{(7C_AC_d)}{8}}\alpha _s^2(m)+\left[{\displaystyle \frac{(7C_AC_d)}{8}}+{\displaystyle \frac{8C_A(C_A+C_d)}{3\beta _0}}\right](z^21)\alpha _s^2(m)`$ (60)
$`+{\displaystyle \frac{8C_A(C_A+C_d)}{\beta _0}}\left[z12\mathrm{ln}(w)\right]\alpha _s^2(m),`$
$`𝒱_k^{(1)}(\nu )`$ $`=`$ $`{\displaystyle \frac{C_1}{2}}\alpha _s^2(m)+\left[{\displaystyle \frac{C_1}{2}}{\displaystyle \frac{16C_AC_1}{\beta _0}}\right](z^21)\alpha _s^2(m){\displaystyle \frac{48C_AC_1}{\beta _0}}\alpha _s^2(m)\left[z12\mathrm{ln}(w)\right],`$ (61)
where
$`z={\displaystyle \frac{\alpha _s(m\nu )}{\alpha _s(m)}},w={\displaystyle \frac{\alpha _s(m\nu ^2)}{\alpha _s(m\nu )}}={\displaystyle \frac{1}{2z}}.`$ (62)
Projecting onto the color singlet channel, $`𝒱^{(s)}=𝒱^{(1)}C_F𝒱^{(T)}`$, and setting $`C_d=8C_F3C_A`$ and $`C_1=C_AC_F/2C_F^2`$ gives
$`𝒱_k^{(s)}(\nu )`$ $`=`$ $`\left({\displaystyle \frac{C_F^2}{2}}C_FC_A\right)\alpha _s^2(m)+\left[{\displaystyle \frac{C_F^2}{2}}C_FC_A{\displaystyle \frac{8C_AC_F(C_A+2C_F)}{3\beta _0}}\right](z^21)\alpha _s^2(m)`$ (64)
$`{\displaystyle \frac{8C_AC_F(C_A+2C_F)}{\beta _0}}\left[z12\mathrm{ln}(w)\right]\alpha _s^2(m).`$
The projection onto the color octet channel is $`𝒱_k^{(o)}=𝒱_k^{(1)}+(C_A/2C_F)𝒱_k^{(T)}`$.
Logarithmic corrections to the color singlet $`V^{(1)}`$ potential were also considered by Brambilla et al. using pNRQCD, but were not resummed. Brambilla et al. have
$`𝒱_k^{(s)}(\nu )`$ $`=`$ $`C_FC_A\alpha _s^2(r){\displaystyle \frac{4C_AC_F(C_A+2C_F)}{3\pi }}\alpha _s(\mu )\alpha _s(r)^2\mathrm{ln}(\mu r).`$ (65)
To compare this expression to ours: $`\alpha (r)\alpha (m\nu )`$, $`r1/(m\nu )`$ and $`\mu =\mu _U=m\nu ^2`$ since $`r`$ corresponds to the soft scale and $`\mu `$ corresponds to the ultrasoft scale. Expanding Eq. (65) in $`\alpha _s(m)`$ gives
$`𝒱_k^{(s)}(\nu )`$ $`=`$ $`C_FC_A\alpha _s^2(m)+{\displaystyle \frac{\beta _0C_FC_A}{\pi }}\alpha _s^3(m)\mathrm{ln}[1/(mr)]`$ (67)
$`{\displaystyle \frac{4C_AC_F(C_A+2C_F)}{3\pi }}\alpha _s^3(m)\mathrm{ln}(\mu r)+\mathrm{}.`$
To compare this to our result we expand the resummed logarithms in Eq. (64):
$`𝒱_k^{(s)}(\nu )`$ $`=`$ $`\left({\displaystyle \frac{C_F^2}{2}}C_FC_A\right)\alpha _s^2(m)+{\displaystyle \frac{\beta _0C_FC_A}{\pi }}\alpha _s^3(m)\mathrm{ln}(\nu ){\displaystyle \frac{\beta _0C_F^2}{2\pi }}\alpha _s^3(m)\mathrm{ln}(\nu )`$ (69)
$`{\displaystyle \frac{4C_AC_F(C_A+2C_F)}{3\pi }}\alpha _s(m)^3\mathrm{ln}(\nu )+\mathrm{}.`$
In Eq. (69), the first and second $`\mathrm{ln}(\nu )`$ terms are entirely from the soft anomalous dimension in Eq. (52), while the third $`\mathrm{ln}(\nu )`$ term is from a combination of the ultrasoft and soft anomalous dimensions. The first $`\mathrm{ln}(\nu )`$ term in Eq. (69) agrees with the $`\mathrm{ln}[1/(mr)]=\mathrm{ln}(\nu )`$ term in Eq. (67), and the third $`\mathrm{ln}(\nu )`$ term in Eq. (69) agrees with the $`\mathrm{ln}(\mu r)=\mathrm{ln}(\nu )`$ term in Eq. (67).
The second term in Eq. (69) does not appear in Brambilla et al.’s expression in Eq. (67). This is because it depends on whether an on-shell or off-shell potential is used for the matching and running<sup>‡‡</sup><sup>‡‡</sup>‡‡The issue of on-shell versus off-shell potential was discussed in Ref. , and in the context of the leading log summation is discussed further in Appendix B.. We have used an on-shell potential, while Ref. uses off-shell Coulomb gauge which includes a $`𝒱_{\mathrm{\Delta }2}^{(T)}=4\pi \alpha _s(r)`$ potential as defined in Eq. (B1) of Appendix B. The second logarithm in Eq. (69) should appear from the coefficient of this potential. Transforming Brambilla et al.’s $`𝒱_{\mathrm{\Delta }2}^{(T)}`$ potential to a $`𝒱_k^{(s)}`$ potential using Eq. (B7) gives
$`𝒱_k^{(s)}(\nu )`$ $`=`$ $`{\displaystyle \frac{C_F^2}{2}}\alpha _s(r)^2`$ (70)
$`=`$ $`{\displaystyle \frac{C_F^2}{2}}\alpha _s^2(m){\displaystyle \frac{\beta _0C_F^2}{2\pi }}\alpha _s^3(m)\mathrm{ln}[1/(mr)]+\mathrm{}.`$ (71)
The term with a $`\mathrm{ln}[1/(mr)]=\mathrm{ln}(\nu )`$ agrees with the second $`\mathrm{ln}(\nu )`$ term in Eq. (69). The sum of terms without logarithms in Eqs. (67) and (70) also agrees with Eq. (69). Note that the next term in the series in Eq. (69) does not give agreement with Eqs. (67) and (70) since Brambilla et al. did not attempt to systematically sum all the logarithms.
To see the effect of the running on the value of the $`V^{(1)}`$ potential, consider the case of top quark production near threshold. Using $`\alpha _s(m_t)=0.108`$ and Eq. (64), the one loop matching value is:
$`𝒱_k^{(s)}(1)=0.0362.`$ (72)
For a Coulombic system we determine the velocity $`v`$ by solving $`\alpha _s(mv)=v`$. Using $`m_t=175\mathrm{GeV}`$ and the one loop running of $`\alpha _s(\mu )`$ with $`n_f=5`$ gives $`v=0.145`$. The NLL color singlet and octet coefficients for the $`1/|𝐤|`$ potential are plotted in Fig. 6. At $`\nu =v`$ the running coupling is
$`𝒱_k^{(s)}(v)=0.0027,`$ (73)
which is a substantial change from Eq. (72). In comparison, the terms shown in Eq. (69) where the combination $`\alpha _s\mathrm{ln}(v)`$ is treated perturbatively give $`𝒱_k^{(s)}(v)=0.0313`$. Thus, the summation of logarithms for the $`1/|𝐤|`$ potential is quite important; it decreases the coefficient by an order of magnitude. Using the unexpanded results in Eqs. (65) and (70) gives $`𝒱_k^{(s)}(v)=0.0201`$, so our resummed result is also quite different from the fixed order result in Ref. .
### B The NLL $`t\overline{t}`$ production current
In terms of the running color singlet potentials the anomalous dimension for the production current in Eq. (6) is:
$`\gamma _{c_1}(\nu )=\nu {\displaystyle \frac{}{\nu }}\mathrm{ln}[c_1(\nu )]`$ $`=`$ $`{\displaystyle \frac{𝒱_c^{(s)}(\nu )}{16\pi ^2}}\left({\displaystyle \frac{𝒱_c^{(s)}(\nu )}{4}}+𝒱_2^{(s)}(\nu )+𝒱_r^{(s)}(\nu )+𝐒^2𝒱_s^{(s)}(\nu )\right)+{\displaystyle \frac{𝒱_k^{(s)}(\nu )}{2}}.`$ (74)
For the vector production current in Eq. (6) only spin-1 states are produced, and we can set $`𝐒^2=S(S+1)=2`$. However, our analysis also applies to the scalar production current
$`c(\nu ){\displaystyle \underset{𝐩}{}}\psi _{𝐩}^{}{}_{}{}^{}\chi _𝐩^{},`$ (76)
which, for example, contributes to the process $`\gamma \gamma t\overline{t}`$ . Therefore, we will keep the factors of $`𝐒^\mathrm{𝟐}`$ explicit in our results. The boundary condition is given by the matching condition at $`\nu =1`$ ($`\mu =m`$) which is known to two loops . For the NLL approximation only the one loop matching condition should be used. For the vector current
$`c_1(1)=1{\displaystyle \frac{2C_F\alpha _s(m)}{\pi }}.`$ (77)
Integrating Eq. (74) with the running potentials in Appendix A, Eq. (A1) and in Eq. (60) we find
$`\mathrm{ln}\left[{\displaystyle \frac{c_1(\nu )}{c_1(1)}}\right]`$ $`=`$ $`a_1\pi \alpha _s(m)\left({\displaystyle \frac{1}{z}}1\right)+a_2\pi \alpha _s(m)(1z)+a_3\pi \alpha _s(m)\mathrm{ln}(z)`$ (81)
$`+a_4\pi \alpha _s(m)\left[1z^{113C_A/(6\beta _0)}\right]+a_5\pi \alpha _s(m)\left[1z^{12C_A/\beta _0}\right]`$
$`+a_6\pi \alpha _s(m)\left[{\displaystyle \frac{\pi ^2}{12}}{\displaystyle \frac{1}{2}}\mathrm{ln}^2(2)\mathrm{ln}(w)\mathrm{ln}\left({\displaystyle \frac{2w}{2w1}}\right)\mathrm{Li}_2\left({\displaystyle \frac{1}{2w}}\right)\right]`$
$`+a_7\pi \alpha _s(m)\left[{\displaystyle \frac{w}{2w1}}\mathrm{ln}(w){\displaystyle \frac{1}{2}}\mathrm{ln}(2w1)\right],`$
where $`z=\alpha _s(m\nu )/\alpha _s(m)`$ and $`w=\alpha _s(m\nu ^2)/\alpha _s(m\nu )`$. The coefficients $`a_i`$ in Eq. (81) are
$`a_1`$ $`=`$ $`{\displaystyle \frac{32C_AC_F(C_A+2C_F)}{3\beta _0^2}},`$ (82)
$`a_2`$ $`=`$ $`{\displaystyle \frac{C_F[3\beta _0(26C_A^2+19C_AC_F32C_F^2)+C_A(208C_A^2+597C_AC_F+716C_F^2)]}{78\beta _0^2C_A}},`$ (83)
$`a_3`$ $`=`$ $`{\displaystyle \frac{C_F}{3\beta _0^2(6\beta _013C_A)(\beta _02C_A)}}\{2C_F^2(66\beta _0179C_A)(\beta _02C_A)`$ (86)
$`C_AC_F\left[6(493𝐒^\mathrm{𝟐})\beta _0^2(112639𝐒^\mathrm{𝟐})\beta _0C_A+1067C_A^2\right]`$
$`24C_A^2(6\beta _013C_A)(\beta _02C_A)\},`$
$`a_4`$ $`=`$ $`{\displaystyle \frac{24C_F^2(3\beta _011C_A)(5C_A+8C_F)}{13C_A(6\beta _013C_A)^2}},a_5={\displaystyle \frac{C_F^2\left[(4𝐒^\mathrm{𝟐}3)\beta _0+(1514𝐒^\mathrm{𝟐})C_A\right]}{6(\beta _02C_A)^2}},`$ (87)
$`a_6`$ $`=`$ $`{\displaystyle \frac{16C_F^2(C_A+2C_F)}{3\beta _0^2}},a_7={\displaystyle \frac{16C_AC_F(C_A+2C_F)}{\beta _0^2}}.`$ (88)
A further check on our result can be made by comparing it to the $`\alpha _s^3\mathrm{ln}^2(\alpha _s)`$ corrections calculated by Kniehl and Penin in Ref. . Near threshold the $`t\overline{t}`$ cross section depends on the product :
$`\sigma |c_1|^2G_C(0,0,E)|\psi _n^C(0)|^2,`$ (89)
where $`\psi ^C`$ is the leading order Coulomb wavefunction, and $`G_C`$ is the Coulomb Green’s function. In the approach used in Ref. , $`G_C`$ embodies corrections to the wavefunction at the origin, $`|\psi (0)|^2|\psi _n^C(0)|^2[1+\mathrm{\Delta }\psi ^2(0)]`$, and includes the large logarithms. They calculate the $`\alpha _s^2\mathrm{ln}(\alpha _s)`$ and $`\alpha _s^3\mathrm{ln}^2(\alpha _s)`$ terms and find:
$`\mathrm{\Delta }\psi ^2(0)`$ $`=`$ $`C_F\alpha _s^2\mathrm{ln}(\alpha _s)\left\{\left[2{\displaystyle \frac{2}{3}}𝐒^\mathrm{𝟐}\right]C_F+C_A\right\}`$ (91)
$`{\displaystyle \frac{C_F}{\pi }}\alpha _s^3\mathrm{ln}^2(\alpha _s)\left\{{\displaystyle \frac{3}{2}}C_F^2+\left[{\displaystyle \frac{41}{12}}{\displaystyle \frac{7}{12}}𝐒^\mathrm{𝟐}\right]C_FC_A+{\displaystyle \frac{2}{3}}C_A^2\right\},`$
for the terms not involving $`\beta _0`$. In our approach the large logarithms in the cross section all appear in the running coefficient $`c_1(\nu )`$, so we expect that the logarithms of Kniehl and Penin will be reproduced by
$`\mathrm{\Delta }\psi ^2(0)`$ $`=`$ $`\left|{\displaystyle \frac{c_1(\nu )}{c_1(1)}}\right|^21=\mathrm{\hspace{0.17em}2}\mathrm{ln}(\alpha _s)\gamma _{c_1}(1)+\mathrm{ln}^2(\alpha _s)\left\{\gamma _{c_1}^{}(1)+2[\gamma _{c_1}(1)]^2\right\}+\mathrm{},`$ (92)
where we have expanded to second order in $`\mathrm{ln}(\nu )=\mathrm{ln}(\alpha _s)`$. The $`[\gamma _{c_1}(1)]^2`$ term does not contribute at order $`\alpha _s^3\mathrm{ln}^2\alpha _s`$, and can be dropped. The remaining terms in Eq. (92) give:
$`\mathrm{\Delta }\psi ^2(0)`$ $`=`$ $`C_F\alpha _s^2\mathrm{ln}(\alpha _s)\left\{\left[2{\displaystyle \frac{2}{3}}𝐒^\mathrm{𝟐}\right]C_F+C_A\right\}`$ (93)
$``$ $`{\displaystyle \frac{C_F}{\pi }}\alpha _s^3\mathrm{ln}^2(\alpha _s)\left\{{\displaystyle \frac{3}{2}}C_F^2+\left[{\displaystyle \frac{41}{12}}{\displaystyle \frac{7}{12}}𝐒^\mathrm{𝟐}\right]C_FC_A+{\displaystyle \frac{2}{3}}C_A^2{\displaystyle \frac{\beta _0}{2}}\left[\left(2{\displaystyle \frac{2}{3}}𝐒^\mathrm{𝟐}\right)C_F+C_A\right]\right\},`$ (94)
which agrees exactly with the result from Ref. in Eq. (91) after setting $`\beta _0=0`$. Thus, we have shown that the logarithms of Kniehl and Penin can indeed be associated with renormalization group logarithms. Expanding Eq. (81) to higher orders gives the $`\alpha _s^4\mathrm{ln}^3(\alpha _s)`$, $`\alpha _s^5\mathrm{ln}^4(\alpha _s)`$, etc. terms.
For QCD with $`n_f=5`$, the coefficients in Eq. (82) are
$`a_1`$ $`=`$ $`4.113,a_2=2.173,a_3=4.3080.417𝐒^\mathrm{𝟐},a_4=5.731,`$ (95)
$`a_5`$ $`=`$ $`2.3471.209𝐒^\mathrm{𝟐},a_6=0.914,a_7=6.170.`$ (96)
The running coefficient $`c_1(\nu )`$ is shown in Fig. 7 for $`𝐒^\mathrm{𝟐}=2`$. For comparison we have also shown by a dotted line the value of the running coefficient obtained in Ref. by an approach based on simply taking $`\alpha _s\alpha _s(m\nu )`$ to approximate the NLL value for the coupling. This approximation misses many of the $`\alpha _s\mathrm{ln}(v)`$ terms, and does not provide a good estimate of the NLL result, as seen in Fig. 7. Some of the more important logarithms in our NLL result include the large $`\mathrm{ln}(m\nu ^2/m)`$ terms that enter through the mixing generated by ultrasoft gluon diagrams in the potential.
We have also shown the NNLO value for $`c_1`$ in Fig. 7. As pointed out in Ref. the NNLO value of $`c_1(1)`$ is fairly large. From Fig. 7 we see that summing the logarithms reduced the size of the NLO matching correction by a factor of two. It would be interesting to see if the running induced at NNLL continues to improve the convergence of the expansion. A consistent calculation at this order requires the three-loop anomalous dimension of $`c_1`$, as well as one loop running of the $`v^2`$ coefficients $`c_2`$ and $`c_3`$. This computation appears quite involved since running of the current at this order will likely depend on the running of higher order terms in the potential.
Several groups have analyzed the $`t\overline{t}`$ cross section predictions at NNLO using effective field theory techniques . It should be straightforward to incorporate the renormalization group improved current and potentials into their analysis by simply choosing a value of $`\nu `$ appropriate to the threshold region. The value of $`\nu `$ only needs to be of order the velocities in this region for the large logarithms to be minimized. Additional logarithms that appear in evaluating matrix elements with the potential will not involve large ratios of scales since they are of the form $`\mathrm{ln}(E/\mu _U)`$ and $`\mathrm{ln}(|𝐤|/\mu _S)`$ where $`\mu _U=m\nu ^2`$ and $`\mu _S=m\nu `$. In a more precise analysis one might wish to use $`\alpha _s(mv)C_F=v`$ to determine the value $`\nu `$ to use.
We would like to thank I. Rothstein and J. Soto for discussions. This work was supported in part by the Department of Energy under grant DOE-FG03-97ER40546, and by the National Science Foundation under NYI award PHY-9457911.
## A LL values for the order $`v^2`$ and $`v^0`$ potentials
The leading log values for the coefficients of the $`V^{(2)}`$ and $`V^{(0)}`$ potentials are :
$`𝒱_c^{(T)}(\nu )`$ $`=`$ $`4\pi \alpha _s(m)z,`$ (A1)
$`𝒱_c^{(1)}(\nu )`$ $`=`$ $`0,`$ (A2)
$`𝒱_r^{(T)}(\nu )`$ $`=`$ $`4\pi \alpha _s(m)z{\displaystyle \frac{32\pi C_A}{3\beta _0}}\alpha _s(m)\left[1z\right]{\displaystyle \frac{64\pi C_A}{3\beta _0}}\alpha _s(m)\mathrm{ln}\left(w\right),`$ (A3)
$`𝒱_2^{(T)}(\nu )`$ $`=`$ $`{\displaystyle \frac{\pi \left[C_A(352C_F+91C_d144C_A)3\beta _0(33C_A+32C_F)\right]}{39\beta _0C_A}}\alpha _s(m)\left[z1\right]`$ (A4)
$`+`$ $`{\displaystyle \frac{8\pi (3\beta _011C_A)(5C_A+8C_F)\alpha _s(m)}{13C_A(6\beta _013C_A)}}\left[z^{(113C_A/(6\beta _0))}1\right]`$ (A5)
$`+`$ $`{\displaystyle \frac{\pi (\beta _05C_A)\alpha _s(m)}{(\beta _02C_A)}}\left[z^{(12C_A/\beta _0)}1\right]{\displaystyle \frac{8\pi (4C_F+C_d3C_A)}{3\beta _0}}\alpha _s(m)\mathrm{ln}\left(w\right),`$ (A6)
$`𝒱_2^{(1)}(\nu )`$ $`=`$ $`{\displaystyle \frac{28\pi C_1}{3\beta _0}}\alpha _s(m)\left(1z\right)+{\displaystyle \frac{32\pi C_1}{3\beta _0}}\alpha _s(m)\mathrm{ln}\left(w\right),`$ (A7)
$`𝒱_s^{(T)}(\nu )`$ $`=`$ $`{\displaystyle \frac{2\pi \alpha _s(m)}{(2C_A\beta _0)}}\left[C_A+{\displaystyle \frac{1}{3}}(2\beta _07C_A)z^{(12C_A/\beta _0)}\right]+{\displaystyle \frac{1}{N_c}}\pi \alpha _s(m),`$ (A8)
$`𝒱_s^{(1)}(\nu )`$ $`=`$ $`{\displaystyle \frac{(N_c^21)}{2N_c^2}}\pi \alpha _s(m),`$ (A9)
$`𝒱_t^{(T)}(\nu )`$ $`=`$ $`{\displaystyle \frac{\pi \alpha _s(m)}{3}}z^{(12C_A/\beta _0)},`$ (A10)
$`𝒱_\mathrm{\Lambda }^{(T)}(\nu )`$ $`=`$ $`2\pi \alpha _s(m)\left[z4z^{(1C_A/\beta _0)}\right],`$ (A11)
where $`z=\alpha _s(m\nu )/\alpha (m)`$, $`w=\alpha _s(m\nu ^2)/\alpha (m\nu )`$, and we have included the $`N_c`$ dependent terms that come from matching the tree level annihilation diagrams. In the color singlet channel, $`𝒱_t(\nu )`$ and $`𝒱_\mathrm{\Lambda }(\nu )`$ were first calculated in Ref. and agree with Eq. (A1).
## B Potential operators that vanish on-shell
It is interesting to consider how the soft and ultrasoft anomalous dimensions for $`V^{(1)}`$ are affected when operators that vanish on-shell are included in the potential Lagrangian, $`_p`$. These new operators can have non-trivial anomalous dimensions. In this Appendix we consider two examples of how including these operators affects intermediate results, but in the end yield the same result as the on-shell potential for observables. Equivalently, it is shown that running the Lagrangian with the new operators from $`\nu =1`$ to $`\nu =v`$ and then removing them by field redefinitions or operator identities reproduces the on-shell running value of $`V^{(1)}`$.
Consider the scattering $`Q(p_1^0,𝐩)+\overline{Q}(p_2^0,𝐩)Q(p_3^0,𝐩^{})+\overline{Q}(p_4^0,𝐩^{})`$. As our first example, consider including in the $`v^0`$ potential in Eq. (14) terms of the form:
$`V^{(0)}`$ $`=`$ $`\left[𝒱_{\mathrm{\Delta }1}^{(T)}(T^A\overline{T}^A)+𝒱_{\mathrm{\Delta }1}^{(1)}(11)\right]{\displaystyle \frac{(p_3^0p_1^0)^2}{𝐤^4}}`$ (B1)
$`+`$ $`\left[𝒱_{\mathrm{\Delta }2}^{(T)}(T^A\overline{T}^A)+𝒱_{\mathrm{\Delta }2}^{(1)}(11)\right]{\displaystyle \frac{(𝐩^2p^2)^2}{4m^2𝐤^4}}.`$ (B2)
We will refer to this potential as an off-shell potential, since on-shell it vanishes by energy conservation, where $`p_3^0=p_1^0`$, $`𝐩_{}^{}{}_{}{}^{2}=𝐩^2`$. Define $`V_\mathrm{\Delta }^{(T,1)}=V_{\mathrm{\Delta }1}^{(T,1)}+V_{\mathrm{\Delta }2}^{(T,1)}`$. In Ref. it was shown that there is an operator identity whereby the time ordered product of a $`𝒱_c`$ and a $`𝒱_\mathrm{\Delta }`$ potential gives $`𝒱_k`$ potentials:
$`\text{}+\text{}`$ $`=`$ $`{\displaystyle \frac{i𝒱_c^{(T)}𝒱_\mathrm{\Delta }^{(T)}}{32mk}}T^AT^B\overline{T}^A\overline{T}^B+{\displaystyle \frac{i𝒱_c^{(T)}𝒱_\mathrm{\Delta }^{(1)}}{32mk}}T^A\overline{T}^A+\mathrm{},`$ (B3)
where the ellipses denote additional terms which vanish on-shell. The identity in Eq. (B3) allows us to remove the potential in Eq. (B1) at an arbitrary velocity scale $`\nu `$ in favor of a purely on-shell potential. Transforming $`𝒱_\mathrm{\Delta }^{(T,1)}(\nu )0`$ gives
$`𝒱_k^{(T)}(\nu )`$ $``$ $`𝒱_k^{(T)}(\nu )+{\displaystyle \frac{1}{32\pi ^2}}𝒱_c^{(T)}(\nu )\left[𝒱_\mathrm{\Delta }^{(1)}(\nu )+{\displaystyle \frac{1}{4}}(C_A+C_d)𝒱_\mathrm{\Delta }^{(T)}(\nu )\right],`$ (B5)
$`𝒱_k^{(1)}(\nu )`$ $``$ $`𝒱_k^{(1)}(\nu ){\displaystyle \frac{1}{32\pi ^2}}𝒱_c^{(T)}(\nu )C_1𝒱_\mathrm{\Delta }^{(T)}(\nu ).`$ (B6)
In a situation where $`𝒱_\mathrm{\Delta }^{(1)}(\nu )=0`$ transforming $`𝒱_\mathrm{\Delta }^{(T)}(\nu )0`$ induces a color singlet $`1/|𝐤|`$ potential of the form
$`𝒱_k^{(s)}(\nu )`$ $``$ $`𝒱_k^{(s)}(\nu ){\displaystyle \frac{C_F^2}{32\pi ^2}}𝒱_c^{(T)}(\nu )𝒱_\mathrm{\Delta }^{(T)}(\nu ).`$ (B7)
(Transformations between order $`1/v`$ and $`v^0`$ potentials were also considered in Refs. .)
To make the implications of Eq. (B5) more clear, we will consider an example and show that including the $`V_\mathrm{\Delta }^{(T,1)}`$ potential will not effect predictions for observables when running below $`m`$. Consider matching offshell in Feynman gauge where the matching coefficients at $`\nu =1`$ are $`V_\mathrm{\Delta }^{(T)}=4\pi \alpha _s(m)`$ and $`V_\mathrm{\Delta }^{(1)}=0`$, and the one loop anomalous dimensions for the potentials in Eq. (B1) are :
$`\nu {\displaystyle \frac{}{\nu }}𝒱_\mathrm{\Delta }^{(T)}`$ $`=`$ $`2\beta _0\alpha _s(m\nu )^2,`$ (B8)
$`\nu {\displaystyle \frac{}{\nu }}𝒱_\mathrm{\Delta }^{(1)}`$ $`=`$ $`0.`$ (B9)
This anomalous dimension has contributions from the one loop graphs with two soft gluons, ghosts or quarks where: the two factors of $`(𝐩^2p^2)`$ are from the soft vertices, or these factors are from two insertions on the soft propagators, or where one factor is from the propagator and one from the vertex (for the ghost loop). The solution of Eq. (B8) is $`V_\mathrm{\Delta }^{(T)}(\nu )=4\pi \alpha _s(m\nu )`$ and $`V_\mathrm{\Delta }^{(1)}(\nu )=0`$. The one loop counterterm which generates this running affects our soft anomalous dimension computation since now there are counterterms of the form
$`{\displaystyle \frac{(𝐩^{}{}_{}{}^{\mathrm{𝟐}}𝐩^\mathrm{𝟐})^2}{m^2(𝐩^{}𝐩)^\mathrm{𝟒}ϵ}}.`$ (B10)
These counterterm give rise to one loop diagrams which exactly cancel the divergences in Eq. (39) for the two loop graph in Fig. 3a. Recalling that the one loop matching values for the $`1/|𝐤|`$ potentials are also different , we find that the running $`𝒱_k^{(T,1)}(\nu )`$ coefficients in Eq. (60) become
$`𝒱_k^{(T)}(\nu )`$ $`=`$ $`{\displaystyle \frac{(3C_AC_d)}{4}}\alpha _s^2(m)+\left[{\displaystyle \frac{(3C_AC_d)}{4}}+{\displaystyle \frac{8C_A(C_A+C_d)}{3\beta _0}}\right](z^21)\alpha _s^2(m)`$ (B12)
$`+{\displaystyle \frac{8C_A(C_A+C_d)}{\beta _0}}\left[z12\mathrm{ln}(w)\right]\alpha _s^2(m),`$
$`𝒱_k^{(1)}(\nu )`$ $`=`$ $`C_1\alpha _s^2(m)+\left[C_1{\displaystyle \frac{16C_AC_1}{\beta _0}}\right](z^21)\alpha _s^2(m){\displaystyle \frac{48C_AC_1}{\beta _0}}\alpha _s^2(m)\left[z12\mathrm{ln}(w)\right],`$ (B13)
This Feynman gauge off-shell potential is consistent with the transformation in Eq. (B5) which transforms the off-shell case, $`V_\mathrm{\Delta }^{(T)}=4\pi \alpha _s(m\nu )`$ and Eq. (B12), into the on-shell case, $`V_\mathrm{\Delta }^{(T,1)}=0`$ and Eq. (60). The comparison in section IV A with Ref. shows that our on-shell analysis is also consistent with an off-shell Coulomb gauge potential.
With the off-shell potential in Eq. (B1) there are two new graphs which contribute to the anomalous dimension for the production current:
(B14)
and the anomalous dimension is therefore
$`\nu {\displaystyle \frac{}{\nu }}\mathrm{ln}[c_1(\nu )]`$ $`=`$ $`{\displaystyle \frac{𝒱_c^{(s)}(\nu )}{16\pi ^2}}\left({\displaystyle \frac{𝒱_c^{(s)}(\nu )}{4}}+𝒱_2^{(s)}(\nu )+𝒱_r^{(s)}(\nu )+𝐒^\mathrm{𝟐}𝒱_s^{(s)}(\nu )\right)+{\displaystyle \frac{𝒱_k^{(s)}(\nu )}{2}}{\displaystyle \frac{𝒱_c^{(s)}𝒱_\mathrm{\Delta }^{(s)}}{64\pi ^2}},`$ (B15)
where $`𝒱_\mathrm{\Delta }^{(s)}(\nu )=𝒱_\mathrm{\Delta }^{(1)}(\nu )C_F𝒱_\mathrm{\Delta }^{(T)}(\nu )`$. Since the last two terms in Eq. (B15) are invariant under the transformation in Eq. (B7) the off-shell potential gives the same prediction for the running of the production current. In calculating observables the operator identity in Eq. (B3) guarantees that the time ordered product of the running $`𝒱_\mathrm{\Delta }^{(T)}`$ and $`𝒱_c^{(T)}`$ produces the same effect as the soft contribution to the running of the on-shell $`𝒱_k^{(T)}`$ in Eq. (39).
As our second example, consider including an operator which would vanish by the lowest order free equations of motion:
$`_{\mathrm{eom}}={\displaystyle \underset{𝐩,𝐩^{}}{}}{\displaystyle \frac{𝒱_F}{(𝐩^{}p)^2}}\left[\psi _𝐩^{}^{}T^A\left(i_0{\displaystyle \frac{𝐩^\mathrm{𝟐}}{2m}}\right)\psi _𝐩\right]\chi _𝐩^{}^{}\overline{T}^A\chi _𝐩+\chi \psi .`$ (B17)
The only feature of the operator in Eq. (B17) that is different from typical operators that vanish by the equations of motion (for example, in HQET) is its non-local nature relative to the scales $`𝐩𝐩^{}mv`$. At one loop there are ultrasoft gluon graphs which mix into the operator in Eq. (B17), for instance from the diagrams:
(B18)
In the first graph the vertex is the first term in Eq. (25) and the ultrasoft gluon is $`𝐀`$, while in the second diagram the cross is a insertion of the ultrasoft $`^2/m`$ operator in Eq. (10) and the gluon is an $`A^0`$. In Feynman gauge both graphs produce a divergence of the form
$`{\displaystyle \frac{(E𝐩^\mathrm{𝟐}/m)}{ϵ(𝐩^{}𝐩)^2}},`$ (B19)
which vanishes by the equations of motion, but off-shell renormalizes the coefficient of the operator in Eq. (B17). To see how this effects the calculation of our ultrasoft anomalous dimension, consider Coulomb gauge, where only Eq. (B18a) is non-zero. In this case there is a one loop counterterm diagram which exactly cancels the divergence in Fig. (5b). After running down to the low scale the value of $`𝒱_k^{(T,1)}(\nu )`$ are different, but our potential also includes the operator in Eq. (B17) with coefficient $`𝒱_F(\nu )`$. If we wish to remove the operator in Eq. (B17) we can do so with a field redefinition
$`\psi _𝐩^{}\psi _𝐩^{}+{\displaystyle \underset{𝐩^{\prime \prime }}{}}{\displaystyle \frac{𝒱_F}{(𝐩^{}𝐩)^2}}\psi _{𝐩^{\prime \prime }}^{}T^A\chi _{𝐩^{\prime \prime }}^{}\overline{T}^A\chi _𝐩,`$ (B20)
which induces a term from the $`\psi _𝐩^{}[i_0𝐩^\mathrm{𝟐}/(2m)]\psi _𝐩`$ Lagrangian which cancels Eq. (B17). This field redefinition also gives other new contributions. For us the important point is that the Coulomb potential induces a six-quark operator of the form:
$`{\displaystyle \underset{𝐩,𝐩^{}}{}}{\displaystyle \underset{𝐩^{\prime \prime }}{}}\psi _{𝐩^{\prime \prime }}^{}\chi _{𝐩^{\prime \prime }}^{}\chi _𝐩^{}\psi _𝐩\chi _𝐩^{}^{}\chi _𝐩{\displaystyle \frac{𝒱_F𝒱_c}{(𝐩^{\prime \prime }𝐩^{})^2(𝐩^{}𝐩)^2}}.`$ (B21)
(The contraction of color indices and factors of $`T^A`$ have been suppressed.) Usually a six quark operator could not possibly effect the running of a four quark operator such as the $`1/|𝐤|`$ potential. However, because of the momentum dependence in the denominator of Eq. (B21) it induces a non-zero tadpole diagram where the fields $`\chi _𝐩^{}`$ and $`\chi _𝐩^{}^{}`$ are contracted. This tadpole graph produces a $`1/|𝐤|`$ and makes up for the running in $`𝒱_k`$ that was removed when the contribution from Eq. (5b) was canceled by a $`𝒱_F`$ counterterm.
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# Some remarks on separability of states
## I Introduction
In the analysis of quantum mechanical properties of composite systems there is a strong demand for a characterization of various subsets of states. Among properties of states, perhaps the most intriguing phenomenon is that of entanglement. We recall that, in physical terms, this notion reflects inability of local preparation of the state. Probably, the most famous example of states with entanglement is that with violation of Bell inequalities (cf. ). Other nice applications can be found in the theory of decoherence or in quantum information theory. This makes clear why one of the problems in foundations of quantum theory is a classification and characterization of the entangled quantum states. In this context it has been shown (cf. , , ) that there should be a strong connection between classification of entanglement of quantum states and the structure of positive maps. The problem is complicated since it is necessary to take into account some basic facts concerning specific properties of tensor product of algebras. Unfortunately, such approach involves mathematically advanced concepts and therefore it is not easily accessible to most physicists working on applications of Quantum Mechanics to Quantum Information or Quantum Optics. Moreover, the theory of positive linear maps of $`C^{}`$ algebras appears in the physical literature in a rather scattered and usually special form, the best known examples are states, dynamical maps, and -representations.
The purpose of this paper is to present, in plain words, state-of-the-art of the relevant parts of theory of positive maps as well as to clarify the present day characterization of entanglement. The structure of the paper is the following: in Section II we review some results in the theory of positive maps and show how these facts can be applied to a characterization of some singled out subsets of states. Section III is devoted to indicate that the family of separable states is really very small. Finally, in Section IV, some remarks concerning recent results on positive maps will be given. Moreover, we argue that the problems associated to the separability follow directly from the rule saying that a composite system is described by tensor product.
Finally, we want to emphasize that we address this paper to physicists working on application of Quantum Mechanics. Therefore, the presented material does not include neither general results nor can be considered as a full account of all mathematical aspects of the theory. However, all results are rigorous.
## II Application of positive maps to the description of separability
We consider a composite system $`𝒜+`$ consisting of two subsystems $`𝒜`$ and $``$ respectively. The algebra of all observables is given by $`(_{𝒜+})`$, $`(_𝒜)`$, and $`(_{})`$, for $`𝒜+`$, $`𝒜`$, and $``$ respectively where $`()`$ stands for the set of all linear bounded operators on a Hilbert space $``$. Moreover, $`_{𝒜+}`$ will be identified with $`_𝒜_{}`$. We note that the particular case $`dim_𝒜=n<\mathrm{}`$, $`dim_{}=m<\mathrm{}`$ (dim stands for dimension) is the standard assumption in an analysis of entanglement states in Quantum Information. Recall that a state (see ) of a system, say $`𝒜`$, is defined as a linear, complex-valued map $`\varphi :(_𝒜)\text{ }\text{ }\mathrm{C}`$ such that $`\varphi (A^{}A)0`$, $`A(_𝒜)`$, and $`\varphi (\widehat{1\mathrm{l}})=1`$ ($`\widehat{1\mathrm{l}}`$ is the identity operator in $`()`$). Usually, in physical literature, a state $`\varphi `$ is assumed to be of the form $`\varphi (A)=\mathrm{Tr}\varrho A`$ where $`\mathrm{Tr}`$ stands for the trace, $`\varrho `$ for a density matrix. This fact is used frequently as a legitimacy for identification of a state with a density matrix. We denote the family of all states of $`𝒜`$ ( $``$, $`𝒜+`$) by $`𝒮_𝒜`$ ($`𝒮_{}`$, $`𝒮_{𝒜+}`$ respectively). Finally, we will need two definitions of positivity. Let $`\mathrm{\Psi }:(_1)(_2)`$ be a linear map ($`_i`$ is a Hilbert space). $`\mathrm{\Psi }`$ is positive when $`\mathrm{\Psi }:(_1)^+(_2)^+`$ where $`(_i)^+`$ denotes all positive operators in $`(_i)`$, $`i=1,2`$. Now let $`id_k`$ be the identity map on $`M_k(\text{ }\text{ }\mathrm{C})`$ (set of all $`k\times k`$ matrices). We define the map $`id_k\mathrm{\Psi }:M_k(\text{ }\text{ }\mathrm{C})()M_k(\text{ }\text{ }\mathrm{C})()`$, for $`k=1,2,..`$ by
$$(id_k\mathrm{\Psi })(\underset{i}{}\sigma _iA_i)=\underset{i}{}\sigma _i\mathrm{\Psi }(A_i),$$
(1)
where $`\sigma _iM_k(\text{ }\text{ }\mathrm{C})`$ and $`A_i()`$. The map $`\mathrm{\Psi }`$ is $`k`$-positive when $`id_k\mathrm{\Psi }`$ is positive. The map $`\mathrm{\Psi }`$ is completely positive when $`\mathrm{\Psi }`$ is $`k`$-positive for all $`k=1,2,\mathrm{}`$
Having the defined notion of state we can ask how $`𝒮_{𝒜+}`$ is related to $`𝒮_𝒜`$ and $`𝒮_{}`$. This is rather deep question and we start its analysis with the following simple exercise (cf. , 11.5.11): let $`dim_𝒜=2=dim_{}`$. Define the state $`\phi `$ in $`𝒮_{𝒜+}`$ as $`\phi (AB)=\frac{1}{2}((AB)(e_1f_1+e_2f_2),e_1f_1+e_2f_2)`$ where $`\{e_i\}_{i=1}^2`$ is a basis in $`_𝒜`$, $`\{f_i\}_{i=1}^2`$ a basis in $`_{}`$. In other words, $`\phi ()`$ is a state determined by the vector $`\frac{1}{\sqrt{2}}(e_1f_1+e_2f_2)_{𝒜+}`$. We denote by $`\phi _0`$ a convex combination of product states, i.e.
$$\phi _0=a_n\omega _n^𝒜\omega _n^{}$$
(2)
where $`\omega _n^𝒜(A)=(Ax_n,x_n)`$ for $`A(_𝒜)`$, $`x_n_𝒜`$ and analogously for $`\omega _n^{}`$. $`\{a_n\}`$ are non-negative numbers such that $`_na_n=1`$. Then one can show (see () that
$$\phi \phi _0\frac{1}{4}.$$
(3)
This clearly shows that the state $`\phi `$ of $`(_𝒜)(_{})`$ lies at norm distance at least $`\frac{1}{4}`$ from the convex hull of the set of product states. Consequently, the family of convex combinations of product states is a proper subset of the set of all states on $`(_𝒜)(_{})`$ and moreover this family is not a norm dense subset in $`𝒮_{𝒜+}`$, i.e. , in general, arbitrary state can not be approximated in norm by convex combination of product states. This result shows the crucial property of the space of states of the composite system. Therefore we define
###### Definition 1
A state $`\phi `$ defined on $`(_𝒜)(_{})`$ is said to be separable if it can be written as
$$\phi =\underset{n}{}a_n\phi _n=\underset{n}{}a_n\phi _n^𝒜\phi _n^{}$$
(4)
where $`\phi _n=\phi _n^𝒜\phi _n^{}`$ with $`\phi _n^𝒜`$ ($`\phi _n^{}`$) a state on $`(_𝒜)`$ ($`(_{})`$), and $`a_n`$ are positive numbers that sum to one. The set of all separable states , for the system $`𝒜+`$ will be denoted by $`𝒮_{𝒜+}^{sep}`$. A state which can not be written as in (4) is called entangled or equivalently non-separable.
Peres () showed that a necessary condition for a density matrix of a bipartite system to be separable is for its partial transpose to be a density matrix where the partial transpose is just the transposition of the matrix representing the state of either subsystem $`𝒜`$ or subsystem $``$. For $`22`$ and $`23`$ systems this condition is also sufficient (). This is the point where a part of theory of positive maps enters to the problem. To make this point more clear we restrict ourselves, for a moment, to finite dimensional case (so $`dim_𝒜=n<\mathrm{}`$, $`(_𝒜)`$ can be taken as $`M_n(\text{ }\text{ }\mathrm{C})M_n`$ and the same for subsystem $``$) and we recall the following scheme which is taken from Choi’s lecture ():
| $`L(M_n,M_k)`$ | $``$ | $`M_nM_k`$ | $``$ | linear functionals on $`M_nM_k`$ |
| --- | --- | --- | --- | --- |
| $``$ | | | | $``$ |
| hermitian preserving linear maps | $``$ | $`(M_nM_k)^\mathrm{h}`$ | $``$ | linear functionals assuming real values on $`(M_nM_k)^\mathrm{h}`$ |
| $``$ | | | | $``$ |
| positive linear maps | $``$ | ? | $``$ | linear functionals assuming positive values on $`M_n^+M_k^+`$ |
| $``$ | | | | $``$ |
| completely positive linear maps | $``$ | $`(M_nM_k)^+`$ | $``$ | linear functionals assuming positive values on $`(M_nM_k)^+`$ |
where $``$ denotes a $`11`$ correspondence, $``$ is an inclusion, $`M_n^h`$ stands for all hermitian matrices in $`M_n`$, $`M_n^+M_k^+=\{A_iB_i,A_iM_n^+,B_i^+M_k^+\}`$, while $`L(M_n,M_k)`$ stands for linear maps from $`M_n`$ to $`M_k`$.
###### Remark 1
Even this finite dimensional case clearly shows strong relations among: states, some singled out subsets of states, positive maps, and finally completely positive maps. This can be taken as a hint that the core of the considered problem is related to some special properties of the order (defined by positive elements) in tensor products. Moreover, one can expect that definition of order in tensor products leads to some serious problems. This is so, and for a general review of that question we refer Wittstock’s lecture in (). Later on, in the next Section, we come back to that point.
Choi’s scheme gives the following simple fact relating separable states with the order in tensor product. Let $`\phi `$ be in $`𝒮_{𝒜+}^{sep}`$, then (we recall that $`M_n^+M_k^+(M_nM_k)^+`$ where $``$ is the proper inclusion)
$$AM_n^+M_k^+implies\phi (A)0.$$
(5)
Also,
$$\phi (\mathrm{\Psi }_n\mathrm{\Psi }_k)|_{M_n^+M_k^+}0$$
(6)
for $`\mathrm{\Psi }_n`$ ($`\mathrm{\Psi }_k`$) positive maps on $`M_n(\text{ }\text{ }\mathrm{C})`$ ($`M_k(\text{ }\text{ }\mathrm{C})`$ respectively). Obviously, the same holds if one replaces $`M_n`$ by $`(_𝒜)`$ with not necessary $`dim_𝒜<\mathrm{}`$, etc.
Thus Choi’s scheme leads to the origin of Peres-Horodeckis characterization of separable states. Namely, let us recall that the transposition map $`\tau :a_{ik}\tau (a_{ik})=a_{ki}`$ is a positive map but not even two-positive, so not completely positive (cf. ). Take a state $`\phi 𝒮_{𝒜+}`$ and $`\psi 𝒮_{𝒜+}^{sep}`$. Denote by $`\tau ^{}`$ the transposition map acting on the set of observables associated with the subsystem $``$. Then, according to Choi’s scheme $`\phi id^𝒜\tau ^{}`$ is a linear functional taking positive values on $`M_n^+M_k^+`$ but not in general on $`(M_nM_k)^+`$. Therefore, separable states can differentiate positivity from completely positivity. This idea was used in
###### Theorem 1
(Horodecki’s. ) A density matrix $`\varrho `$ on $`_𝒜_{}`$ is entangled iff there exists a positive linear map $`𝒮:_𝒜_{}`$ such that
$$(id_𝒜𝒮)\varrho $$
(7)
is not positive semidefinite. Here $`id_𝒜`$ denotes the identity map on $`(_𝒜)`$.
Let us turn again to Choi’s scheme. To complete partly the scheme we recall that one can endow the set of all matrices, say $`n\times n`$, with the so called Hadamard product (cf. ). Let $`A,BM_n(\text{ }\text{ }\mathrm{C})`$, then the matrix $`H=ABM_n(\text{ }\text{ }\mathrm{C})`$ is called Hadamard product of $`A`$ and $`B`$ if its $`(i,j)`$ entry is equal to $`a_{ij}b_{ij}`$. Let $`\mathrm{\Psi }:M_nM_n`$ be a positive map, then $`A\mathrm{\Psi }(A)`$ is also positive for all $`AM_n^+`$. Conversely, for a fixed $`AM_n^+`$ the map $`S_A:M_nM_n`$ defined by
$$S_A(B)=AB.$$
(8)
leads (as it will be shown) to a positive map. Thus we get a nice correspondence:
positive maps in $`L(M_n,M_n)`$ $``$ $`M_n^+M_n^+`$.
To see the positivity of $`S_A`$ and to understand this result in the context of separability problem let us reconsider another Kadison-Ringrose exercise (cf. 2.8.41). Namely, let $``$ be a Hilbert space (not necessary of finite dimension) with an orthonormal basis $`\{e_1,e_2,\mathrm{}\}`$. We denote by $`[a_{ij}]`$ and $`[b_{ij}]`$ the matrices of $`A,B()`$ with respect to the basis $`\{e_i\}`$. Next, let us consider the orthonormal system $`\{e_ie_i\}_i`$ in $``$ and we denote by $`𝒦`$ the subspace of $``$ spanned by the system $`\{e_ie_i\}`$. Finally, the projection from $``$ onto $`𝒦`$ will be denoted by $`P`$. Then solving the Kadison-Ringrose exercise 2.8.41 we find that the matrix elements $`\{t_{ij}\}`$ of the operator $`(P(AB)P|_𝒦)T`$ with respect to the basis $`\{e_ie_i\}`$ of $`𝒦`$ are
$$t_{ij}=a_{ij}b_{ij}([a_{ij}][b_{ij}])_{ij}.$$
(9)
Clearly, $`\{t_{ij}\}0`$ provided that $`\{a_{ij}\}0`$ and $`\{b_{ij}\}0`$.
Now, let $`A`$ be a density matrix $`\varrho `$, $`B`$ a density matrix $`\sigma `$, both acting on a Hilbert space $``$. Then, dropping for simplicity the normalizating factor, we got a hint that there are density (nonseparable) matrices satisfying Peres condition (note that $`P`$ is not in $`()^+()^+`$; cf. , exercise 11.5.10).
## III How generic separable states are?
Among questions concerning the separability, great effort has been done to determine degree of largeness of the family of separable states in the set of all states (cf. , ). We have seen, in Section II, that $`𝒮^{sep}`$ is not a norm-dense subset of $`𝒮`$. However, it appear highly desirable to study this question with a concept of convergence which would be more accessible to experimental verifications. Namely, from physical point of view, it would be reasonable to say that a sequence of separable states $`\{\phi _n\}`$ approximates a state $`\phi `$ if, given any observable $`A`$ in $`()`$ and any $`ϵ>0`$ one can find an integer $`N(A,ϵ)`$ such that
$$|\phi _n(A)\phi (A)|<ϵforallnN(A,ϵ).$$
(10)
In mathematical terms, it would mean: is the set $`𝒮^{sep}`$ weakly -dense in $`𝒮`$? In Section II it was shown that a description of $`𝒮^{sep}`$ (as well as of $`𝒮`$) is closely related to distinguished subsets of positive observables in the tensor product $`(_𝒜)(_{})`$. So to solve the just posed question we should have a characterization of positiveness in the tensor products. The desired result follows from old Tomiyama theorem (see ). To quote his result, unfortunately, some additional definitions are necessary. We define (cf. ) n-state of $`()`$ to be a matrix $`[\vartheta _{ij}]`$ of linear functionals on $`()`$ such that the matrix $`[\vartheta _{ij}(A_{ij})]`$ is semi-positive defined when $`[A_{ij}]M_n(())^+`$ (so the matrix $`[A_{ij}]`$ having entries in $`()`$, is also nonnegative) and $`\vartheta _{ii}(\widehat{1\mathrm{l}})=1`$ for $`i=1,2,..`$. If the normalization is omitted we get the notion of n-positive functional.
To get an example of such $`n`$-state let us pick $`\{x_1,x_2,\mathrm{}\}`$ a set of unit vectors in $``$. Denote by $`\omega _{x_i,x_j}`$ the following functional $`\omega _{x_i,x_j}(A)=(Ax_i,x_j)`$, $`A()`$. Then, one can easily check that $`\{\omega _{x_i,x_j}\}`$ is an $`n`$-state on $`()`$.
###### Remark 2
We have seen, in Section II, the great significance of complete positive maps in the analysis of separable states. In that context the following observation seems to be important: a linear map $`\mathrm{\Phi }:(_1)(_2)`$, such that $`\mathrm{\Phi }(\widehat{1\mathrm{l}})=\widehat{1\mathrm{l}}`$, is completely positive if and only if $`\{\vartheta _{ik}\mathrm{\Phi }\}`$ is n-state for each n-state $`\varrho _{ik}`$ of $`(_2)`$.
Let $`\phi `$ be a functional on $`(_1)(_2)`$. $`\phi `$ is called a positive functional of order $`n`$ if $`\phi `$ is expressed as
$$\phi =\underset{i,j=1}{\overset{n}{}}\varphi _{ij}\psi _{ij}$$
(11)
where $`\{\varphi _{ij}\}`$ and $`\{\psi _{ij}\}`$ are $`n`$-positive functionals on $`(_1)`$ and $`(_2)`$ respectively.
Now, in this setting, the Tomiyama theorem asserts that the order in $`(_1)(_2)`$ (so the distinguished family of positive observables of the composite system) is determined by the set of positive functionals of order n where n is equal to $`min\{dim_1,dim_2\}`$.
Now to get the promised estimation of degree of largeness of $`𝒮^{sep}`$ we have to recall another Kadison’s result. Namely, a family of states which determines the order is weak dense and conversely a weak dense family of states determines the order in the set of observables. Combining Kadison and Tomiyama results one can see that $`𝒮^{sep}`$ can not be weakly dense subset in $`𝒮`$ unless one of the subsystems is one dimensional (so classical one, cf. also ).
To make more clear the above conclusion let us consider the favorite model in Quantum Information. We assume that, say, $`dim_1=2`$, $`dim_2`$ can be arbitrary. Then the order of $`(_1)(_2)`$ is determined by functionals of the form
$$\phi =\underset{i,j=1}{\overset{2}{}}\varphi _{i,j}\psi _{ij}.$$
(12)
Take for $`\{\varphi _{ij}\}`$ the 2-state with $`\varphi _{ij}`$ being the same state $`\eta _1`$ of $`(_1)`$ and analogously for $`\{\psi _{ij}\}`$. Then, it is clear that the family $`𝒮_0`$ determining the order in the tensor product contains as a proper subset the family of separable states, $`𝒮^{sep}𝒮_0`$ but not $`𝒮^{sep}𝒮_0`$. Therefore, in general, one can not approximate a state $`\phi `$ by separable states in a way which is accessible to experimental verifications. By the way, the same argument leads to similar conclusion for the corresponding property of entangled states.
There is a remarkable relation between $`k`$-functionals and a certain class of states with Schmidt rank k which were considered recently (see for example ). Namely, let $`|\psi >=_{i=1}^k\sqrt{\lambda _i}|x_i>|y_i>`$ where $`\{|x_i>\}`$ and $`\{|y_i>\}`$ are orthonormal systems. We observe
$`Tr(|\psi ><\psi |ab)=<\psi |ab|\psi >={\displaystyle \underset{ij}{}}\sqrt{\lambda _i\lambda _j}\omega _{x_i,x_j}(a)\omega _{y_i,y_j}(b)`$ (13)
$`={\displaystyle \sqrt{\lambda _i\lambda _j}(\omega _{x_i,x_j}\omega _{y_i,y_j})(ab)}.`$ (14)
Clearly, $`Tr\{|\psi ><\psi |\}`$ is a positive functional of order $`k`$ while $`\{\omega _{x_i,x_j}\}`$ and $`\{\omega _{y_i,y_j}\}`$ are $`k`$-states. Let $`\mathrm{\Phi }`$ be a positive normalized map and $`\omega _{x_i,x_j}`$ a $`k`$-state. Then it is well known that $`\mathrm{\Phi }`$ is $`k`$-positive iff $`\omega _{x_i,x_j}\mathrm{\Phi }`$ is again $`k`$-state (cf. Remark 2). This implies that
$$Tr(|\psi ><\psi |a\mathrm{\Phi }(b)$$
(15)
is another positive functional of order $`k`$ if $`\mathrm{\Phi }`$ is $`k`$-positive. Consequently, this remark and Kadison and Tomiyama results provides the general context of the recent characterization of $`k`$-positivity given in ().
To get another approximation procedure as well as another characterization of separable states one can proceed as follows. Let $`\varrho `$ be a density matrix on $`_1_2`$ and let $`\{e_i\}_i`$ be a basis in $`_2`$. Define a linear isometry $`U_i:_1_1_2`$ by $`U_ix=xe_i^i`$ for $`x_1`$ where $`\{^i\}`$ is a family of pairwise orthogonal subspaces in $`_1_2`$. Then $`U_i^{}`$ is a linear map of $`_1_2`$ in $`_1`$ such that $`U_k^{}(_1_2^l)=0`$ for $`kl`$. For $`\varrho (_1_2)`$ we define
$$\widehat{\varrho _{ik}}=U_i^{}\varrho U_k(_1).$$
(16)
In other words, the density matrix $`\varrho `$ can be represented as a semipositive matrix $`\{\widehat{\varrho _{ik}}\}`$ with operator-valued entries.
Assume $`\varrho `$ to be of the form $`\varrho =\varrho _1\varrho _2`$ with $`\varrho _i(_i)`$, $`i=1,2.`$ The the corresponding matrix has the very special form:
$$[\widehat{\varrho _{ik}}]=[\lambda _{ij}\varrho _1]$$
(17)
where $`[\lambda _{ik}][(e_i,\varrho _2e_k)]`$ is semipositive defined matrix with $`\mathrm{C}`$-valued entries, $`\varrho _1`$ is a density matrix in $`(_1)`$. Therefore, such matrix has commutative entries. Furthermore, any separable density matrix is just a convex combination of matrices of such special form.
Now we in position to study another approximation procedure. Assume $`\varrho `$ can be approximated by separable density matrices. It would mean that, in just described matrix representation, one can find a family of semipositive defined matrices $`[\lambda _{ik}^\alpha ]`$ with $`\mathrm{C}`$-valued entries, a family of density matrices $`\{\varrho ^\alpha \}`$ (on the Hilbert space associated with the arbitrary but fixed subsystem) such that
$$\widehat{\varrho _{ik}}\underset{\alpha }{}\lambda _{ik}^\alpha \varrho _1^\alpha $$
(18)
is small. Clearly, for high dimension this is a hopeless task. For low dimensions, for some subsets of states and (or) for some finite accuracy one can get such approximation.
## IV General positive maps and concluding remarks
In previous Sections, general positive maps are frequently used for an analysis of separable and non-separable states. On the other hand, completely positive maps are admittedly the ”physical” ones (cf. ). So one can suspect us for using ”non-physical” tools for a description of some singled out subsets of ”physical” states. To argue that this is not the case we again use another Kadison result (cf. ). Namely, let us restrict ourselves to Hamiltonian flows and consider equivalence between Schrödinger and Heisenberg picture. Physical maps sending states into states can be taken to be affine (it takes into account the superposition principle) bijections (this property reflects the reversibility of Hamiltonian dynamics). Their counterparts in Heisenberg picture are linear maps, preserving hermitian conjugation and anticommutators (in mathematical terms they are Jordan morphisms). But any Jordan morphism can be splitted into -homomorphism (so a completely positive map) and -anti-homomorphism (only positive, not even two positive map, e.g. the transposition $`\tau `$ used in Section II is an example of anti-homomorphism). So this result shows that there is a room for physical maps which are not completetely positive.
Having a motivation for positiveness we want to present a general construction of a very general linear, positive map $`S:(_𝒜)(_{})`$ (cf. and ). For simplicity we assume that $`_𝒜`$ and $`_{}`$ are finite dimensional Hilbert spaces. However, nearly all results can be straightforwardly generalized to infinite dimensional case. Let us take a hermitian operator $`H(_𝒜_{})`$ such that
$$H\{A(_𝒜_{});\varphi (A)0forall\varphi 𝒮_{𝒜+}^{sep}\}𝒮_{sep}^d.$$
(19)
and we assume that $`H(_𝒜_{})^+`$. This is always possible since $`(_𝒜)^+(_{})^+`$ ($`𝒮_{𝒜+}^{sep}`$) is a proper subset of $`(_𝒜_{})^+`$ ($`𝒮_{𝒜+}`$ respectively). For some further explanation an an example of such a choice see (, another example can be found in (). In fact, here, we are partly following Terhal and Takasaki-Tomiyama idea (, ).
Let $`\{e_i\}`$ be an orthonormal basis in $`_𝒜`$, $`z,x`$ arbitrary vectors in $`_𝒜`$, $`v,y`$ arbitrary vectors in $`_{}`$. Let us observe (we take the scalar product to be linear in the second factor)
$`(zv,Hxy)={\displaystyle \underset{ij}{}}\overline{(e_i,z)}(e_j,x)(e_iv,He_jy)={\displaystyle \underset{ij}{}}(z,E_{ij}x)(v,V_i^{}HV_jy)`$ (20)
$`=(zv,\left({\displaystyle \underset{ij}{}}E_{ij}V_i^{}HV_j\right)xy)`$ (21)
where $`E_{ij}|e_i><e_j|`$, $`V_iy=e_iy`$ (so $`V_i`$ is defined analogously to $`U_i`$, see definition given prior to (16)). Define the map $`S`$
$$S(E_{ij})=V_i^{}HV_j.$$
(22)
Let us note
$$S(E_{ij})^{}=(V_i^{}HV_j)^{}=V_j^{}HV_i=S(E_{ji})=S(E_{ij}^{}).$$
(23)
Therefore to prove positivity of the map $`S`$ it is enough to show that $`S`$ maps any projector $`|f><f|`$ ($`f_𝒜`$) into a positive operator. To this end we note
$`(y,S\left(|f><f|\right)y)=(y,{\displaystyle \underset{ij}{}}\overline{\lambda _i}\lambda _jS(E_{ij})y)={\displaystyle \underset{ij}{}}\overline{\lambda _i}\lambda _j(y,V_i^{}HV_jy)`$ (24)
$`{\displaystyle \underset{ij}{}}(\lambda _ie_iy,H\lambda _je_jy)=(fy,Hfy)0`$ (25)
where the last inequality follows from (19). Consequently, $`S`$ is a positive map. Further we note that (20) implies
$$H=\underset{ij}{}E_{ij}S(E_{ij})=\underset{ij}{}(\widehat{1\mathrm{l}}_𝒜S)(E_{ij}E_{ij}).$$
(26)
Let us assume additionally that $`_𝒜_{}`$. Then
$`Tr__𝒜_{}(\mu H)={\displaystyle \underset{ij}{}}Tr__𝒜_{}\left(\mu (\widehat{1\mathrm{l}}_𝒜S)(E_{ij}E_{ij})\right)`$ (27)
$`={\displaystyle \underset{ij}{}}Tr__𝒜_{}\left((\widehat{1\mathrm{l}}_𝒜S^{})(\mu )E_{ij}E_{ij}\right)=Tr__𝒜_{}\left((\widehat{1\mathrm{l}}_𝒜S^{})(\mu )P_{_ie_ie_i}\right)`$ (28)
where $`P_{_ie_ie_i}`$ is the projector onto $`_ie_ie_i`$. Since $`H(_𝒜_{})^+`$ then one can find a density matrix, say, $`\mu `$ such that $`Tr__𝒜_{}(\mu H)<0`$. Thus
$$(\underset{i}{}e_ie_i,(\widehat{1\mathrm{l}}_𝒜S^{})(\mu )\underset{i}{}e_ie_i)<0.$$
(29)
Consequently, $`S^{}`$ so the map $`S`$ is not a completely positive map. We want to emphasize that this result stems from the fact that $`𝒮^{sep}`$ is a proper subset of the set of all states $`𝒮`$; see a remark prior to definition of separable states. Now we want to show that the above construction provides positive maps which are not decomposable ones. We recall that a linear positive map $`S`$ is decomposable if it can be written as
$$S=S_1+S_2\tau $$
(30)
where $`S_1,S_2`$ are completely positive maps while $`\tau `$ is a transposition. To prove the claim let us take $`H𝒮_{sep}^d`$, $`H^{}=H`$ and a state $`\varrho 𝒮`$ such that
$$TrH\varrho <0$$
(31)
and $`(\widehat{1\mathrm{l}}\tau )\varrho `$ is semipositive defined operator. Here we assume that the Hilbert spaces are of dimension larger than 2. Namely, it can be shown that for the case $`dim_1=2`$ and $`dim_2=n`$, $`n2`$, it is impossible to find density matrix $`\varrho `$ such that the assumed conditions are satisfied. For Hilbert spaces of higher dimension such a choice is possible, for an example see ().
Now we want to argue that the condition: for a fixed state $`\varrho `$
$$Tr((\widehat{1\mathrm{l}}\tau )\varrho )B)0foranyB(_𝒜_{})^+$$
(32)
implies
$$Tr\left((\widehat{1\mathrm{l}}\tau T)(\varrho )B\right)0foranycompletelypositivemapTandB(_𝒜_{})^+.$$
(33)
To show this we observe $`B(_𝒜_{})^+`$ if and only if $`B=C^{}C`$ with $`C(_𝒜_{})`$. So $`C=_iE_iF_i`$ and $`B=_{ij}E_i^{}E_jF_i^{}F_j`$. But (32) implies
$$\underset{ij}{}Tr\left(\varrho (\widehat{1\mathrm{l}}\tau ^{})(E_i^{}E_jF_i^{}F_j)\right)0.$$
(34)
Thus
$$\underset{ij}{}Tr\left(\varrho E_i^{}E_j\tau ^{}(F_i^{}F_j)\right)0.$$
(35)
On the other hand (33) can be rewritten as
$`{\displaystyle \underset{ij}{}}Tr((\widehat{1\mathrm{l}}T)(\varrho )E_i^{}E_j\tau ^{}(F_i^{}F_j)={\displaystyle \underset{m}{}}{\displaystyle \underset{ij}{}}Tr(\varrho E_i^{}E_jV_m\tau ^{}(F_i^{}F_j)V_m^{})`$ (36)
$`={\displaystyle \underset{m}{}}{\displaystyle \underset{ij}{}}Tr\left(\varrho E_i^{}E_j\tau ^{}(\tau ^{}(V_m)F_i^{}F_j\tau ^{}(V_m))\right)`$ (37)
where we used the fact that the transposition $`\tau `$ is antihomomorphism and the fact that any completely positive map $`T(A)`$ can be written as $`_mV_m^{}AV_m`$ (cf. ). But (35) implies
$$\underset{ij}{}Tr\left(\varrho E_i^{}E_j\tau ^{}(V_m^{}F_i^{}F_jV_m)\right)=\underset{ij}{}Tr\left(\varrho E_i^{}E_j\tau ^{}(G_i^{}G_j)\right)0$$
(38)
where we put $`G_iF_i\tau ^{}(V_m)`$ for any $`i`$, and this completes the proof of our statement.
Let us turn to the question of nondecomposable maps. We assume that $`(\widehat{1\mathrm{l}}\tau )\varrho 0`$ and the map $`S`$ defined in (22) is decomposable. So $`S=S_1+S_2\tau `$ where $`S_1,S_2`$ are completely positive maps and $`\tau `$, as usually, is the transposition. We observe
$$0>TrH\varrho =<\underset{i}{}e_ie_i|\widehat{1\mathrm{l}}(S_1^{}+\tau ^{}S_2)(\varrho )|\underset{i}{}e_ie_i>.$$
(39)
But this is a contradiction since the just proved statement implies
$$(S_1^{}+\tau ^{}S_2^{})(\varrho )0$$
(40)
Consequently, $`S`$ can not be a decomposable map.
To close our remarks on positive maps let us mention that there is a nice relationship between such approach to a construction of a class of positive maps and unextedible systems (cf. , ). To do this we have to recall some Woronowicz results. Namely, as the general theory of positive maps is very complicated (cf. Choi’s scheme), the investigation of positive maps based on analysis of extreme positive maps did not provide satisfactory results (cf. ), Woronowicz () proposed to study nonextendible positive maps - a map $`\mathrm{\Phi }`$ is nonextendible if it can not be written in the form
$$\mathrm{\Phi }(a)=P\stackrel{~}{\mathrm{\Phi }}(a)P$$
(41)
where $`\stackrel{~}{\mathrm{\Phi }}`$ is another positive map, $`P`$ is a projection, in a nontrivial way (see for a precise definition). Now assume $`dim_𝒜=n<\mathrm{}`$, $`dim_{}=n<\mathrm{}`$. Then the map $`S`$ defined with $`H=_{i=1}^n|\alpha _i><\alpha _i||\beta _i><\beta _i|`$ for a properly chosen orthonormal systems leads to a nonextendible map. To prove this statement it is enough to apply Theorem 3.3 in (). To appreciate the notion of nonextendible positive maps we recall some Woronowicz results: nonextendible normalized positive map is extreme in the cone of all positive maps, and any normalized positive map admits a nonextendible extension. Moreover any Jordan map (normalized positive map $`\mathrm{\Phi }`$ such that $`\mathrm{\Phi }(a^2)=(\mathrm{\Phi }(a))^2`$) is nonextendible.
To end this paper, we want to stress that nearly all given results follow, more or less, directly from mathematical features of tensor products. But using the Quantum Mechanics rule saying that composite systems are described by tensor products we should expect that some properties of tensor product have a great impact on a description of composite systems. Therefore, all peculiarities related to the order in tensor products should be taken into account. In particular, we should remember that the order in a $`C^{}`$-algebra (so in $`()`$) is not susceptible to the order of multiplication. The same can be said about the notion of positive map. On the contrary the notion of completely positive map and the notion of $`m`$-positive map ($`m>1`$) are susceptible to the order of multiplication. This explains the role of transposition $`\tau `$, its relation to $`M_n^+M_k^+`$, $`(M_nM_k)^+`$ (in Choi’s scheme), and finally the significance of the proper chosen family of functionals in definition of the order in the tensor product. This concluding “mathematical” remark can be also considered as a partial answer to the question related to some experiments of the Bell type, like those studied in (). Namely, there are problems associated with propagation properties (variables) of photons as well as with corellations between their spins. But such a system is an example of a composite system. Moreover we note that any real experiment involves time evolution. So, a description of such a system deals with dynamical maps preserving an order structure of tensor product. Furthermore, an analysis of such the systems involves spin variables as well as other ones necessary for full description. Consequently, in general, the considered maps should be sensitive to the order properties of various variables. Taking the separability problem as an illustration we have argued that basic features of the problem really follow from such the general theory.
Acknowledgments: We are grateful to M. Zukowski and M. Marciniak for remarks. W.A.M. acknowledges the support from K.B.N. (project PB/0273/PO3/99/16 ).
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# Prompt Photon Production and Observation of Deeply Virtual Compton Scattering
## 1 Introduction
Two types of real photon production processes in $`e^+p`$ collisions which are currently studied using the ZEUS detector at HERA are prompt-$`\gamma `$ production in photoproduction ($`\gamma p`$) and Deeply Virtual Compton Scattering (DVCS) in deep inelastic scattering (DIS).
Prompt-$`\gamma `$ production in $`\gamma p`$ interactions is the production of a real $`\gamma `$ directly from the hard interaction of a quasi-real $`\gamma `$ (invariant mass $`Q^20`$ <sup>1</sup><sup>1</sup>1$`Q^2=(ee^{})^2`$, where $`e`$ and $`e^{}`$ are the initial and final positron four-momenta, respectively.) with the proton. At leading order, two kinds of $`\gamma p`$ processes can be defined: direct, where the $`\gamma `$ participates entirely in the hard interaction, and resolved, where the $`\gamma `$ first fluctuates into a hadronic system and a parton from this system then enters into the hard interaction. Examples are depicted in Fig. 1. Prompt-$`\gamma `$ production
is less influenced by hadronisation effects than, e.g., dijet production, although its production rate is down by a factor $`\alpha /\alpha _S`$ in comparison. It is interesting because of its sensitivity to the parton density of the $`\gamma `$. The measurements can also be used to test next-to-leading order (NLO) perturbative QCD (pQCD).
DVCS , depicted in Fig. 2, is the hard diffractive scattering of a $`\gamma `$ off a proton, $`e^+pe^+\gamma p`$. This is a process which has never been seen before at high energies, but which is predicted by FFS to have a fairly high counting rate.
DVCS is an exciting process because of its potential for accessing the skewed parton distributions (SPD’s) of the proton. SPD’s, which quantify two-particle correlations in the proton, are a generalisation of the usual proton parton distributions to the case where the squared momentum transfer to the proton is non-zero. The advantage of DVCS over processes where a hadron is diffractively produced is two-fold: the theoretical uncertainty in the hadronic wave function is avoided, and the DVCS rate is predicted to be less suppressed by a factor of $`Q^2`$ . Furthermore, its final state is identical to that of QED Compton scattering (QEDC), and interference of the two processes potentially allows the measurement of the real part of a QCD amplitude.
A common feature of these two studies at ZEUS is that each involves the detection of a $`\gamma `$ in the barrel part of the calorimeter (BCAL), which has its electromagnetic (EM) section segmented into $`5\times 20`$ cm<sup>2</sup> cells. A potential background contribution to $`\gamma `$ candidates, i.e. those tagged EM clusters not associated with a track in the central tracking detector, arises from $`\pi ^0`$ and $`\eta `$ production. However, these particles tend to produce broader clusters in the BCAL than single $`\gamma `$’s. Two shower shape variables are employed to exploit this difference: $`Z_{width}`$, which is the energy-weighted average of the width of the EM cluster in the $`Z`$-direction <sup>2</sup><sup>2</sup>2 The ZEUS coordinate system is defined as right-handed with the $`Z`$-axis pointing in the forward (proton beam) direction. The origin is at the nominal $`ep`$ interaction point, and the polar angle $`\theta `$ is defined with respect to the positive $`Z`$-direction. (the direction in which the BCAL is most finely segmented); and $`f_{max}`$, which is the fraction of the EM cluster energy carried by the most energetic cell in the cluster.
## 2 Prompt Photon Production
The $`f_{max}`$ distribution for the 1996/97 prompt-$`\gamma `$ photoproduction analysis is shown in Fig. 3.
Overlaid are the expected contributions from $`\eta `$, $`\pi ^0`$, and $`\gamma `$. A $`\gamma `$ signal at high $`f_{max}`$ values is evident. The prompt-$`\gamma `$ production cross section for $`\gamma `$ transverse energies $`E_T^\gamma >5`$ GeV and $`\gamma `$ rapidity interval $`0.7\eta ^\gamma 0.9`$ is shown in Fig. 4 as a function of $`E_T^\gamma `$ and $`\eta ^\gamma `$.
Overlaid are the NLO predictions of two groups of theorists, LG and KZ , each using two different $`\gamma `$ parton density parametrisations: GRV and GS . The predictions are in reasonable agreement with the data, although those based on the GS parametrisation tend to be low. This demonstrates the sensitivity of this type of analysis to the $`\gamma `$ parton density. With more data and further theoretical progress the prompt-$`\gamma `$ process will provide a valuable tool for studying NLO pQCD and measuring the $`\gamma `$ parton density.
## 3 Deeply Virtual Compton Scattering
For the DVCS analysis , events with only two EM clusters and at most one track (which must be matched to one of the clusters) are selected. The first (second) candidate, corresponding to the scattered $`e^+`$ ($`\gamma `$) in the DVCS case, must have polar angle $`\theta _1>2.8`$ ($`\theta _2<2.4`$) radians and $`E_1>10`$ GeV ($`E_2>2`$ GeV). To suppress the QEDC process the polar angle difference must satisfy $`|\theta _1\theta _2|>0.8`$ radians. Among the remaining requirements are a cut of $`Q^2>6`$ GeV<sup>2</sup>, calculated using the first EM candidate, and a cut on the invariant mass of the two EM candidates, $`M_{12}<30`$ GeV. From $`37\mathrm{p}\mathrm{b}^1`$ of $`e^+p`$ data, 1954 events remain for further study after application of all cuts.
As an aid to studying DVCS, a MC generator GenDVCS based on the DVCS, QEDC, and interference term (int) cross sections provided by FFS was developed at ZEUS. Samples of DVCS+QEDC+int events (elastic only) were generated according to the predicted cross sections, run through a full detector and trigger simulation, and processed using the same reconstruction program as the real data. Additional samples were similarly generated using the QEDC generator Compton2.0 (elastic only) for comparison, as well as RAPGAP (diffractive events) and DJANGOH (inclusive DIS events) samples for studying background from $`\pi ^0/\eta `$ contamination.
Shown in Fig. 5 is the polar angle distribution of the
second EM candidate, $`\theta _2`$, for (a) all candidates, and (b) only those candidates without a matched track ($`\gamma `$ candidates). The third plot (c) shows the ratio of (a) and (b). Overlaid are the predictions from Compton2.0 and GenDVCS, normalised to the same luminosity as the data, based on their calculated cross sections. (The QEDC ratio is unchanged when the inelastic component is included.) There is a clear deficit in the number of small angle $`\gamma `$ candidates predicted by the QEDC simulation. The inclusion of DVCS brings the prediction into reasonable agreement with the data.
A potential source of background arises from processes like $`e^+pe^+\pi ^0p`$, etc., where a $`\pi ^0`$/$`\eta `$ fakes a $`\gamma `$ signal in the calorimeter. Redisplayed in Figs. 6(a) and (b) is the $`\theta _2`$ distribution for $`\gamma `$ candidates. Overlaid are the RAPGAP and DJANGOH predictions. It may thus appear possible to explain the data as being due to such events. However, these two generators are not expected to predict accurate rates for low-multiplicity $`\pi ^0/\eta `$ production. In fact, as shown in Figs. 6(c)-(f), the $`Z_{width}`$ and $`f_{max}`$ distributions for BCAL $`\gamma `$ candidates having $`\theta _2<1.6`$ radians, where the QEDC contribution is predicted to be small, indicate that the $`\pi ^0/\eta `$ hypotheses cannot account for the EM shower shapes, and so the data cannot be explained as hadronic background from low-multiplicity DIS events. A repetition of the analysis with a harder energy cut ($`E_2>5`$ GeV) on the second candidate further supports this conclusion. This, then, is first evidence for DVCS at high energy.
## 4 Acknowledgements
I gratefully acknowledge the many useful discussions about DVCS I have had with A. Freund and M. Strikman. I thank the conveners of my session for their kind and considerate handling of the scheduling of my talk.
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# Introduction
## Introduction
Dimensional regularization and minimal subtraction are the most convenient tools for multi-loop computations. However the accuracy of the computations of chiral quantities is much lower than for non-chiral quantities, due to the difficulties encountered in restoring the chiral symmetries in the BMHV scheme , the only one known until recently to deal consistently with $`\gamma ^5`$ in dimensional regularization.
Few chiral quantities have been computed in the BMHV scheme; among the most accurate computations made up to now there are the renormalization at two loops of the non-singlet axial current , of the singlet axial current and of $`F\stackrel{~}{F}`$ in QCD, and the corresponding calculation of the three-loop anomalous dimension of the singlet axial current .
The techniques used in these papers to guarantee the validity of the chiral Ward identities in the BMHV scheme are indirect, and not suitable to be generalized to chiral gauge theories. The most impressive trick used is the determination of two-loop finite counterterms in the fermionic sector of the singlet axial current operator by making a three-loop computation in its gluonic sector .
Recently it has been shown that there is a consistent extension of the BMHV scheme, called semi-naive dimensional regularization (SNDR), in which after minimal subtraction only few graphs can produce terms breaking the chiral Ward identities. The SNDR scheme has been applied in that paper to the case of the renormalization at two loops of the Yukawa model in presence of external gauge fields, previously studied in the BMHV scheme .
As a preliminary investigation in the use of the techniques developed in and for renormalizing gauge theories with chiral interactions, in this letter we apply the SNDR scheme to renormalize at two loops the axial current and $`F\stackrel{~}{F}`$ in QED; to check the chiral Ward identities and to compute the anomalous dimensions we will use the Wilsonian methods introduced in .
We find that using the minimal subtraction prescription in SNDR (MS-SNDR) for the axial current the anomalous axial Ward identity is preserved by choosing a non-minimal subtraction for $`F\stackrel{~}{F}`$. Vice-versa, one can satisfy the anomalous axial Ward identity by choosing the minimal subtraction for the operator $`F\stackrel{~}{F}`$ and a non-minimal subtraction for the axial current. With the latter renormalization prescription we find the same three-loop anomalous dimension of the axial current as in the BMHV scheme , where minimal subtraction is made on $`F\stackrel{~}{F}`$. Notice that in the BMHV it is not possible to choose minimal subtraction on the axial current, due to the presence of the same kind of spurious anomalies appearing in the non-singlet axial current.
In the first section we review the Adler-Bardeen theorem , following to a large extent the regularization-independent derivation in .
In the second section we review the SNDR scheme and we perform the two-loop renormalization of the axial current and of $`F\stackrel{~}{F}`$ in QED.
In the third section we discuss our methods of computation, we compare our results with those in and we discuss in this context the relation between SNDR and BMHV.
## 1 Review of the Adler-Bardeen theorem
The classical action of QED with $`N_f`$ massless fermions in presence of sources for $`J_\mu ^5`$ and $`K_\mu `$ is
$`S^{(0)}={\displaystyle \overline{\psi }\gamma _\mu _\mu \psi }+{\displaystyle \frac{1}{4}}F_{\mu \nu }^2+{\displaystyle \frac{1}{2\alpha }}(_\mu V_\mu )^2+ieV_\mu \overline{\psi }\gamma _\mu \psi +A_\mu J_\mu ^5+\chi _\mu K_\mu `$ (1)
where we define
$`J_\mu ^5i\overline{\psi }\gamma _\mu \gamma ^5\psi ;K_\mu 4iϵ_{\mu \nu \rho \sigma }V_\nu _\rho V_\sigma `$ (2)
We use Euclidean space conventions.
On the functional generator $`\mathrm{\Gamma }=\mathrm{\Gamma }[V,\psi ,\overline{\psi },A,\chi ]`$ of $`1PI`$ vertex functions the renormalization group equation reads
$`𝒟\mathrm{\Gamma }=0`$ (3)
$`𝒟\mu {\displaystyle \frac{}{\mu }}+\beta {\displaystyle \frac{}{e}}+\delta \alpha {\displaystyle \frac{}{\alpha }}{\displaystyle \underset{i}{}}\gamma _iN_i{\displaystyle (\gamma _{\chi A}\chi _\mu \frac{\delta }{\delta A_\mu }+\gamma _{A\chi }A_\mu \frac{\delta }{\delta \chi _\mu })}`$
where the index $`i`$ runs over $`V_\mu ,\psi ,\overline{\psi },A_\mu ,\chi `$; $`\gamma _i`$ and $`N_i`$ are the corresponding anomalous dimensions and number operators; $`\gamma _{A\chi }`$ and $`\gamma _{\chi A}`$ are the mixing anomalous dimensions for the operators (2).
The vectorial Ward identity reads
$`G_v\mathrm{\Gamma }={\displaystyle \frac{1}{\alpha }}^2_\mu V_\mu +4iϵ_{\mu \nu \rho \sigma }_\mu \chi _\nu _\rho V_\sigma `$
$`G_v_\mu {\displaystyle \frac{\delta }{\delta V_\mu }}+ie\psi {\displaystyle \frac{\delta }{\delta \psi }}ie\overline{\psi }{\displaystyle \frac{\delta }{\delta \overline{\psi }}}`$ (4)
Due to the linearity of the breaking terms, they do not need to be renormalized. The proof of this fact in the case of the gauge-fixing term can be straightforwardly extended to the $`\chi `$ term.
The axial Ward identity is
$`G_a\mathrm{\Gamma }=\rho _\mu {\displaystyle \frac{\delta \mathrm{\Gamma }}{\delta \chi _\mu }}`$ (5)
$`G_a_\mu {\displaystyle \frac{\delta }{\delta A_\mu }}+i\gamma ^5\psi {\displaystyle \frac{\delta }{\delta \psi }}+i\overline{\psi }\gamma ^5{\displaystyle \frac{\delta }{\delta \overline{\psi }}}`$
The system of constraints (3,1,5) satisfies the following consistency conditions:
$`[𝒟,G_v]=i\beta (\psi {\displaystyle \frac{\delta }{\delta \psi }}\overline{\psi }{\displaystyle \frac{\delta }{\delta \overline{\psi }}})+\gamma _V_\mu {\displaystyle \frac{\delta }{\delta V_\mu }}`$ (6)
$`[𝒟,G_a\rho _\mu {\displaystyle \frac{\delta }{\delta \chi _\mu }}]=(\gamma _A\rho \gamma _{\chi A})_\mu {\displaystyle \frac{\delta }{\delta A_\mu }}+(\gamma _{A\chi }\rho \gamma _\chi \beta {\displaystyle \frac{\rho }{e}})_\mu {\displaystyle \frac{\delta }{\delta \chi _\mu }}`$
which, together with eq.(1) imply respectively
$`\beta =e\gamma _V;\delta =\gamma _\chi =2\gamma _V;\gamma _{A\chi }=0`$ (7)
$`\gamma _A=\rho \gamma _{\chi A};\gamma _\chi =\beta {\displaystyle \frac{ln\rho }{e}}`$ (8)
from which it follows that, since $`\beta 0`$,
$`\rho =ca`$ (9)
where $`a\frac{e^2}{16\pi ^2}`$ and $`c`$ is a constant in $`e`$ which is fixed by the one-loop anomaly computation to be
$`c=N_f`$ (10)
This proof of the Adler-Bardeen theorem , which is close to the one in , relies on the validity of (1); this relation can be imposed independently of the regularization scheme. In a regularization which respects the vectorial Ward identities, like Pauli-Villars or the BMHV and the SNDR dimensional regularization schemes, this relation implies that the term $`\chi _\mu K_\mu `$ of the bare action does not need to be renormalized.
The first relation in (8) gives $`\gamma _A`$ at $`l`$ loops in terms of $`\gamma _{\chi A}`$ at $`(l1)`$-loops.
Allowing for finite renormalization $`f_\alpha `$ and $`f_\chi `$ respectively for the gauge-fixing and $`\chi _\mu K_\mu `$ terms, the relations (7) are modified in a trivial way
$`\delta =2\gamma _V+\beta {\displaystyle \frac{lnf_\alpha }{e}};\gamma _\chi =2\gamma _V+\beta {\displaystyle \frac{lnf_\chi }{e}}`$ (11)
Using (8) one gets
$`\rho ={\displaystyle \frac{ca}{f_\chi }}`$ (12)
## 2 Axial current in SNDR
Let us review the BMHV and the SNDR dimensional regularization schemes.
In the BMHV scheme one considers the Lorentz covariants $`\delta _{\mu \nu }`$, $`\gamma _\mu `$, $`p_\mu `$, etc. as formal objects, satisfying the usual tensorial rules. $`\delta _{\mu \nu }`$ is the Kronecker delta in $`d=4ϵ`$ dimensions; a formal rule for summed indices is given:
$$\delta _{\mu \nu }\delta _{\nu \rho }=\delta _{\mu \rho };\delta _{\mu \nu }p_\nu =p_\mu ;\delta _{\mu \mu }=d$$
(13)
The gamma ‘matrices’ $`\gamma _\mu `$ satisfy the relation
$$\{\gamma _\mu ,\gamma _\nu \}=2\delta _{\mu \nu }I$$
(14)
where $`I`$ is the identity. The trace is cyclic and satisfies
$$trI=4$$
(15)
In the BMHV scheme additional ‘$`(d4)`$-dimensional’ or ‘evanescent’ tensors $`\widehat{\delta }_{\mu \nu },\widehat{p}_\mu `$ and $`\widehat{\gamma }_\mu `$ are introduced; the Kronecker delta in the $`(d4)`$-dimensional space is $`\widehat{\delta }_{\mu \nu }`$ , satisfying
$`\widehat{\delta }_{\mu \nu }\delta _{\nu \rho }=\widehat{\delta }_{\mu \nu }\widehat{\delta }_{\nu \rho }=\widehat{\delta }_{\mu \rho };\widehat{\delta }_{\mu \mu }=ϵ`$
$`\widehat{p}_\mu \widehat{\delta }_{\mu \nu }p_\nu ;\widehat{\gamma }_\mu \widehat{\delta }_{\mu \nu }\gamma _\nu `$ (16)
The Kronecker delta in four dimensions in $`\overline{\delta }_{\mu \nu }`$, satisfying
$`\overline{\delta }_{\mu \nu }\delta _{\mu \nu }\widehat{\delta }_{\mu \nu };\overline{p}_\mu \overline{\delta }_{\mu \nu }p_\nu ;\overline{\gamma }_\mu \overline{\delta }_{\mu \nu }\gamma _\nu `$ (17)
The Levi-Civita antisymmetric tensor has no evanescent component:
$$\widehat{\delta }_{\mu \nu }ϵ_{\nu \rho \sigma \tau }=0$$
(18)
$`\gamma ^5`$ is defined by
$$\gamma ^5=\frac{1}{4!}ϵ_{\mu \nu \rho \sigma }\gamma _\mu \gamma _\nu \gamma _\rho \gamma _\sigma $$
(19)
which implies
$`\{\gamma ^5,\gamma _\mu \}=2\widehat{\gamma }_\mu \gamma ^5;(\gamma ^5)^2=I`$ (20)
The fact that $`\gamma ^5`$ is not-anticommuting with $`\gamma _\mu `$ leads to violations of the (non-anomalous) axial Ward identities. On the other hand the possibility of having anomalous axial currents must be contemplated, so that it is necessary to define $`\gamma ^5`$ so that it is not anticommuting. However (20) introduces many breaking terms which are not related to anomalies; these terms are sometimes called spurious anomalies .
The SNDR scheme is an extension of this scheme. Add to the BMHV Dirac algebra the objects $`\eta `$ and $`\eta _1`$ satisfying the following defining relations:
$`\{\eta ,\gamma _\mu \}=\{\eta _1,\gamma _\mu \}=2\widehat{\gamma }_\mu \eta _1`$
$`\eta ^2=I;\eta \eta _1=\eta _1\eta =\eta _1^2;\eta _1^3=\eta _1`$
$`tr\eta \gamma _\mu \gamma _\nu \gamma _\rho \gamma _\sigma =tr\eta _1\gamma _\mu \gamma _\nu \gamma _\rho \gamma _\sigma =4ϵ_{\mu \nu \rho \sigma }`$ (21)
$`tr\eta _1=tr\eta _1^2=0;tr\eta _1^2\gamma _{\mu _1}\mathrm{}\gamma _{\mu _r}=0`$
$`tr\eta _1\gamma ^5\gamma _{\mu _1}\mathrm{}\gamma _{\mu _r}=tr\eta \gamma ^5\gamma _{\mu _1}\mathrm{}\gamma _{\mu _r}=tr\gamma _{\mu _1}\mathrm{}\gamma _{\mu _r}`$
The trace is cyclic on this enlarged algebra; $`\eta _1^2`$ is a projector.
The idea of this regularization is that using $`\eta `$ instead of $`\gamma ^5`$ in the tree-level chiral vertices, the number of spurious anomalies is greatly reduced. Let us discuss separately the cases of open and closed fermionic lines belonging to a $`1PI`$ Feynman graph.
i) Open fermionic lines
In an open fermionic line $`\eta `$ can be anticommuted naively modulo terms belonging to the $`\eta _1^2`$ subspace, i.e. monomials containing $`\eta _1`$; algebraic manipulations on a monomial containing $`\eta _1`$ give again terms in the $`\eta _1^2`$ subspace, which cannot be confused with those belonging to the orthogonal subspace. The minimal subtraction on a $`1PI`$ Feynman graph consists in subtracting all the poles in $`ϵ`$ and all the finite terms containing $`\eta _1`$. In it is explained why the subtraction of these finite terms is necessary and why it can be considered to be minimal. The idea is that in the process of removing the regulator $`\eta `$ is homomorphically mapped in $`\gamma ^5`$; since $`tr(\gamma ^5\eta )\eta _1=4`$ , there is no way to extend this trace-preserving homomorphism to $`\eta _1`$. Therefore the $`\eta _1`$ terms must be subtracted before applying this homomorphism. Notice that the necessity of subtracting the terms in the $`\eta _1^2`$ subspace has nothing to do a priori with the requirement of satisfying the chiral Ward identities, but it has to do simply with the inner consistency of SNDR.
In the BMHV scheme $`\gamma ^5`$ can be anticommuted naively modulo terms containing again $`\gamma ^5`$; the latter terms can generate spurious anomalies, which cannot be distinguished from the non-anomalous contributions, so that their subtraction is non-trivial, unlike in SNDR.
ii) Closed fermionic lines
If in a trace there is no $`\eta _1`$ and an even number of $`\eta `$, these $`\eta `$ can be anticommuted naively. If in a fermionic trace there is at least one $`\eta _1`$, all the $`\eta `$’s in that trace can be replaced with $`\eta _1`$; if at this point there is an even number of $`\eta _1`$’s, the trace vanishes; these properties reduce greatly the possibility of occurrence of spurious anomalies.
If in a fermionic trace there is an odd number of $`\eta `$ and $`\eta _1`$, they can be replaced with $`\gamma ^5`$. The true anomalies originate from these traces.
The bare action in SNDR corresponding to the classical action (1) is
$`S={\displaystyle Z_\psi \overline{\psi }\gamma _\mu [_\mu +ieV_\mu +iA_\mu \eta ]\psi }+{\displaystyle \frac{1}{4}}Z_VF_{\mu \nu }^2+{\displaystyle \frac{1}{2\alpha }}(_\mu V_\mu )^2+`$
$`iA_\mu \overline{\psi }\gamma _\mu (Z_{A2}\eta _1+Z_{A3}\gamma ^5)\psi +\chi _\mu (K_\mu +Z_{\chi A}i\overline{\psi }\gamma _\mu \gamma ^5\psi )`$ (22)
where the sources satisfy $`\widehat{A}_\mu =0=\widehat{\chi }_\mu `$. The bare action in MS-SNDR is vector gauge invariant apart from tree-level terms; in this scheme the vectorial Ward identities are automatically satisfied.
We have computed the two-loop renormalization constants such that the axial Ward identity (5) is satisfied. The only axial Ward identity which is not trivially satisfied is the one involving Figure 1, leading to the non-minimal renormalization
$`Z_{A3}^{nonMS}\rho Z_{\chi A}^{nonMS}=3N_fa^2`$ (23)
All the other renormalization constants can be chosen minimal in SNDR. At the first two loops the renormalization constants are, for $`\alpha =1`$,
$`Z_V=Z_V^{MS}=1{\displaystyle \frac{8}{3ϵ}}N_fa{\displaystyle \frac{4}{ϵ}}N_fa^2`$
$`Z_\psi =Z_\psi ^{MS}=1{\displaystyle \frac{2a}{ϵ}}+({\displaystyle \frac{2}{ϵ^2}}+{\displaystyle \frac{3}{2ϵ}}+{\displaystyle \frac{2N_f}{ϵ}})a^2`$
$`Z_{A2}=Z_{A2}^{MS}=4a+({\displaystyle \frac{16}{3ϵ}}+{\displaystyle \frac{4}{9}})N_fa^2+({\displaystyle \frac{8}{ϵ}}+22)a^2`$ (24)
$`Z_{A3}={\displaystyle \frac{12}{ϵ}}N_fa^2+Z_{A3}^{nonMS}`$
$`Z_{\chi A}={\displaystyle \frac{24a}{ϵ}}+({\displaystyle \frac{64}{ϵ^2}}+{\displaystyle \frac{16}{3ϵ}})N_fa^2+({\displaystyle \frac{48}{ϵ^2}}+{\displaystyle \frac{84}{ϵ}})a^2+Z_{\chi A}^{nonMS}`$
where $`Z_{A3}^{nonMS}`$ is of order $`a^2`$, whereas $`Z_{\chi A}^{nonMS}`$ is considered at order $`a`$ only, since its $`a^2`$ term is related by the axial Ward identity to the three-loop $`a^3`$ term of $`Z_{A3}^{nonMS}`$.
The corresponding anomalous dimensions are given by
$`\gamma _V={\displaystyle \frac{4}{3}}N_fa+4N_fa^2`$
$`\gamma _\psi =a(2N_f+{\displaystyle \frac{3}{2}})a^2`$ (25)
$`\gamma _{\chi A}=24a72a^2{\displaystyle \frac{32}{3}}N_fa^2+{\displaystyle \frac{16}{3}}N_faZ_{\chi A}^{nonMS}`$
the remaining anomalous dimensions being fixed by eqs.(7, 8). In particular one determines in this way indirectly the three-loop anomalous dimension of the axial current
$`\gamma _A=24N_fa^2+72N_fa^3+{\displaystyle \frac{32}{3}}N_f^2a^3{\displaystyle \frac{16}{3}}N_f^2a^2Z_{\chi A}^{nonMS}`$ (26)
## 3 Discussion of the results
To obtain the results in the previous section we have used the Wilsonian method devised in for computing in a systematic way the finite counterterms needed to restore chiral Ward identities in dimensional regularization schemes.
In this approach the renormalization of the theory is obtained imposing renormalization conditions on a Wilsonian functional $`\mathrm{\Gamma }^\mathrm{\Lambda }`$, which is perturbatively defined with the same Feynman rules for the vertices as in the usual $`1PI`$ functional generator $`\mathrm{\Gamma }`$, but with the usual propagators $`D`$ replaced by ‘hard’ propagators
$`D^H=(1K^\mathrm{\Lambda })D;K^\mathrm{\Lambda }(p)=({\displaystyle \frac{\mathrm{\Lambda }^2}{p^2+\mathrm{\Lambda }^2}})^2`$ (27)
Let us denote by $`S_W`$ the local functional whose tree-level part is equal to the classical action (1), and which for $`l1`$ is equal to the marginal part of the Wilsonian functional $`\mathrm{\Gamma }^\mathrm{\Lambda }`$; one has
$`S_W={\displaystyle a_\psi \overline{\psi }\gamma _\mu _\mu \psi }+{\displaystyle \frac{1}{2}}a_{V1}(_\mu V_\nu )^2+{\displaystyle \frac{1}{2}}a_{V2}(_\mu V_\mu )^2+{\displaystyle \frac{1}{4}}a_{V4}V_\mu ^2V_\nu ^2+`$
$`ia_{\psi V}V_\mu \overline{\psi }\gamma _\mu \psi +A_\mu (a_{A\chi }K_\mu +a_AJ_\mu ^5)+\chi _\mu (a_\chi K_\mu +a_{\chi A}J_\mu ^5)`$ (28)
The renormalization group equation and the Ward identities assume now the ‘effective’ form
$`𝒟S_W=𝒯_\gamma `$ (29)
$`{\displaystyle ϵ_vG_vS_W}={\displaystyle ϵ_v(\frac{1}{\alpha }^2_\mu V_\mu +4iϵ_{\mu \nu \rho \sigma }_\mu \chi _\nu _\rho V_\sigma )}+𝒯_v`$ (30)
$`{\displaystyle ϵ_aG_aS_W}={\displaystyle ϵ_a\rho _\mu \frac{\delta S_W}{\delta \chi _\mu }}+𝒯_a`$ (31)
where the $`𝒯`$ terms are contributions to the number operators and to the contact term operators at the Wilsonian scale. The Green functions of the $`𝒯`$ terms can be computed using Feynman rules as described in .
We use the hard propagator $`D_{\mu \nu }^H=D^H\delta _{\mu \nu }`$ for the photon in the Feynman gauge.
To compute the renormalization group functions we use (29). The derivative with respect to the gauge-fixing parameter in (29) in $`\alpha =1`$ is treated as the insertion of the operator $`\frac{1}{2}(_\mu V_\mu )^2`$.
As discussed before, the vectorial Ward identities are trivially satisfied in MS-SNDR, so that we have to discuss only the axial Ward identities (31).
One has
$`𝒯_a={\displaystyle _\mu ϵ_a(b_\psi J_\mu ^5+b_\chi K_\mu )}`$ (32)
The axial Ward identity gives the following relations
$`a_Aa_\psi +b_\psi \rho a_{\chi A}=0`$ (33)
$`a_{A\chi }+b_\chi \rho a_\chi =0`$ (34)
The marginal part of $`S_W`$ exists for $`\mathrm{\Lambda }>0`$; we will compute its coefficients at $`\mathrm{\Lambda }=\mu `$, the dimensional regularization scale.
At zero and one loops one has, in the MS-SNDR scheme,
$`a_\psi ^{MS}=1+{\displaystyle \frac{a}{6}};a_A^{MS}=1{\displaystyle \frac{43a}{60}}`$
$`a_{\chi A}^{MS}={\displaystyle \frac{77a}{5}};a_{A\chi }^{MS}={\displaystyle \frac{2N_fa}{3}};a_\chi =1`$ (35)
$`b_\psi ={\displaystyle \frac{53a}{60}};b_\chi ={\displaystyle \frac{N_fa}{3}}`$
The axial Ward identities (33,34) are manifestly satisfied in the SNDR scheme as long as $`\gamma ^5`$ or the Levi-Civita tensor does not appear in the Feynman rules. This is the case at one loop, with the exception of the $`\mathrm{\Gamma }_{\overline{\psi }\psi \chi \mu }`$ vertex, which is related by the axial Ward identity to the two-loop insertions of $`\mathrm{\Gamma }_{\overline{\psi }\psi A\mu }`$, so that we will discuss it later.
To illustrate this renormalization procedure, consider for instance the relevant part of the unrenormalized Wilsonian axial vertex $`\mathrm{\Gamma }_{\overline{\psi }\psi A\mu }^{unren\mathrm{\Lambda }}`$ at one loop:
$`ia[({\displaystyle \frac{2}{ϵ}}{\displaystyle \frac{43}{60}})\overline{\gamma }_\mu \eta +4\overline{\gamma }_\mu \eta _1]`$ (36)
The pole term and the $`\eta _1`$ terms are minimally subtracted, giving the renormalization constants $`Z_\psi `$ and $`Z_{A2}`$ at one loop (2); the remaining finite part, evaluated in four dimensions (i.e. $`\eta \gamma ^5`$), gives the contribution to $`S_W`$ in (3,3).
At two loops the MS-SNDR scheme satisfies manifestly the axial Ward identities, apart from the contribution of Fig.1. The two-loop Wilsonian Green function $`\mathrm{\Gamma }_{\overline{\psi }\psi A\mu }^{unren\mathrm{\Lambda }}`$, unrenormalized at two loops but renormalized at one loop, gives the following relevant contributions
$`ia^2\overline{\gamma }_\mu \{({\displaystyle \frac{2}{ϵ^2}}+{\displaystyle \frac{3}{2ϵ}}{\displaystyle \frac{806251}{204120}}+{\displaystyle \frac{12937v}{4374}})\eta ({\displaystyle \frac{8}{ϵ}}+22)\eta _1`$
$`N_f({\displaystyle \frac{2}{ϵ}}{\displaystyle \frac{1363}{5670}}+{\displaystyle \frac{1136v}{243}})\eta +N_f({\displaystyle \frac{16}{3ϵ}}{\displaystyle \frac{4}{9}})\eta _1N_f({\displaystyle \frac{12}{ϵ}}+{\displaystyle \frac{2449}{270}}+{\displaystyle \frac{1072v}{81}})\gamma ^5\}`$
where the $`\overline{\gamma }_\mu \gamma ^5`$ contribution, which comes from the graph in Figure 1, is the only one for which minimal subtraction is not automatically sufficient to preserve the axial Ward identity. The pole terms and the $`\eta _1`$ terms are minimally subtracted, giving the renormalization constants $`Z_\psi `$ ,$`Z_{A2}`$ and $`Z_{A3}^{MS}`$ at two loop (2).
Let us consider for instance the contributions to the axial Ward identity due to the graphs in Figure 2
and the corresponding counterterms:
$`a_A^{MS}(fig.2)=({\displaystyle \frac{1363}{5670}}{\displaystyle \frac{1136v}{243}})N_fa^2`$
$`b_\psi ^{MS}(fig.2)=({\displaystyle \frac{28393}{17010}}+{\displaystyle \frac{2272v}{729}})N_fa^2`$
$`a_\psi ^{MS}(fig.2)=({\displaystyle \frac{1736}{1215}}+{\displaystyle \frac{5680v}{729}})N_fa^2`$
There is no $`K_\mu `$ insertion corresponding to the graphs in Figure 2, so that $`a_{\chi A}^{MS}(fig.2)=0`$ and the axial Ward identity (33) is satisfied as expected.
Let us consider finally the only axial Ward identity which is not automatically satisfied in this scheme. For the graph in Fig. 1 one has
$`a_A^{MS}(fig.1)=({\displaystyle \frac{2449}{270}}+{\displaystyle \frac{1072v}{81}})N_fa^2`$
$`b_\psi ^{MS}(fig.1)=({\displaystyle \frac{2519}{270}}+{\displaystyle \frac{1072v}{81}})N_fa^2`$
$`a_\psi ^{MS}(fig.1)=0`$
The corresponding $`K_\mu `$ insertion gives $`a_{\chi A}^{MS}(fig.1)=\frac{77a}{5}`$ (see eq.(3)) so that in the minimal scheme the axial Ward identity (33) is not satisfied
$`(a_Aa_\psi +b_\psi \rho a_{\chi A})^{MS}(fig.1)=3N_fa^2`$ (37)
and the non-minimal counterterms in (23) must be added to restore (33).
To obtain the bare action in the BMHV scheme, one can use the values found for $`S_W`$ as renormalization conditions at the Wilsonian scale. The result is simply phrased: it is sufficient to replace $`\eta `$ and $`\eta _1`$ in (2) with $`\gamma ^5`$. In fact, if $`\eta `$ or $`\eta _1`$ belongs to an open fermionic line of a $`1PI`$ graph, anticommuting it through the gamma matrices of the fermionic lines one uses only the relations in the first line of (2), which agree with (20) after replacing $`\eta `$ and $`\eta _1`$ with $`\gamma ^5`$. If $`\eta `$ or $`\eta _1`$ belongs to an open fermionic line of a $`1PI`$ graph, it gives the same as in the case in which it is replaced by $`\gamma ^5`$, due to the last line in (2).
The correspondence between the bare actions in these two schemes is not always so simple; in presence of more than one chiral vertex the BMHV scheme produces many finite counterterms which are absent in the SNDR scheme; the difference is due to the fact that $`tr\eta _1^2=0`$ in SNDR (see (2)) whereas $`tr(\gamma ^5)^2=4`$ in BMHV; see for instance the Yukawa model in the BMHV scheme and in the SNDR scheme .
Our results agree with those in , after the changes due to conventions, provided one makes the minimal choice for $`Z_{\chi A}`$. Similarly for the anomalous dimensions. In the BMHV scheme it is possible to make this minimal choice, whereas one cannot make minimal subtraction of the axial current, since the one-loop finite counterterm $`4aA_\mu J_\mu ^5`$ in the bare action necessary to satisfy the axial Ward identity cannot be replaced by the term $`\frac{4}{N_f}\chi _\mu J_\mu ^5`$ without violating the classical limit of the operator $`K_\mu `$.
On the other hand in the SNDR scheme it is possible to make the minimal subtraction either on $`K_\mu `$ or on $`J_\mu ^5`$, to all orders in perturbation theory.
Let us review the tricks used in to perform this computation. To compute the one-loop finite counterterm for the axial current, comparison between the axial and the vector vertices is made, as suggested in . To compute the two-loop finite counterterm the same trick cannot be used in the case of the singlet axial current, due to fact that the graph in Figure 1 has no counterpart in the vector vertex. To fix this finite counterterm, instead of checking directly the axial Ward identity on the axial vertex at two loops, the three-loop computation of $`<_\mu J_\mu ^5V_\nu V_\rho >`$ has been made in , obtaining these two-loop finite terms by consistency with the Adler-Bardeen theorem.
## 4 Conclusion
The SNDR scheme is a consistent extension of the BMHV dimensional regularization and renormalization scheme, which has been introduced to reduce the number of spurious anomalies present in the latter scheme.
As a preliminary investigation in gauge theories with chiral couplings, in this letter we have applied the SNDR scheme to the renormalization of the axial current in QED. In this case there is only one chiral vertex, so that there are few spurious anomalies, which have been determined in the BMHV scheme in , together with the three-loop anomalous dimension of the singlet axial current in QCD. We find agreement with these results in the QED case. The correspondence between the SNDR and the BMHV scheme is in this case so easy that it can be made even at the bare action level. As expected, we found that it is easier to satisfy the axial Ward identity in the SNDR scheme than in the BMHV scheme. These computational advantages are expected to be much greater in the case of chiral gauge theories, which have been renormalized systematically in the BMHV scheme only at one loop .
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# Non-orthogonal preferred projectors for modal interpretations of quantum mechanics
## I Introduction
In operational quantum mechanics, theoretical predictions take the form ‘if such-and-such a measurement is made after such-and-such a preparation, such-and-such an outcome will be found with such-and-such a probability’. In contrast, a realist interpretation is an attempt to understand quantum mechanics as making stronger claims of the form ‘such-and-such a variable has such-and-such a value with such-and-such a probability’. The ‘elements of reality’ of Einstein, Podolsky and Rosen, and Bell’s ‘be-ables’ are two ways of referring to the variables that possess definite values in a realist interpretation. We will simply refer to them as the determinate variables. Since the set of determinate variables in some sense specifies ‘what exists’, we call it the ontology for the system. The specification of the values of the determinate variables will be called the value ascription to the ontology. Within this approach, one assigns a property to a system by assigning a value to a determinate variable. Since the ontology and the value ascription together constitute a complete specification of the properties of a quantum system, they will jointly be referred to as the property ascription.
Within ‘orthodox’ interpretations of this type, a variable is determinate if and only if it is associated with an operator for which the state vector is an eigenstate, and its value is the corresponding eigenvalue. It is also assumed that a variable $`V`$ defined on a subsystem is determinate only if $`VI`$ is determinate on the total system, where $`I`$ is the identity operator for the part of the total system that is not included in the subsystem. According to these rules, the ontologies and the value ascriptions for all systems are uniquely defined by the state vector. It is widely recognized that this view, together with the assumption that the evolution of the state vector is unitary for all time, leads to the quantum measurement problem, namely, the failure to ensure the determinateness of macroscopic variables such as the pointer reading of an apparatus.
One approach to the problem is to introduce a non-unitary dynamics for the state vector into the formalism of the theory (the ‘collapse’ of the state vector). A different approach is to preserve the unitary dynamics, but to reject the notion that a variable is determinate only if it has the state vector as an eigenstate of the associated operator. In the latter type of approach the property ascription need not be fixed at a given time by the state vector. Rather, it may be that the state vector describes only the set of possible property ascriptions, in which case it describes what is possible and what, if anything, is necessary. Since the logic of possibility and necessity is modal logic, realist no-collapse interpretations of this type have been called modal interpretations of quantum mechanics.
Modal interpretations typically impose many constraints on the form of the property ascription for a system. Given these constraints, there is always a unique ‘most elementary’ possessed property defined by a property ascription. We call this the preferred property for that property ascription. Since at a given time the property ascription may be one of several possibilities, each of which define a different preferred property, there is in general a set of preferred properties associated with a system. An example may serve to clarify these concepts. Suppose the system is a digital display on an apparatus. The property ascription for the display may include such properties as ‘the digital display shows a number between $`1`$ and $`3`$’, ‘the digital display shows a number smaller than $`5`$’, etc., while the preferred property may be ‘the digital display shows the number $`2`$’. The set of preferred properties may consist of a list of properties each of the form ‘the digital display shows the number $`k`$’, but differing in the value of $`k.`$ In previous modal interpretations, the preferred properties have been associated with orthogonal projectors.
There are three contributions made in this paper. First, we demonstrate that any modal interpretation which adopts the standard constraints upon the property ascription and which seeks to satisfy a particular criterion of faithful measurement must allow for the set of preferred properties to be associated with non-orthogonal projectors. Second, we introduce a framework for modal interpretations that incorporates such preferred properties. Third, we present a novel proposal within this framework wherein the preferred properties are fixed by a principle of entropy minimization.
The paper is organized as follows. In section 2, we present a review of the constraints upon the property ascription that are standard among modal interpretations, and we provide a rigorous definition of the notion of a preferred property. In section 3, we present the argument that the preferred properties must be associated with non-orthogonal projectors if one hopes to explain the outcomes of perfectly predictable measurements in terms of pre-existing properties of the system under investigation, that is, if perfectly predictable measurements are to be faithful. The argument relies on a particular kind of experiment, involving a sequence of two measurements which have the following critical features: (1) the first measurement disturbs the state of the system differently for different outcomes, resulting in the preparation of non-orthogonal states; and (2) the variable measured by the second device depends on the outcome of the first measurement in such a way that the outcome of the second measurement is always perfectly predictable.
In order to accomodate non-orthogonal preferred properties, we require a new framework for modal interpretations, which is the subject of section 4. We preserve most of the standard constraints on the property ascription, in particular constraints involving the functional relations between the values of variables. However, we show that one must abandon the assumption that different property ascriptions share a common ontology. We assume Healey’s so-called ‘weakening condition’, and adopt Clifton’s rule for relating the properties of composites to the properties of the subsystems of which they are composed. Moreover, we follow previous authors in requiring that the dynamics of the property ascription be Markovian and satisfy certain constraints of analyticity, while also reproducing the standard quantum statistics. Guided by these constraints, we introduce a framework for modal interpretations that incorporates non-orthogonal preferred projectors. This constitutes a generalization of the framework introduced by Bub and Clifton.
With this framework in hand, we proceed in section 5 to present a novel proposal for a modal interpretation. We begin by assuming that there is a distinguished division of the universe into elementary subsystems, or equivalently a distinguished factorization of the total Hilbert space. A preferred decomposition of the state vector is singled out by the minimization of a particular entropic quantity that quantifies the degree of entanglement of the state vector with respect to the distinguished factorization. The actual property ascription is assumed to be fixed by a single element of this decomposition, whose identity evolves by a stochastic process with specified statistical properties. Within this proposal, we demonstrate that the quantum measurement problem is avoided for several models of measurement interactions, and a large class of perfectly predictable measurements are shown to reveal pre-existing properties. In section 6, we present our concluding remarks.
## II The Modal Approach
### A Review of Constraints
We begin by considering the notion of a property of a physical system. The type of properties in which we are interested are those of the form ‘having a value of the variable $`V`$ in the range $`\mathrm{\Delta }`$’. For every such property, one can associate an idempotent variable that has value $`1`$ if the value of $`V`$ is in the range $`\mathrm{\Delta }`$ and $`0`$ if the value lies outside this range. Whether a property is possessed or not is given by the value of this idempotent variable; it is possessed if the value is $`1,`$ and it is not possessed if the value is $`0.`$
In classical mechanics, a variable is represented by a function on phase space, and the possible values of this variable are just the values in the range of this function. The property of ‘having a value of the variable $`V`$ in the range $`\mathrm{\Delta }`$’ is associated with the subset of phase space containing all points for which $`V`$ is in the range $`\mathrm{\Delta }`$. For instance, if the system is a one-dimensional harmonic oscillator, the property ‘having energy between $`E_1`$ and $`E_2`$’ is associated with an elliptical ring in phase space, while the property ‘having position between $`x_1`$ and $`x_2`$’ is associated with a vertical band. Suppose the property $`s`$ is associated with a subset $`\mathrm{\Omega }`$ of phase space. The idempotent variable associated with the property $`s`$ is the function $`F_\mathrm{\Omega }`$ that takes the value 1 for every point in $`\mathrm{\Omega }`$ and $`0`$ for every point outside $`\mathrm{\Omega }`$. If two properties $`s`$ and $`s^{}`$ are represented by subsets $`\mathrm{\Omega }`$ and $`\mathrm{\Omega }^{}`$ of phase space, then the disjunction of $`s`$ and $`s^{}`$ is represented by the union of $`\mathrm{\Omega }`$ and $`\mathrm{\Omega }^{},`$ the conjunction of $`s`$ and $`s^{}`$ is represented by the intersection of $`\mathrm{\Omega }`$ and $`\mathrm{\Omega }^{},`$ and the negation of $`s`$ is represented by the complement of $`\mathrm{\Omega }.`$ The truth tables appropriate for conjunction, disjunction and negation in classical logic place the following constraints on the values of the idempotent variables:
$`F_{\mathrm{\Omega }\mathrm{\Omega }^{}}`$ $`=`$ $`F_\mathrm{\Omega }+F_\mathrm{\Omega }^{}F_\mathrm{\Omega }F_\mathrm{\Omega }^{}`$
$`F_{\mathrm{\Omega }\mathrm{\Omega }^{}}`$ $`=`$ $`F_\mathrm{\Omega }F_\mathrm{\Omega }^{}`$
$`F_{\overline{\mathrm{\Omega }}}`$ $`=`$ $`1F_\mathrm{\Omega },`$
where $`,`$ and an over-bar respectively denote union, intersection and complementation of subsets of phase space.
We now consider an approach to realist interpretations of quantum mechanics that parallels those features of classical theories described above, except of course that the mathematical structure relevant for the description of a system is no longer phase space but Hilbert space. This approach has its origin in the field of quantum logic, and is adopted by most modal interpretations. A variable is represented by a Hermitian operator over the Hilbert space, and its possible values are the eigenvalues of this operator. Properties are associated with subspaces of Hilbert space. The idempotent variable associated with a given property is represented by the projector onto the corresponding subspace. If two properties $`s`$ and $`s^{}`$ are represented by subspaces $`𝒮`$ and $`𝒮^{}`$, then the disjunction of $`s`$ and $`s^{}`$ is represented by the linear span (direct sum) of $`𝒮`$ and $`𝒮^{}`$, the conjunction of $`s`$ and $`s^{}`$ is represented by the intersection of $`𝒮`$ and $`𝒮^{}`$, and the negation of $`s`$ is represented by the orthogonal complement of $`𝒮`$. These assumptions will be called the constraints on logical connectives.
For simplicity, we denote both the projector onto the subspace $`𝒮`$ and the associated idempotent variable by $`P_𝒮`$, and denote the value of this idempotent variable by $`[P_𝒮].`$ Analogously to the classical case, we adopt the following constraint on the values of the idempotent variables:
Functional relation constraint
$`[P_{𝒮𝒮^{}}]`$ $`=`$ $`[P_𝒮]+[P_𝒮^{}][P_𝒮][P_𝒮^{}]`$ (1)
$`[P_{𝒮𝒮^{}}]`$ $`=`$ $`[P_𝒮][P_𝒮^{}]`$ (2)
$`[P_𝒮^{}]`$ $`=`$ $`1[P_𝒮],`$ (3)
where $`,`$ $``$ and denote respectively linear span, intersection, and orthogonal complement.
As it turns out, it is impossible to associate values with all the projectors in a Hilbert space in a way that is consistent with the functional relation constraint and the constraints on logical connectives . The response of modal interpretations is to associate definite values with only a subset of all the projectors in the Hilbert space. Thus, in contrast with classical mechanics, only a subset of all idempotent variables correspond to well-defined properties at any given time. The projectors that are associated with definite values are labelled determinate, as are the corresponding idempotent variables. The functional relation constraint is only required to hold among the determinate projectors.
As regards non-idempotent variables, we adopt the convention that a variable is denoted by the same symbol as the associated hermitian operator, for instance $`V,`$ and that its value is denoted by $`[V].`$ Moreover, following other modal interpreters, we adopt the attitude that if the spectral resolution of a Hermitian operator is $`V=_kv_kP_{𝒮_k},`$ where $`v_k0`$ for all $`k,`$ and $`v_kv_k^{}`$ for $`kk^{},`$ then $`V`$ is determinate if and only if all of the projectors in the set $`\{P_{𝒮_k}\}_k`$ are determinate, and in this case its value is $`[V]=_kv_k[P_{𝒮_k}]`$. We call this the spectral constraint.
It follows from this constraint that the set of determinate idempotent variables and their values are sufficient to specify the set of all determinate variables and their values. As noted in the introduction, we will refer to the set of determinate variables as the ontology, and the values of these variables as the value ascription. The ontology and the value ascription together define the property ascription.
We turn now to constraints on the nature of the ontology. It is typically assumed that logical combinations of well-defined properties are also well-defined. Thus, if property $`s`$ is well-defined, then so too should be the property ‘not $`s`$. In other words, if $`P_𝒮`$ is determinate, then $`P_𝒮^{}`$ ($`=IP_𝒮)`$ should be determinate as well. Similarly, if properties $`s`$ and $`s^{}`$ are well-defined, so that the associated projectors $`P_𝒮`$ and $`P_𝒮^{}`$ are determinate, then the properties ‘$`s`$ or $`s^{}`$’ and ‘$`s`$ and $`s^{}`$’ should also be well-defined, and the associated projectors $`P_{𝒮𝒮^{}}`$ and $`P_{𝒮𝒮^{}}`$ should be determinate. In summary, we require
Closure constraint:
$`\text{If }P_𝒮Ont\text{ then }P_𝒮^{}Ont`$
$`\text{If }\{P_𝒮,P_𝒮^{}\}`$ $``$ $`Ont\text{ then }\{P_{𝒮𝒮^{}},P_{𝒮𝒮^{}}\}Ont,`$
where $`Ont`$ denotes the ontology. Assuming that for every system there is at least one projector that is determinate, it follows from the closure constraint that the identity operator and the projector onto the null space are determinate for every system.
We refer to all of the constraints on the ontology and value ascription that have been presented thus far as the algebraic constraints.
### B Preferred Projectors
We now derive a few consequences of the algebraic constraints. For this purpose, it is useful to introduce a notational convenience: ‘$`P_𝒮>P_𝒮^{}`$’ denotes that $`𝒮^{}`$ is a subspace of $`𝒮,`$ and ‘$`P_𝒮P_𝒮^{}`$’ denotes that $`𝒮`$ and $`𝒮^{}`$ are orthogonal subspaces. First, we show that in every property ascription there is a projector $`P_𝒮`$ that receives the value $`1,`$ and for which no projector onto a proper subspace of $`𝒮`$ receives the value 1; that is, there is a subspace $`𝒮`$ such that
$$[P_𝒮]=1\text{ and there is no subspace }𝒯\text{ such that }P_𝒯<P_𝒮\text{ and }[P_𝒯]=1.$$
(4)
Such a subspace always exists since at least one projector, namely the identity operator, always receives the value 1. Such a subspace is unique because if $`𝒮`$ and $`𝒮^{}`$ were two distinct subspaces satisfying this definition, then both would receive the value $`1,`$ and by the functional relation constraint their intersection would also receive the value $`1,`$ which implies that either $`𝒮`$ or $`𝒮^{}`$ has a proper subspace that receives the value $`1.`$ We call the unique projector satisfying Eq. (4) the preferred projector for the property ascription. The property associated with this projector is also called preferred.
A modal interpretation must account for the fact that a measurement device may have one of several different properties at the end of a measurement despite there being a single state vector for the universe. This is accomplished by assuming that the property ascription is not fixed by the state vector, but may be one of several possibilities at a given time. There are two notions of possibility that are adopted by modal interpreters in this context. In the first, the different possibilities for the final property ascription to the measurement device are attributed to differences in the initial properties(possibly hidden) of the system. In the second, the different possibilities for the final property ascription to the measurement device arise from an objective stochasticity in the evolution of the property state<sup>*</sup><sup>*</sup>*An example of the first approach is the Bub-Clifton interpretation when the preferred variable has a continuous unbounded spectrum, and the dynamics is given by a guidance equation analogous to the one used in Bohmian mechanics. An example of the second approach is the Bub-Clifton interpretation when one adopts a different guidance equation, or when the preferred variable has a discrete spectrum. See, for instance, section 5.2 of Ref. .. Although we adopt the latter view in subsequent sections, for the present it suffices to note that in all modal interpretations one associates with a system, at every time, one out of a set of several possible property ascriptions. Since each such property ascription defines a preferred projector, there is in general a set of preferred projectors associated with a system at a given time. We now consider the relation between the elements of this set.
We begin with a definition: two properties are said to be mutually exclusive if their conjunction is a contradiction. This is stronger than simply being distinct, as is illustrated by the properties ‘red’ and ‘red or blue’, which are distinct, but not mutually exclusive. We also use the term ‘mutually exclusive’ to describe two property ascriptions if there is a property obtaining in one that is mutually exclusive to a property obtaining in the other. We assume that for all systems at all times, the different possible property ascriptions are mutually exclusive. This assumption can be recast as a constraint upon the set of preferred projectors. Recalling the constraints on logical connectives, for two property ascriptions to be mutually exclusive there must be two projectors $`P_{}`$ and $`P_{^{}}`$ such that $`[P_{}]=1`$ in the first property ascription and $`[P_{^{}}]=1`$ in the second, but such that the intersection of $``$ and $`^{}`$ is the null space. By definition, the preferred projector for the first property ascription, call it $`P_𝒮,`$ must be such that $`P_𝒮P_{},`$ and the preferred projector for the second, call it $`P_𝒮^{},`$ must be such that $`P_𝒮^{}P_{^{}},`$ from which it follows that the intersection of $`𝒮`$ and $`𝒮^{}`$ must also be the null space. Thus we conclude that the preferred projectors for mutually exclusive property ascriptions are associated with subspaces whose intersection is the null space. Note that we have not concluded that the preferred projectors are orthogonal. Indeed, the possibility of preferred projectors that are non-orthogonal will be the focus of much of this paper.
## III Non-orthogonal Preferred Projectors
### A The Faithfulness Criterion
One can debate the merits of assuming orthogonal preferred projectors in the context of a macroscopic system such as the pointer on a measurement device. On the one hand, the distinguishability of the different physical states of a pointer suggest that they must be associated with orthogonal projectors; on the other hand, the requirement of quantum-classical correspondence suggests that the alternative positions of a pointer should be associated with projectors onto a set of coherent states, or some similar over-complete and non-orthogonal basis. We do not take a stand on this issue here. We do provide an argument for adopting a non-orthogonal set of preferred projectors, but it appeals to the properties that should be assigned to microscopic rather than macroscopic systems. In particular, we consider a microscopic system that is the object of investigation in a quantum measurement.
Since operational quantum mechanics only makes reference to the properties of macroscopic preparation and measurement devices, the requirement of agreement with the operational theory does not by itself constrain the properties that are assigned to the microscopic systems under investigation. But if the properties of such microscopic systems are to play any explanatory role in the theory, then one would expect their role to be in determining the outcomes of measurements upon them. In particular, we consider the following criterion for assigning determinate status to a variable:
If it can be predicted with probability $`1`$ that a measurement of the variable $`V`$ will yield the result $`v`$, then immediately prior to the measurement the variable $`V`$ is determinate with value $`v.`$
The motivation for adopting this criterion is presented in the next subsection. It is nonetheless worth emphasizing at this point its importance in this paper: the particular form of the framework for modal interpretations that is presented in section 4 and the particular proposal presented in section 5 are both to a large extent attempts to satisfy the faithfulness criterion.
This criterion is applicable to experiments involving a sequence of measurements, the first of which may be considered a preparation. Within the context of operational quantum mechanics, we consider a sequence of two measurements, associated with distinct Hermitian operators, $`V`$ and $`V^{},`$ both belonging to the Hilbert space $`^S.`$ For simplicity, we assume that these operators have non-degenerate eigenvalues, denoted by $`\{v_k\}_{k=1}^m`$ and $`\{v_j^{}\}_{j=1}^m`$ respectively, and eigenvectors denoted by $`\{|\phi _k\}_{k=1}^m`$ and $`\{|\phi _j^{}\}_{j=1}^m`$ respectively. In order to predict the outcome of the second measurement, it is necessary to also specify how, if at all, the first measurement disturbs the state. Suppose then that upon obtaining outcome $`k`$ for the first measurement, the state $`|\stackrel{~}{\phi }_k`$ is prepared, and that the set of vectors $`\{|\stackrel{~}{\phi }_k\}_{k=1}^m,`$ although normalized and non-collinear, are non-orthogonal. Since it follows that $`|\stackrel{~}{\phi }_k|\phi _k`$ for one or more values of $`k,`$ we say that the measurement is disturbing. For simplicity, the second measurement is assumed to be non-disturbing, and the two measurements are assumed to be immediately consecutive. Finally, we take the preparation procedure that precedes the first measurement to be associated with a state vector $`_kc_k|\phi _k,`$ where $`_k\left|c_k\right|^2=1,`$ and $`c_k0`$ for all $`k.`$
Operational quantum mechanics predicts, via the generalized Born rule, that the probability of the second apparatus indicating outcome $`j`$ given that the first apparatus indicates outcome $`k`$ is
$`\text{Prob}(j|k)=\left|\stackrel{~}{\phi }_k|\phi _j^{}\right|^2.`$
It follows that in order for the outcome of the second measurement to be predictable with probability 1 given the outcome of the first measurement, it must be the case that $`\stackrel{~}{\phi }_k|\phi _j^{}=1`$ for some values of $`k`$ and $`j.`$ Hence, the variable $`V^{}`$ measured by the second apparatus must have at least one of the states in the set $`\{|\stackrel{~}{\phi }_k\}_{k=1}^m`$ as an eigenstate. We define a set of variables $`\{V_{(k)}\}_{k=1}^m,`$ such that the variable $`V_{(k)}`$ has $`|\stackrel{~}{\phi }_k`$ as an eigenstate. In particular, we denote the eigenvalues of $`V_{(k)}`$ by $`\{v_{(k),j}\}_{j=1}^m,`$ and the associated eigenvectors by $`\{|\phi _{(k),j}\}_{j=1}^m,`$ and take $`|\phi _{(k),1}|\stackrel{~}{\phi }_k.`$ It then follows that if the first measurement has outcome $`k,`$ so that $`|\stackrel{~}{\phi }_k`$ is prepared, and if $`V^{}=V_{(k)},`$ that is, the second apparatus measures $`V_{(k)},`$ then with probability $`1,`$ the second measurement has the outcome $`1.`$ Hence the faithfulness criterion is applicable in this case, and implies that $`V_{(k)}`$ is determinate with value $`v_{(k),1}`$ at time $`t,`$ immediately prior to the second measurement.
We now introduce a critical assumption about the sequence of measurements: the nature of the second measurement is taken to depend on the outcome of the first. In particular, we imagine a set-up where if the outcome of the first measurement is $`k,`$ then the second apparatus measures the variable $`V_{(k)}`$; we imagine that this is done mechanically by the measurement apparatus, without the intervention of a physicist. In this case, the faithfulness criterion is applicable for all possible outcomes of the first measurement.
We now show that for the faithfulness criterion to be satisfied for such a sequence of measurements, the preferred projectors must be non-orthogonal. Since there is a non-zero probability for the first measurement to have the outcome $`k`$ for every $`k,`$ it follows from the faithfulness criterion that there is a non-zero probability for the system to possess the property $`[V_{(k)}]=v_{(k),1}`$ immediately prior to the second measurement, for every $`k.`$ If a property has non-zero probability of being possessed, it is a possible property. Since the projector $`P_{\stackrel{~}{\phi }_k}`$ associated with the property $`[V_{(k)}]=v_{(k),1}`$ is one-dimensional, it has no non-null proper subspaces, and thus by definition it is the preferred projector for the $`k`$th property ascription. Thus, the set of preferred projectors is $`\{P_{\stackrel{~}{\phi }_k}\}_{k=1}^m.`$ Finally, since the set of vectors $`\{|\stackrel{~}{\phi }_k\}_{k=1}^m`$ is by hypothesis non-orthogonal, the set of preferred projectors must also be non-orthogonal.
For clarity, we briefly repeat this argument in the context of a simple example, illustrated in Fig. 1. Suppose the variables being measured correspond to the components along different spatial axes of the spin operator, $`𝐒,`$ for a spin 1/2 particle. Denote the component along axis $`\widehat{𝐧}`$ by $`𝐒\widehat{𝐧},`$ and the eigenstate associated with eigenvalue $`\pm \mathrm{}/2`$ by $`|\pm \widehat{𝐧}.`$ Suppose the first measurement is along $`\widehat{𝐳},`$ so that $`V=𝐒\widehat{𝐳},`$ $`|\phi _1=|+\widehat{𝐳},`$ and $`|\phi _2=|\widehat{𝐳}.`$ Suppose moreover that the state $`|+\widehat{𝐱}`$ is prepared if the outcome of the first measurement is $`\widehat{𝐳},`$ while no disturbance occurs if the outcome is $`+\widehat{𝐳},`$ so that $`|\stackrel{~}{\phi }_1=|+\widehat{𝐳},`$ and $`|\stackrel{~}{\phi }_2=|+\widehat{𝐱}.`$ Now suppose that the manner in which the nature of the second measurement depends on the outcome of the first is the following: if the first measurement has outcome $`+\widehat{𝐳},`$ then the second measurement is of $`𝐒\widehat{𝐳}`$, while if it has outcome $`\widehat{𝐳}`$, then the second measurement is of $`𝐒\widehat{𝐱}.`$ Thus, $`V_{(1)}=𝐒\widehat{𝐳}`$ and $`V_{(2)}=𝐒\widehat{𝐱}.`$ Assume the initial state of the spin is $`c_1|+\widehat{𝐳}+c_2|\widehat{𝐳},`$ where $`\left|c_1\right|^2+\left|c_2\right|^2=1`$ and $`c_1,c_20.`$
Using the generalized Born rule, it is straightforward to verify that the result of the second measurement is predictable with probability 1 given the result of the first measurement. It follows from the faithfulness criterion that if the outcome of the first measurement is $`+\widehat{𝐳}`$, then $`𝐒\widehat{𝐳}`$ is subsequently determinate with value $`+\mathrm{}/2`$, and if the outcome of the first measurement is $`\widehat{𝐳}`$, then $`𝐒\widehat{𝐱}`$ is subsequently determinate with value $`+\mathrm{}/2`$. Since both of these options occur with non-zero probability, the properties $`[𝐒\widehat{𝐳}]=+\mathrm{}/2`$ and $`[𝐒\widehat{𝐱}]=+\mathrm{}/2`$ are both possible. It follows that the preferred projectors are $`P_{+\widehat{𝐳}}`$ and $`P_{+\widehat{𝐱}},`$ which are non-orthogonal.
### B Motivation
We now consider the reasons for adopting the faithfulness criterion. As is argued by Redhead, in seeking a realist interpretation of quantum mechanics one is seeking an explanation of the successes of the operational version of the theory. One way to secure an explanation of a measurement outcome is to demand that the properties of the systems involved ensure this outcome. The faithfulness criterion goes beyond this however, in that it specifies the form that such an explanation must take. Specifically, it is assumed that the reason a variable $`V`$ is found to have value $`v`$ in a measurement that is predictable with probability 1 is because immediately prior to the measurement $`V`$ is determinate and has value $`v`$. Although this is perhaps the simplest form the explanation could take, it is not the only form, as is evidenced by Bohm’s theory and Bell’s be-able interpretation, where only the outcomes of measurements of certain variables (position in Bohm’s case and lattice fermion number in Bell’s case), are taken to reveal pre-existing values of these variables.
So we see that the faithfulness criterion is not a necessary feature of a realist interpretation. Nonetheless, there have been many attempts to ensure that the outcomes of perfectly predictable measurements do reveal pre-existing values of these variables.
This tradition dates back to von Neumann, who assumed that the determinate variables of a system and their values are fixed by the density operator for the system, $`\rho (t),`$ by what we shall call the orthodox rule, namely,
$`Ont(t)`$ $`=`$ $`\{V|V\rho (t)\rho (t)\}`$ (5)
$`[V]_t`$ $`=`$ $`Tr(V\rho (t)),`$ (6)
where $`Ont(t)`$ indicates the ontology at time $`t,`$ and $`[V]_t`$ indicates the value of $`V`$ at time $`t`$ (this is simply the rule adopted by the ‘orthodox’ realist interpretations discussed in the introduction). Given that after a non-disturbing measurement of the variable $`V`$ with outcome $`v,`$ one can predict, with probability $`1,`$ that the outcome of an immediately consecutive ideal measurement of $`V`$ will also be $`v`$, the faithfulness criterion demands that $`V`$ be determinate with value $`v`$ prior to the second measurement. However, given the orthodox rule, this can only occur if the density operator after the first measurement is a projector onto an eigenstate of $`V`$ associated with eigenvalue $`v.`$ This must be the case regardless of the density operator prior to the first measurement. Thus, in order to satisfy the faithfulness criterion, von Neumann assumed that upon measurement the state vector undergoes a non-unitary evolution (the so-called ‘collapse’) to the eigenvector associated with the outcome of the measurement. As a realist interpretation of quantum mechanics this proposal is at best incomplete since it fails to specify, in terms of the primitives of the theory, the conditions under which a collapse occurs.
Many modal interpretretations also attempt to satisfy the faithfulness criterion, but unlike von Neumann, they abandon the orthodox rule rather than assuming collapse. For instance, it has occurred to many authors, including Kochen, Healey and Dieks, that by assigning determinate status to the projectors in the spectral resolution of the density operator one can satisfy the faithfulness criterion for ideal measurements. Modern versions of this approach include the proposals of Vermaas and Dieks, and Bacciagaluppi and Dickson. However, it was noted by Bacciagaluppi and Hemmo that the Vermaas and Dieks proposal failed to satisfy the faithfulness criterion for certain non-ideal measurements, specifically, disturbing measurements. The same argument can be applied against the Bacciagaluppi and Dickson proposal.
These results do not rule out the possibility that some new proposal involving a different, but still orthogonal, choice of preferred projectors might satisfy the faithfulness criterion for non-ideal measurements. However, by considering an experiment wherein the nature of the second measurement depends on the outcome of the first, we have shown that the faithfulness criterion fails to be satisfied for any modal interpretation that adopts orthogonal preferred projectors.
Thus any modal proposal seeking to satisfy the faithfulness criterion must allow for non-orthogonal preferred projectors. However, a satisfactory proposal must provide an unambiguous rule for identifying the set of preferred projectors for every system at every time, and it remains to be seen whether there exists any rule that consistently satisfies faithfulness. This rule must also satisfy other constraints, such as predicting properties for macroscopic systems that are in accord with our everyday perceptions of them. In particular, it must yield a solution to the measurement problem. It may be that the preferred set cannot be chosen to satisfy faithfulness for all measurements while also satisfying these other constraints. If this were true, it would certainly remove some of the motivation for pursuing a modal interpretation in the tradition of the authors specified above. We are not able to rule out this possibility here. Nonetheless, the range of measurements for which faithfulness is satisfied can at least be expanded if one assumes non-orthogonal preferred projectors, as we demonstrate in section 5 by a specific proposal,
### C Consequences for the ontology
In a modal interpretation, the property ascription to a system at a given time can be one of several possibilities. A question which we now address is whether or not these possibilities should differ with respect to the ontology they ascribe. Most previous modal interpreters have assumed that they should not. In such interpretations, the possible property ascriptions differ only with respect to the value ascription to a single common ontology. However, as we now prove, such an approach is unable to accomodate non-orthogonal preferred projectors.
It is not possible for there to be, at a given time, several possible mutually exclusive property ascriptions which (1) satisfy the algebraic constraints, (2) do not differ with respect to ontology, and (3) are associated with preferred projectors that are non-orthogonal.
Proof. The proof is by contradiction. Suppose the ontology and preferred projector for the $`k`$th property ascription (in the set of possible property ascriptions at a given time) are denoted respectively by $`Ont_k`$ and $`P_k.`$ Since by hypothesis the possible property ascriptions do not differ with respect to ontology, there exists a single set of determinate variables, denoted by $`Ont,`$ such that $`k:Ont_k=Ont`$. Since $`k:P_kOnt_k`$, it follows that $`\{P_k\}_{k=1}^mOnt.`$ In other words, if there is only a single possible ontology at a given time, the preferred projectors for all the different property ascriptions must simultaneously be part of this ontology. Since one of the possible property ascriptions must actually obtain, one of these projectors must receive the value $`1.`$ Moreover, since by hypothesis the preferred projectors are non-orthogonal, it follows that the common ontology includes several non-orthogonal projectors, one of which receives the value 1. However, this is in contradiction with the algebraic constraints, as we now demonstrate.
Suppose that $`P_𝒮`$ and $`P_𝒮^{}`$ are two non-orthogonal preferred projectors, and that $`[P_𝒮]=1.`$ Since the two property ascriptions associated with these are by assumption mutually exclusive, the intersection of $`𝒮`$ and $`𝒮^{}`$ is the null space. Moreover, by the definition of a preferred projector, no non-null proper subspace of $`𝒮`$ or $`𝒮^{}`$ can be determinate. By closure and the fact that $`P_𝒮^{}`$ is determinate, $`P_{(𝒮^{})^{}}`$ is also determinate. By closure and the fact that $`P_𝒮`$ and $`P_{(𝒮^{})^{}}`$ are determinate, the projector onto the intersection of $`(𝒮^{})^{}`$ and $`𝒮`$ must also be determinate. Since no non-null proper subspaces of $`𝒮`$ can be determinate, the intersection of $`(𝒮^{})^{}`$ and $`𝒮`$ cannot be a non-null proper subspace of $`𝒮.`$ Moreover, this intersection cannot be $`𝒮`$ itself, since then $`𝒮`$ and $`𝒮^{}`$ would be orthogonal, contradicting our initial assumption. Thus, the intersection of $`(𝒮^{})^{}`$ and $`𝒮`$ must be the null space. It then follows from the functional relation constraint that $`[P_{(𝒮^{})^{}}][P_𝒮]=0,`$ and since $`[P_𝒮]=1,`$ this implies that $`[P_{(𝒮^{})^{}}]=0.`$ It also follows from the functional relation constraint that $`[P_{(𝒮^{})^{}}]=1[P_𝒮^{}]`$, so that $`[P_𝒮^{}]=1.`$ Thus, both $`P_𝒮`$ and $`P_𝒮^{}`$ receive the value $`1`$. But, this is in contradiction with $`[P_𝒮][P_𝒮^{}]=0`$ which follows from the fact that the intersection of $`𝒮`$ and $`𝒮^{}`$ is the null space. QED.
## IV An Interpretive Framework Incorporating Non-orthogonal Preferred Projectors
### A Preliminaries
In the previous section it was established that in order to consider a modal interpretation with non-orthogonal preferred projectors, the different possible property ascriptions to a system must differ with respect to the ontology they ascribe. The precise form of the ontology associated with a particular property ascription has not yet been specified. It turns out that this form is fixed if an additional constraint on the property ascription is adopted, namely,
If $`P_𝒮Ont`$ and $`[P_𝒮]=1`$ then for all $`P_{}`$ such that $`P_{}P_𝒮,`$ $`P_{}Ont`$ and $`[P_{}]=1.`$
This condition was introduced by Healey and was resurrected recently by Vermaas. It is called ‘weakening’ since in Healey’s terminology $`P_{}`$ is said to be weaker than $`P_𝒮`$ if $`P_{}P_𝒮`$. It is motivated by the same sorts of considerations that lead one to adopt the closure constraint and the functional relation constraint; it is an attempt to preserve the logical structure of classical mechanics. In the language of properties, the weakening condition states that if property $`s`$ is well-defined and holds for the system, then any property implied by $`s,`$ namely any property of the form ‘$`s`$ or $`s^{}`$’ should also be well-defined and hold for the system. It should be noted that the weakening condition is unlike previous constraints, insofar as the nature of the ontology is made to depend on features of the value ascription.
We now demonstrate the form of property ascription that results from adopting the weakening condition.
The algebraic constraints and the weakening condition imply that the set of determinate projectors and the value ascription to these must respectively have the forms
$`\{P_{}|P_{}`$ $``$ $`P_𝒮\text{ or }P_{}P_𝒮\}`$ (7)
$`[P_{}]`$ $`=`$ $`\{\begin{array}{c}1\text{ if }P_{}P_𝒮\\ 0\text{ if }P_{}P_𝒮\end{array},`$ (10)
where $`P_𝒮`$ is the preferred projector for the property ascription.
Proof. Recall that the preferred projector $`P_𝒮`$ is the unique projector in the property ascription satisfying Eq. (4), so that $`[P_𝒮]=1`$ and there is no subspace $`𝒯`$ such that $`P_𝒯<P_𝒮`$ and $`[P_𝒯]=1.`$ By the weakening condition, $`[P_𝒮]=1`$ implies that the set of projectors $`\{P_{}|P_{}P_𝒮\}`$ is determinate. Moreover, for any projector $`P_𝒰`$ orthogonal to $`P_𝒮,`$ $`P_𝒰+P_𝒮`$ is determinate, since $`(P_𝒮+P_𝒰)\{P_{}|P_{}P_𝒮\}.`$ It then follows from the constraint of closure that $`(IP_𝒮)(P_𝒮+P_𝒰)=P_𝒰`$ is determinate. Thus, all projectors orthogonal to $`P_𝒮,`$ namely the set $`\{P_𝒰|P_𝒰P_𝒮\},`$ are also determinate. In summary, all the projectors in the set $`\{P_{}|P_{}P_𝒮`$ or $`P_{}P_𝒮\}`$ must be determinate. We now show that the projectors in this set are the only projectors that are determinate.
Suppose the contrary, namely that there exists a determinate projector $`P_𝒱`$ such that $`P_𝒱P_𝒮`$ and $`P_𝒱⟂̸P_𝒮.`$ If $`P_𝒱P_𝒮`$ then the intersection of $`𝒱`$ and $`𝒮`$ is not equal to $`𝒮`$, and must therefore be a proper subspace of $`𝒮.`$ It then follows from the functional relation constraint and the assumption that all proper subspaces of $`𝒮`$ receive the value $`0`$ that $`[P_𝒱][P_𝒮]=0.`$ Since $`[P_𝒮]=1,`$ we conclude that $`[P_𝒱]=0.`$ Moreover, since $`P_𝒱⟂̸P_𝒮`$ is equivalent to $`(IP_𝒱)P_𝒮,`$ it follows by the same argument that $`[IP_𝒱]=0.`$ But $`[IP_𝒱]=0`$ implies $`[P_𝒱]=1,`$ thereby yielding a contradiction.
Finally, we demonstrate that the value ascription must be of the form of (10). Given that $`[P_𝒮]=1,`$ it follows trivially from the weakening condition that $`[P_{}]=1`$ if $`P_{}P_𝒮.`$ Moreover, $`[P_{}]=0`$ if $`P_{}P_𝒮`$ since otherwise the intersection of $``$ and $`𝒮,`$ which is the null space, would receive the value $`1.`$ QED.
Theorem 2 identifies the set of idempotent variables that are determinate given the weakening condition and the algebraic constraints. The set of non-idempotent variables that are determinate then follows from the spectral constraint. Specifically, we have
If the set of determinate projectors and their values are given by Eqs.(7) and (10), then the spectral constraint implies that the ontology and its value ascription are given by
$`Ont`$ $`=`$ $`\{V|VP_𝒮P_𝒮\}`$ (11)
$`[V]`$ $`=`$ $`Tr\left(VP_𝒮\right),`$ (12)
where $`P_𝒮`$ is the preferred projector for the property ascription.
Proof. The spectral constraint states that a non-idempotent variable $`V`$ is determinate if and only if all the elements of its spectral resolution are determinate. Thus every variable $`V`$ that is determinate has a spectral resolution $`V=_k\lambda _kP__k,`$ where $`k:`$ $`(P__kP_𝒮`$ or $`P__kP_𝒮).`$ But the latter condition is equivalent to $`k:`$ $`P__kP_𝒮P_𝒮,`$ from which it follows that $`VP_𝒮P_𝒮.`$ The spectral constratint also states that $`[V]=_k\lambda _k[P__k]`$. Since the $`P__k`$ are orthogonal, only one can satisfy $`P__kP_𝒮,`$ and thereby receive the value $`1`$ by Eq.(10). Labelling this projector by $`k^{},`$ we have $`[V]=\lambda _k^{},`$ and $`VP_𝒮=\lambda _k^{}P_𝒮,`$ from which Eq.(12) follows. QED.
The corollary to theorem 2 states that a variable is determinate if it has the subspace associated with the preferred projector as an eigenspace, and the value of this variable is the associated eigenvalue. This has the form of the orthodox rule, defined in Eq. (6), but where the role of the density operator is played by the preferred projector.
Note that for systems of dimensionality 3 or greater, theorem 2 implies that the ontologies associated with mutually exclusive property ascriptions are necessarily distinct. This holds true for such systems even if the preferred projectors for the property ascriptions are orthogonalThis is not true for a 2-dimensional Hilbert space, since two distinct property ascriptions can be associated with the same ontology. This occurs when the preferred projectors for these property ascriptions are orthogonal.. In this sense, the weakening condition provides another reason, independent of the one provided in section 3.3, for allowing the possible property ascriptions to differ in ontology. Such an argument was in fact made by Vermaas in the context of the Vermaas and Dieks version of the modal interpretation .
Thus far, we have focused upon the property ascriptions for individual systems, and nothing has been said concerning the relationship between the property ascriptions to composite systems and the subsystems of which they are formed. Clifton has argued for the following constraint on this relationship, which we call the reductionist rule:
$$V^AOnt^A\text{ if and only if }V^AI^BOnt^{AB},\text{ and }[V^A]=[V^AI^B],$$
(13)
where $`Ont^A`$ is the ontology of system $`A,`$ and $`AB`$ is the composite of systems $`A`$ and $`B`$ (i.e. $`A(B)`$ denotes the system associated with Hilbert space $`^A(^B),`$ and $`AB`$ denotes the system associated with $`^A^B).`$ Denying this constraint leads to what Clifton has called ontological perpectivalism, the view that what exists depends on the level of compositeness of the description. For instance, to deny the ‘only if’ half of the rule amounts to claiming that it is possible for part $`A`$ of a composite to have the property $`s`$, while the composite itself does not have the property that part $`A`$ has property $`s`$. Clifton has characterized such a position as ‘metaphysically untenable’.
If one adopts the reductionist rule, the property ascriptions for all subsystems are uniquely fixed by the property ascription for the composite. However, we also wish to assume that the property ascriptions for every system satisfy the algebraic constraints and the weakening condition. It has yet to be demonstrated that these constraints are consistent with the reductionist rule. In fact they are. Specifically, if the property ascription for the composite has the form given in Eqs. (11) and (12), with a preferred projector denoted by $`P_S^{AB},`$ then the property ascription for subsystem $`A`$ also has the form given in Eqs. (11) and (12), where the preferred projector, denoted by $`P_𝒮^A,`$ is the unique projector satisfying
$$P_𝒮^{AB}P_𝒮^AI^B\text{ and there is no subspace }𝒯^A\text{ such that }P_𝒮^{AB}P_𝒯^AI^B<P_𝒮^AI^B.$$
(14)
(In other words, $`P_𝒮^A`$ is the ‘smallest’ projector satisfying $`P_𝒮^AI^BP_𝒮^{AB})`$ This can be shown to be a limiting case of a result by Dickson and Clifton , however for clarity we prove it directly. It suffices to demonstrate the following equivalences:
$`\{P_{}^A|P_{}^AI^BP_𝒮^{AB}\}=\{P_{}^A|P_{}^AP_𝒮^A\}`$
and
$`\{P_{}^A|\text{ }P_{}^AI^BP_𝒮^{AB}\}=\{P_{}^A|P_{}^AP_𝒮^A\},`$
where $`P_𝒮^A`$ and $`P_𝒮^{AB}`$ are related as above. We first demonstrate that the right hand sides imply the left. $`P_{}^AP_𝒮^A`$ trivially implies $`P_{}^AI^BP_𝒮^AI^B,`$ and since by definition, $`P_𝒮^AI^BP_𝒮^{AB},`$ it follows that $`P_{}^AI^BP_𝒮^{AB}.`$ Similarly, $`P_{}^AP_𝒮^A`$ trivially implies $`P_{}^AI^BP_𝒮^AI^B,`$ and together with $`P_𝒮^AI^BP_𝒮^{AB},`$ this implies that $`P_{}^AI^BP_𝒮^{AB}.`$ To show that the left hand sides imply the right, we make use of the fact that $`P_𝒮^AI^B`$ is the ‘smallest’ projector satisfying $`P_𝒮^AI^BP_𝒮^{AB}.`$ This implies that any projector $`P_{}^A`$ satisfying $`P_{}^AI^BP_𝒮^{AB}`$ must also satisfy $`P_{}^AI^BP_𝒮^AI^B,`$ and hence $`P_{}^AP_𝒮^A.`$ In addition, any projector $`P_{}^A`$ satisfying $`P_{}^AI^BP_𝒮^{AB}`$ (equivalently $`(I^AP_{}^A)I^BP_𝒮^{AB})`$ must also satisfy $`(I^AP_{}^A)I^BP_𝒮^AI^B,`$ which implies $`P_{}^AP_𝒮^A.`$ This concludes the proof.
Clearly, the preferred projectors for a subsystem can be non-orthogonal if the preferred projectors for the composite are non-orthogonal. What is perhaps more surprising is that the preferred projectors for a subsystem can be non-orthogonal even if the preferred projectors for the composite are not! For example, suppose the preferred projectors for the composite are two orthogonal projectors, $`P_1^{AB}=P_1^AP_1^B`$ and $`P_2^{AB}=P_2^AP_2^B,`$ where $`P_1^B`$ and $`P_2^B`$ are orthogonal projectors, but $`P_1^A`$ and $`P_2^A`$ are not. It then follows from Eq. (14) that the preferred projectors for $`A`$ are simply $`P_1^A`$ and $`P_2^A,`$ which are non-orthogonal.
In the next subsection, we will introduce a framework for interpretation wherein the possible property ascriptions for the universe are defined first, in accordance with the algebraic constraints and the weakening condition, and the possible property ascriptions for all subsystems are then inferred using the reductionist rule. The preferred projectors for the universe will be assumed to be orthogonal, but as shown above, this is consistent with the preferred projectors for a subsystem being non-orthogonal and hence does not rule out the possibility of satisfying the faithfulness criterion. A more general approach would be to assume a non-orthogonal preferred set for the universe as well. However, the faithfulness criterion does not necessitate this assumption, and indeed, as we will demonstrate in section 5, one can satisfy this criterion for a wide variety of measurements without it. The case of an orthogonal preferred set for the universe is in any event a natural place to begin such an investigation.
The framework that emerges is similar to the one proposed by Bub and Clifton. The most significant difference is in the form of the property ascription, since the latter do not assume the weakening condition. Another difference is in the dynamics of the property ascription. Bub and Clifton defined a dynamics following Vink and Bell. This approach was subsequently generalized in two respects by Bacciagaluppi and Dickson, and Dickson. First, the preferred projectors were allowed to be time-dependent, and secondly it was shown that there is a plurality of possible dynamics consistent with the quantum statistics. We follow the latter, generalized approach.
Since many of the ingredients of the framework derive from a number of sources, and since we introduce some novel terminology, we have written the rest of this section in such a way that it constitutes a self-contained description of the framework.
### B Details of the framework
It is assumed that the universe is associated with a Hilbert space $``$ and a vector $`|\psi (t)`$ that evolves deterministically over time in accordance with the Schrödinger equation,
$$\frac{d}{dt}|\psi (t)=iH|\psi (t),$$
(15)
where $`H`$ is the total Hamiltonian, and where the units are chosen such that $`\mathrm{}=1`$. Since, as will be demonstrated shortly, the role of the vector $`|\psi (t)`$ in the framework is to determine the probabilities of various different property ascriptions as well as their dynamics, it will be dubbed the dynamical state vector.
Define a decomposition $`D`$ of a vector $`|\psi `$ as a set $`\{(c_k,|\varphi _k)\}_{k=1}^m`$ of non-zero coefficients $`c_k`$ and orthonormal vectors $`|\varphi _k`$ such that $`|\psi =_{k=1}^mc_k|\varphi _k.`$ It is assumed that every interpretation within the framework selects a preferred decomposition of the dynamical state $`|\psi (t)`$ at every time $`t.`$ The projectors onto the elements of the preferred decomposition constitute the preferred projectors for the possible property ascriptions to the universe.
We also introduce a new ‘state vector’ that we denote by $`|\mathrm{\Phi }(t).`$ It can be any one of the vector elements of the preferred decomposition. The projector onto this vector is the preferred projector for the property ascription to the universe that obtains at time $`t`$. Assuming the algebraic constraints and the weakening condition, it follows from Theorem 2 that the property ascription for the universe has the form of Eqs. (11) and (12), which may be rewritten in terms of $`|\mathrm{\Phi }(t)`$ as
$`Ont(t)`$ $`=`$ $`\{V|V|\mathrm{\Phi }(t)|\mathrm{\Phi }(t)\},`$
$`[V]_t`$ $`=`$ $`\mathrm{\Phi }(t)\left|V\right|\mathrm{\Phi }(t).`$
The property ascription to any subsystem of the universe is then fixed by the reductionist rule, defined in Eq. (13). Since $`|\mathrm{\Phi }(t)`$ determines the property ascription to every system, we call it property state vector.
Next, we introduce a restriction on the manner in which the elements of the preferred decomposition can evolve over time. Suppose the set of vectors $`\{|\varphi _k(t)\}_{k=1}^d`$ at every time $`t`$ is a complete orthogonal basis for $``$ that includes as a subset the vector elements of the preferred decomposition. We require that there is an indexing of the basis vectors such that every vector with a given index is an analytic function of time. We call this the constraint of analyticity. It can be satisfied by requiring that the time-dependent vectors in the set $`\{|\varphi _k(t)\}_{k=1}^d`$ each define a path through Hilbert space obeying the equation
$$\frac{d}{dt}|\varphi _k(t)=i\stackrel{~}{H}(t)|\varphi _k(t),$$
(16)
for some Hermitian operator $`\stackrel{~}{H}(t).`$ It is convenient to refer to these vectors, considered as functions of time, as the preferred paths.
It is assumed that the property state vector evolves according to a Markovian stochastic dynamics that permits hopping among the preferred paths. We require that at every time $`t`$ the probability $`p_k(t)`$ that the property state vector lies on the $`k`$th preferred path is given by
$$p_k(t)=\left|\varphi _k(t)|\psi (t)\right|^2.$$
(17)
The latter requirement is called the Born rule constraint. Although the basis $`\{|\varphi _k(t)\}_{k=1}^d`$ that is defined at time $`t`$ by the preferred paths may include elements that are not part of the preferred decomposition, these elements have no overlap with $`|\psi (t),`$ so that the probability associated with them is zero. It follows therefore that the property state vector always corresponds to one of the vector elements of the preferred decomposition. There are many dynamics that satisfy the Born rule constraint; these will be considered in the next subsection.
We refer to this entire interpretive structure as a ‘framework’ for modal interpretations, since there are a plurality of possible interpretations that have this form. Specifically, there is a different interpretation for every choice of rule for determining the preferred decomposition and every choice of dynamics that satisfies the Born rule constraint.
### C The general form of the dynamics
We now recall the general form of a Markovian stochastic dynamics that satisfies the Born rule constraint. This constraint, articulated in Eq.(17), can be recast as constraints upon the initial conditions and the dynamics:
$`p_k(0)=\left|\varphi _k(0)|\psi (0)\right|^2,`$
and
$$\frac{d}{dt}p_k(t)=\frac{d}{dt}\left|\varphi _k(t)|\psi (t)\right|^2.$$
(18)
Using Eqs.(15) and (16), the latter becomes
$$\frac{d}{dt}p_k(t)=2Im\left[\psi (t)|\varphi _k(t)\varphi _k(t)|H\stackrel{~}{H}(t)|\psi (t)\right].$$
(19)
Since we assume Markovian dynamics, it is sufficient to specify the probability $`T_{kj}(t)dt`$ of a transition from path $`j`$ to path $`k`$ during the infinitesimal interval between $`t`$ and $`t+dt,`$ for all $`j`$ and $`k.`$ The evolution of a probability distribution $`p_k(t)`$ over the paths is then given by the master equation
$`{\displaystyle \frac{d}{dt}}p_k(t)={\displaystyle \underset{j}{}}\left[T_{kj}(t)p_j(t)T_{jk}(t)p_k(t)\right].`$
In what follows, we consider the problem of finding a set of functions $`T_{kj}`$$`(t)`$ that satisfy the master equation given $`p_k(t).`$ Following Bell, it is useful to define a new set of functions, namely a set of probability currents, $`J_{kj}(t),`$ as follows:
$$J_{kj}(t)=T_{kj}(t)p_j(t)T_{jk}(t)p_k(t).$$
(20)
The current $`J_{kj}(t)`$ describes the net flow of probability from path $`j`$ to $`k`$ at time $`t.`$ This definition implies that the current is antisymmetric with respect to an interchange of its indices
$$J_{kj}(t)=J_{jk}(t).$$
(22)
In terms of these currents, the master equation becomes a continuity equation:
$$\frac{d}{dt}p_k(t)=\underset{j}{}J_{kj}(t).$$
(23)
Following Bacciagaluppi and Dickson, one can solve for the $`T_{kj}(t)`$ in two steps. First, one finds a set of currents $`J_{kj}(t)`$ that satisfy Eq.(22) and that solve Eq.(23) with $`dp_k(t)/dt`$ given by Eq.(19). Next, one finds a set of functions $`T_{kj}(t)`$ that solve Eq.(20) given a particular solution for $`J_{kj}(t).`$ It turns out that there an infinite number of sets of antisymmetric currents which solve the continuity equation. Moreover, for a given set of currents, there are an infinite number of solutions for the $`T_{kj}(t),`$ specifically, any set of functions that satisfy
$$T_{kj}\mathrm{max}\{0,\frac{J_{kj}}{p_j}\},$$
(24)
and
$$T_{jk}=\frac{(T_{kj}p_jJ_{kj})}{p_k}.$$
(26)
for every pair of indices $`k>j`$.
So we see that there is a large number of solutions for the dynamics which satisfy the constraints introduced. It is possible that additional constraints, such as a requirement of quantum-classical correspondence, might eliminate the ambiguity in the choice of dynamics, but this has yet to be demonstrated and some authors argue that it is unlikely.
## V The Minimal Entropy Proposal
### A Details of the proposal
We begin by introducing some terminology. A factorization $`F`$ of a Hilbert space $``$ is defined to be a set of Hilbert spaces each of dimensionality greater than one, the direct product of which is $`,`$ that is, $`F`$ $`=\{^{(p)}\}_{p=1}^n,`$ such that $`=^{(1)}^{(2)}\mathrm{}^{(n)},`$ and $`dim(^{(p)})>1`$. A more precise definition of this concept is supplied by Bacciagaluppi, but this is not required for our purposes. A factorization containing $`n`$ elements is called $`n`$-partite, and the elements themselves are called factor spaces. A factorization $`F`$ is said to be a coarse-graining of a factorization $`F^{},`$ and $`F^{}`$ a fine-graining of $`F,`$ if $`F^{}`$ can be generated from $`F`$ by factorizing one or more of the elements of $`F.`$ Finally, a product decomposition of $`|\psi `$ with respect to the factorization $`F=\{^{(p)}\}_{p=1}^n`$ is any decomposition $`\{(c_k,_{p=1}^n|\varphi _k^{(p)})\}_{k=1}^m`$ of $`|\psi `$ every element of which is a product state over $`F.`$
The first element of the proposal is to assume that there is a factorization of the Hilbert space of the universe that is more physically relevant than the others; we call it the distinguished factorization . There is a precedent for such an assumption, specifically, in the modal interpretations of Healey, Bacciagaluppi and Dickson, and Dieks. Such interpretations have been called ‘atomic’, since the factor spaces of the distinguished factorization represent the most elementary physical systems. Some restrictions on what the distinguished factorization could be will be discussed briefly in section 5.2.
The first constraint upon the preferred decomposition is that it be a product decomposition with respect to the distinguished factorization. This constraint is not sufficient to uniquely specify a decomposition. Indeed, the number of product decompositions of any state vector with respect to a given factorization is infinite. In order to distinguish between these, we turn our attention towards the coefficients in the decomposition. Since these coefficients define a probability distribution, different decompositions can be ordered with respect to the uniformity of the associated distributions. This uniformity can be quantified by several ‘entropic’ quantities. The most obvious candidate is the Shannon entropy, defined for a probability distribution $`𝐩=(p_1,p_2,\mathrm{},p_m)`$ as
$$H(𝐩)=\underset{k=1}{\overset{m}{}}p_k\mathrm{log}p_k.$$
(27)
Thus, we can associate with every decomposition $`D=\{(c_k,|\varphi _k)\}_{k=1}^m`$ of the state vector $`|\psi `$ the entropy
$$S_{|\psi }(D)=\underset{k=1}{\overset{m}{}}\left|c_k\right|^2\mathrm{log}\left|c_k\right|^2.$$
(28)
We refer to this quantity as the $`\mathrm{I}\mathrm{U}`$ entropy of the state vector $`|\psi `$ for the decomposition $`D,`$ since it has previously been considered by Ingarden and Urbanik, albeit in a very different context.
It is now possible to state our choice of preferred decomposition:
> Given a distinguished factorization $`F`$, the preferred decomposition of the dynamical state vector $`|\psi `$ is the one that minimizes the IU entropy of $`|\psi `$ from among all product decompositions with respect to $`F`$.
By choosing the product decomposition that minimizes the IU entropy, we are choosing the interpretation where the probability distribution over the possible property state vectors is as narrow as possible. Moreover, since the minimum IU entropy (from among IU entropies for product decompositions) is zero if and only if $`|\psi `$ is a product state, it can be thought of as a measure of the entanglement of $`|\psi `$ with respect to the distinguished factorization. The strongest motivation for such a choice of preferred decomposition is that it appears very promising in securing a solution to the measurement problem and in satisfying the faithfulness criterion, as will be demonstrated in sections 5.2 and 5.3. We do not however offer any a priori justification of the principle.
Implementing the proposal requires solving the minimization problem for a given dynamical state vector and a given choice of distinguished factorization. If the distinguished factorization is bi-partite, the solution is given by the following theorem.
Suppose $`|\psi `$ is any vector in a Hilbert space with a bi-partite distinguished factorization $`F_{\text{bi}}.`$ Any decomposition of $`|\psi `$ that is bi-orthogonal with respect to $`F_{\text{bi}}`$ minimizes the IU entropy from among all product decompositions of $`|\psi `$ with respect to $`F_{\text{bi}}.`$
The proof of this theorem is relegated to appendix A. In the case of an $`n`$-partite distinguished factorization, with $`n>2,`$ we have not yet found a solution to the minimization problem for all state vectors. However, the bi-partite result can be used to identify the preferred decomposition for some state vectors, as follows.
Suppose $`|\psi `$ is a vector in a Hilbert space with an $`n`$-partite distinguished factorization, $`F_n,`$ where $`n>2.`$ If there exists a decomposition of $`|\psi `$ that is a product decomposition with respect to $`F_n`$ and that is a bi-orthogonal decomposition with respect to some bi-partite coarse-graining of $`F_n`$, then this decomposition minimizes the IU entropy from among all product decompositions with respect to $`F_n`$.
Proof. Suppose $`F_{\text{bi}}`$ is a bi-partite coarse-graining of $`F_n.`$ The set $`S_n`$ of decompositions of $`|\psi `$ that are product decompositions with respect to $`F_n`$ is a subset of the set $`S_{\text{bi}}`$ that are product decompositions with respect to $`F_{\text{bi}}.`$ We can denote this by $`S_nS_{\text{bi}}.`$ Moreover, suppose $`D_n^{\mathrm{min}}(D_{\text{bi}}^{\mathrm{min}})`$ is the decomposition that minimizes the IU entropy from among all the elements of $`S_n(S_{\text{bi}}).`$ Theorem 3 shows that for every state vector $`|\psi `$, $`D_{\text{bi}}^{\mathrm{min}}`$ is the bi-orthogonal decomposition of $`|\psi `$. For certain state vectors, it may happen that $`D_{\text{bi}}^{\mathrm{min}}`$ lies among the elements of $`S_n.`$ Since we know that $`D_{\text{bi}}^{\mathrm{min}}`$ minimizes the IU entropy from among all the elements of $`S_{\text{bi}},`$ and $`S_nS_{\text{bi}},`$ it follows that in this case $`D_{\text{bi}}^{\mathrm{min}}`$ also minimizes the IU entropy from among all the elements of $`S_n.`$ Thus, in this case $`D_n^{\mathrm{min}}=D_{\text{bi}}^{\mathrm{min}}.`$ QED.
Theorem 4 is not a complete solution to the minimization problem because there exist state vectors for which $`D_{\text{bi}}^{\mathrm{min}}`$ does not lie among the elements of $`S_n.`$ Further work is required to determine the decomposition that minimizes the IU entropy in such cases.
We note that in the proof of theorem 3, presented in appendix A, the only relevant feature of the IU entropy is that it has the form $`_{k=1}^mf(\left|c_k\right|^2)`$ for some concave function $`f.`$ It follows that one would obtain the same results if, instead of minimizing the IU entropy, one minimized any other entropic quantity having this form. However, there is no guarantee that this insensitivity to the choice of entropic quantity persists in the more general case of state vectors for which theorem 4 does not apply.
A possible difficulty with the minimal entropy proposal as it stands has to do with the uniqueness of the preferred decomposition. It is well known that the bi-orthogonal decomposition of a state vector is not unique when the eigenvalues of the reduced density operator for one of the subsystems are degenerate. It follows from theorem 3 that if the distinguished factorization is bi-partite, then the decomposition that minimizes the IU entropy may not be unique, and the minimal entropy proposal may fail to uniquely specify a preferred decomposition. For instance, this occurs if the dynamical state vector is the EPR-Bell state for two spins $`|\psi =2^{1/2}(|+𝐚|𝐚+|𝐚|+𝐚).`$ This difficulty persists in the case of an $`n`$-partite distinguished factorization, $`F_n,`$ where $`n>2`$, since there are dynamical state vectors for which theorem 4 applies and the decomposition that minimizes the IU entropy is non-unique; an example being a tensor product of EPR-Bell states. It should be noted however that a degeneracy among the eigenvalues of the reduced density operator for one of the factor spaces of $`F_n`$ does not always lead to a non-unique preferred decomposition. For instance, if the dynamical state vector has a decomposition that is $`n`$-orthogonal with respect to the factorization $`F_n,`$ then it follows from theorem 4 that this decomposition minimizes the IU entropy, and since the $`n`$-orthogonal decomposition is unique for $`n>2`$ , so is the preferred decomposition. It is an open question whether the minimization of the IU entropy leads to a unique preferred decomposition when the dynamical state vector is such that theorem 4 does not apply.
It is useful to distinguish two cases of non-uniqueness of the preferred decomposition: an instantaneous non-uniqueness, occurring at an isolated moment in time, and an extended non-uniqueness, occurring over a finite interval of time. If the constraint of analyticity (defined in Eq.(16)) holds for the minimal entropy proposal, then the instantaneous non-uniqueness problem can be solved easily: the preferred paths at the moment of non-uniqueness are simply taken to be the limit of the preferred paths at adjoining times. This is the same solution as was proposed in the context of the atomic modal interpretation by Bacciagaluppi and Dickson. The extended non-uniqueness problem is not so easily solved. One possible approach to the problem is to argue that cases wherein there is an extended non-uniqueness have negligible probability. Since such an argument has been made for the occurrence of a non-unique bi-orthogonal decomposition by Bacciagaluppi, Donald and Vermaas , this result can be applied to the minimal entropy proposal in cases where theorem 4 applies.
Finally, we turn to the issue of dynamics. Given theorem 4, it is possible to show that the minimal entropy proposal satisfies the constraint of analyticity in some cases. In particular, if the dynamical state vector evolves in such a way that it has a bi-orthogonal decomposition with respect to some coarse-graining of the distinguished factorization for a finite interval of time, then Eq.(16) can be satisfied for that interval. The reason is that the vector elements of a bi-orthogonal decomposition are analytic functions of time, as has been shown by Bacciagaluppi and Dickson. It remains an open question whether for arbitrary dynamical state vectors the decomposition that minimizes the IU entropy, considered as a function of time, satisfies the analyticity constraint. If this is indeed the case, then the entropy minimization rule defines a set of preferred paths.
Given such a set of paths, denoted by $`\{|\varphi _k(t)\}_{k=1}^d,`$ we must choose the form of the dynamics from among all possible solutions for $`J_{kj}(t)`$ and $`T_{kj}(t)`$ in Eqs. (20), (22) and (23). We follow Bacciagaluppi and Dickson in choosing:
$$J_{kj}(t)=2Im\left[\psi (t)|\varphi _k(t)\varphi _k(t)|H\stackrel{~}{H}(t)|\varphi _j(t)\varphi _j(t)|\psi (t)\right],$$
(29)
and
$$T_{kj}(t)=\mathrm{max}\{0,\frac{J_{kj}(t)}{p_j(t)}\}.$$
(30)
This is a generalization to time-dependent preferred decompositions of the choice made by Bell, Vink and Bub. Since the inequality in Eq.(24) is saturated, this choice of $`T_{kj}(t)`$ minimizes the degree of stochasticity for a given form of the current. Such a choice is motivated by the fact that classical mechanics, which is deterministic, must be obtained as a limit of quantum mechanics.
### B The quantum measurement problem
We now consider whether the minimal entropy proposal solves the quantum measurement problem. Although this term is often taken to refer to the whole cluster of conceptual difficulties surrounding measurement, we shall use it to refer to the particular problem of deriving operational quantum mechanics from a realist no-collapse interpretation. To consider the problem, we must introduce a quantum mechanical model of the measurement procedure, that is, a model of the interaction between the degrees of freedom of the system under investigation, the apparatus, and the environment. We discuss both single measurements and sequences of measurements.
#### 1 Single measurements
Following the notation introduced in section 3.1, we consider the measurement of a Hermitian operator $`V,`$ belonging to a Hilbert space $`^S,`$ the eigenvalues of which are non-degenerate and the eigenvectors of which are denoted by $`\{|\phi _k\}_{k=1}^m`$. Assuming the preparation procedure is associated with a state vector $`_kc_k|\phi _k,`$ where $`_k\left|c_k\right|^2=1`$, operational quantum mechanics predicts, via the Born rule, that the measurement will have outcome $`k`$ with probability $`\left|c_k\right|^2.`$
We now consider a quantum mechanical model of the measurement process. The system under investigation is called the object system and is assumed to be microscopic. This is made to interact with a macroscopic apparatus, associated with a Hilbert space $`^A,`$ which in turn interacts with a macroscopic environment, associated with a Hilbert space $`^E`$. Given an initial state vector in $`^S^A^E,`$ one could in principle determine the evolution of the total system using the full microscopic Hamiltonian.
In practice of course the problem is far too complex to be solved exactly. Nonetheless, there is a set of standard toy models of measurement that are commonly used to investigate realist interpretations. These models adopt some simplifying assumptions about the initial state and the form of the evolution. Specifically, it is assumed that the object system, apparatus and environment are all initially uncorrelated, so that the initial dynamical state vector has the form $`|\phi _k|A_0|E_0,`$ a product state with respect to the factorization $`\{^S,^A,^E\}`$ of the Hilbert space. The dynamics is assumed to be such that
$$|\phi _k|A_0|E_0|\stackrel{~}{\phi }_k|A_k|E_k,$$
(31)
where {$`|A_k\}_{k=1}^m`$ is a set of orthonormal vectors for the apparatus, {$`|E_k\}_{k=1}^m`$ is a set of orthonormal vectors for the environment, and {$`|\stackrel{~}{\phi }_k\}_{k=1}^m`$ is a set of normalized but possibly non-orthogonal vectors for the object system, and where ‘$``$’ denotes the mapping corresponding to the unitary evolution.
If the initial state for the object system is $`_kc_k|\phi _k`$, the final dynamical state vector for the total system, given Eq.(31) and the assumption that the evolution is linear, is
$$|\psi _{\text{final}}=\underset{k=1}{\overset{m}{}}c_k|\stackrel{~}{\phi }_k|A_k|E_k.$$
(32)
We are now in a position to ask whether a given realist no-collapse interpretation falls prey to the quantum measurement problem within this model. We begin by illustrating the problem in the traditional manner, specifically, in the context of the simplest realist no-collapse interpretation one can imagine: one where the property ascriptions for systems are fixed by the orthodox rule, defined in Eq.(6). Such an interpretation has been called the ‘bare theory’ by Albert. Within the framework of section 4, it corresponds to adopting the trivial decomposition of the dynamical state vector as preferred (the trivial decomposition of $`|\psi (t)`$ is simply $`\{(1,|\psi (t))\}`$).
Consider first a case where $`c_k0`$ for only a single value of $`k,`$ that is, where the initial state vector of the object system is an eigenstate of $`V.`$ The final state vector is then of the form $`|\stackrel{~}{\phi }_k|A_k|E_k`$. By the orthodox rule and the reductionist rule, the preferred projector for the property ascription to the apparatus is $`P_{A_k}.`$ If the bare theory is to reproduce the predictions of operational quantum mechanics in this case, then the property associated with the projector $`P_{A_k}`$ must be such that the apparatus can be accurately described as ‘indicating outcome $`k`$’ (for instance, if the apparatus indicates the outcome by a digital display, $`P_{A_k}`$ could correspond to the property of displaying the number $`k`$). We refer to this as the assumption of ontological correspondence.
If, on the other hand, the initial state is such that $`c_k0`$ for more than one value of $`k,`$ then the final state vector is of the form $`_k^{}c_k|\stackrel{~}{\phi }_k|A_k|E_k,`$ where $`_k^{}`$ indicates a sum over values of $`k`$ for which $`c_k0.`$ In this case, the preferred projector for the property ascription to the apparatus is $`_k^{^{}}P_{A_k},`$ while no projector of the form $`P_{A_k}`$ receives the value 1. Thus, even given the assumption of ontological correspondence, the bare theory does not predict that the apparatus indicates the outcome $`k`$ for any value of $`k`$ for which $`c_k0.`$ Hence the bare theory does not reproduce the predictions of operational quantum mechanics. This is the quantum measurement problem.
We now specify the assumptions under which the minimal entropy proposal solves this problem. These involve the nature of the distinguished factorization, which we have not yet specified. Whatever it might be, the distinguished factorization should be defined in terms of primitives of the theory and selected by physical principles, for instance, from considerations of symmetry. We do not here present an argument for the identity of the distinguished factorization, however a discussion of the issue can be found in Dieks, wherein it is argued that a necessary condition on this choice is that the factor spaces carry an irreducible representation of the space-time group (the Galilei group in nonrelativistic quantum mechanics). For the present, we insist only that the distinguished factorization, which we denote by $`F,`$ has the factorization $`\{^S,^A,^E\}`$ as a coarse-graining, and that its elements correspond to microscopic degrees of freedom (for instance, they could correspond to degrees of freedom of elementary particles).
Now, suppose that $`F`$ is such that all the vectors in the sets $`\{|\stackrel{~}{\phi }_k\}_{k=1}^m,`$ $`\{|A_k\}_{k=1}^m,`$ and $`\{|E_k\}_{k=1}^m`$ are product states with respect to it. One can then determine that the preferred decomposition of $`|\psi _{\text{final}}`$ is
$`D_{\text{final}}=\left\{(c_k,|\stackrel{~}{\phi }_k|A_k|E_k)\right\}_{k=1}^m.`$
This follows from theorem 4 and the fact that $`D_{\text{final}}`$ is a product decomposition with respect to $`F`$ that is bi-orthogonal with respect to the coarse-graining $`F_{\text{bi}}=\{^S^A,^E\}`$ of $`F.`$ It then follows from the Born rule constraint that, with probability $`\left|c_k\right|^2,`$ the property state vector is
$`|\mathrm{\Phi }_{\text{final}}=|\stackrel{~}{\phi }_k|A_k|E_k.`$
Using the reductionist rule we find that the projector $`P_{A_k}`$ is determinate and receives the value $`1`$ with probability $`\left|c_k\right|^2.`$ Finally, by the assumption of ontological correspondence, the apparatus has the property of indicating outcome $`k`$ with probability $`\left|c_k\right|^2.`$ This is in agreement with the predictions of operational quantum mechanics.
Thus, we have obtained a solution to the measurement problem within the standard model of measurement. In so doing, we have had to assume that when the initial state vector of the object system is an eigenstate of $`V,`$ the final dynamical state vector for the total system is unentangled with respect to the distinguished factorization.
The assumption of no entanglement between the distinguished factor spaces of the apparatus and the environment is not particularly realistic, given that these factor spaces are taken to correspond to microscopic degrees of freedom, and typical interactions between the apparatus and the environment are likely to entangle these degrees of freedom. However, this assumption can be relaxed somewhat without changing any of our conclusions, as we now demonstrate.
We consider a model of measurement wherein the evolution is of the form:
$$|\phi _k|A_0|E_0|\stackrel{~}{\phi }_k\underset{\mu =1}{\overset{M}{}}f_{k,\mu }\left(|A_{k,\mu }|E_{k,\mu }\right).$$
(33)
where $`_\mu \left|f_{k,\mu }\right|^2=1`$ and where $`\{\{|A_{k,\mu }\}_{\mu =1}^M\}_{k=1}^m`$ and $`\{\{|E_{k,\mu }\}_{\mu =1}^M\}_{k=1}^m`$ are orthonormal sets of vectors that are product states with respect to $`F.`$ The assumption of ontological correspondence in this case becomes the assumption that for every value of $`\mu ,`$ the projector $`P_{A_{k,\mu }}`$ corresponds to the apparatus indicating outcome $`k`$. That there can be more than one projector corresponding to indicating a particular outcome is not unreasonable since there can be many different microscopic configurations of the apparatus leading to the same overall macroscopic appearance.
An arbitrary initial state vector for the object system, $`_kc_k|\phi _k,`$ leads, via Eq.(33), to the following final dynamical state vector for the total system
$`|\psi _{\text{final}}^{}={\displaystyle \underset{k=1}{\overset{m}{}}}c_k|\stackrel{~}{\phi }_k{\displaystyle \underset{\mu =1}{\overset{M}{}}}f_{k,\mu }\left(|A_{k,\mu }|E_{k,\mu }\right).`$
The preferred decomposition of this state vector is
$`D_{\text{final}}^{}=\left\{\left\{(c_kf_{k,\mu },|\stackrel{~}{\phi }_k|A_{k,\mu }|E_{k,\mu })\right\}_{\mu =1}^M\right\}_{k=1}^m,`$
since this is a product decomposition with respect to $`F`$ that is bi-orthogonal with respect to $`F_{\text{bi}}.`$ It follows that the property state vector is
$`|\mathrm{\Phi }_{\text{final}}^{}=|\stackrel{~}{\phi }_k|A_{k,\mu }|E_{k,\mu }`$
with probability $`\left|c_kf_{k,\mu }\right|^2.`$ By the reductionist rule, the projector $`P_{A_{k,\mu }}`$ is determinate and receives value $`1`$ with probability $`\left|c_kf_{k,\mu }\right|^2.`$ Finally, by the assumption of ontological correspondence, the apparatus has the property of indicating outcome $`k`$ with probability $`_\mu \left|c_kf_{k,\mu }\right|^2=\left|c_k\right|^2,`$ in agreement with operational quantum mechanics.
Note that the model of measurement provided by Eq.(33) can also describe error-prone measurements. This occurs if for some values of $`\mu ,`$ $`P_{A_{k,\mu }}`$ corresponds to the property of indicating an outcome $`k^{}k,`$ or to the property of indicating a malfunction. Furthermore, this model can incorporate measurements described by positive operator-valued measures(POVMs). This follows from the fact that such measurements are implemented by adjoining an ancilla to the system under investigation and measuring a projector-valued measure(PVM) on the composite. By including the ancilla in our definition of the object system, the model presented above can describe these measurements. Note however that we are restricted to PVMs whose eigenvectors are product states with respect to the distinguished factorization.
Despite the possibility of incorporating some error-prone and POVM measurements, the model of measurement provided by Eq.(33) is still not the most general or realistic. Although it is true that an arbitrary state vector has many decompositions into product states with respect to the distinguished factorization, it is not necessarily the case that any of these decompositions are bi-orthogonal with respect to a coarse-graining of the distinguished factorization. For instance, if any of the vectors in the set $`\{|\stackrel{~}{\phi }_k\}_{k=1}^m`$ are entangled with respect to the distinguished factor spaces of $`^S,`$ then theorem 4 fails to apply if the final dynamical state vector is of the form of $`|\psi _{\text{final}}^{}`$. Since the problem of minimizing the IU entropy for arbitrary state vectors has not yet been solved, it is not clear what the preferred decomposition will be in this case and whether the measurement problem is resolved or not.
It is nonetheless interesting to consider one particular type of modification of the evolution where the only change from the model considered above is that the set of vectors $`\{\{|E_{k,\mu }\}_{\mu =1}^M\}_{k=1}^m`$ (describing the states of the environment that are relative to the apparatus states $`\{\{|A_{k,\mu }\}_{\mu =1}^M\}_{k=1}^m`$) is only approximately orthogonal. This is an instance where theorem 4 may fail to apply. However, the difference between $`|\psi _{\text{final}}^{}`$ when the elements of $`\{\{|E_{k,\mu }\}_{\mu =1}^M\}_{k=1}^m`$ are orthogonal and when they are very nearly orthogonal, is not significant. Thus, if the preferred decomposition does not depend sensitively on small variations in the dynamical state vector, the preferred decomposition in the nearly orthogonal case should be ‘close to’ $`D_{\text{final}}^{},`$ and it is then likely that the apparatus will be assigned an ontology that is ‘close to’ the one it receives for the orthogonal case. We see therefore that whether or not there is a measurement problem in this case depends on whether or not there is such sensitive dependence. The answer to this question must await further progress on the problem of the minimization of the IU entropy.The analagous question in the Vermaas-Dieks version of the modal interpretation is whether the spectral resolution of a density operator is sensitive to small changes in the density operator. Bacciagaluppi, Donald and Vermaas have shown that this does in fact occur when the density operator has nearly degenerate eigenvalues.
Finally, we note that the assumption that the apparatus and environment are initially unentangled is also an unrealistic feature of the standard model of measurement. For that matter, the assumption that the composite of system, apparatus and environment is unentangled with the rest of the universe may not be realistic either. However, this difficulty is not unique to the minimal entropy proposal. Every realist no-collapse interpretation must contend with the fact that the dynamical state vector for the universe is in general not factorizable with respect to subsystems that have interacted in the past, even if this interaction is quite weak. Further work is required to determine whether the predictions of the minimal entropy proposal remain satisfactory when these assumptions are relaxed.
#### 2 Sequences of measurements
We now demonstrate the extent to which the minimal entropy proposal is in agreement with operational quantum mechanics for sequences of measurements. Consider in particular the sequence of two measurements described in section 3.1. Recall that the first measurement is of a variable $`V`$ with eigenstates $`\{|\phi _k\}_{k=1}^m`$, the second measurement is of a variable $`V^{}`$ with eigenstates $`\{|\phi _k^{}\}_{k=1}^m,`$ and the state prepared by the first apparatus given outcome $`k`$ is denoted by $`|\stackrel{~}{\phi }_k.`$
In the last subsection we considered two distinct models of measurement which differed in the extent to which the apparatus and the environment became entangled due to their interaction. In this subsection, we consider only the simpler of the two models. The reader can verify that the more realistic model leads to the same conclusions.
There are now two apparatuses, and an environment for each. We denote their Hilbert spaces by $`^{A1},^{A2},^{E1},`$ and $`^{E2}`$ respectively, and we distinguish state vectors for the two apparatuses(environments) by a superscript. It is again assumed that the object system, the two apparatuses and the two environments are all initially uncorrelated. The distinguished factorization $`F`$ is assumed to have $`\{^S,^{A1},^{A2},^{E1},`$ $`^{E2}\}`$ as a coarse-graining.
We assume that the first measurement is well described by Eq.(31) with the exception of a change of notation: $`|A_0,|E_0,|A_k`$ and $`|E_k`$ become $`|A_0^1,|E_0^1,|A_k^1`$ and $`|E_k^1`$ in order to specify that the object system interacts with the first rather than the second apparatus. We assume that the second apparatus and its environment remain uncorrelated with the rest of the system and each other during this first measurement. It follows that the dynamical state vector for the total system after the first measurement is
$$|\psi _{\text{final 1}}=\left(\underset{k=1}{\overset{m}{}}c_k|\stackrel{~}{\phi }_k|A_k^1|E_k^1\right)|A_0^2|E_0^2.$$
(34)
Suppose that $`|A_0^2`$ and $`|E_0^2`$ are product states with respect to $`F.`$ If we make all the same assumptions about $`|A_k^1`$ and $`|E_k^1`$ as were made for $`|A_k`$ and $`|E_k`$ in the previous subsection, and if we use the bi-partite factorization $`\{^S^{A1}^{A2},^{E1}^{E2}\}`$ in place of the bi-partite factorization $`\{^S^A,^E\}`$ in the arguments found therein, then it is straightforward to show that the preferred decomposition of $`|\psi _{\text{final 1}}`$ is
$$D_{\text{final 1}}=\left\{(c_k,|\stackrel{~}{\phi }_k|A_k^1|E_k^1|A_0^2|E_0^2)\right\}_{k=1}^m.$$
(35)
We conclude that with probability $`\left|c_k\right|^2,`$ the first apparatus indicates outcome $`k,`$ while the second apparatus remains ready to measure.
Now assume that the second measurement is also well described by Eq.(31) with the notational change that $`|A_0,|E_0,|A_k`$ and $`|E_k`$ become $`|A_0^2,|E_0^2,|A_k^2`$ and $`|E_k^2`$ since the object system is now interacting with the second apparatus, and where the vectors for the object system acquire a prime since the second measurement is of $`V^{}`$ rather than $`V.`$ For simplicity, we take this second measurement to be non-disturbing, so that $`|\stackrel{~}{\phi }_k^{}=|\phi _k^{}.`$ Assume also that the first apparatus and its environment have no interactions during this measurement. It then follows that the dynamical state vector after the second measurement is
$$|\psi _{\text{final 2}}=\underset{k=1}{\overset{m}{}}\underset{j=1}{\overset{m}{}}c_kd_j^k|\phi _j^{}|A_k^1|E_k^1|A_j^2|E_j^2,$$
(36)
where the coefficients $`\{d_j^k\}_{k=1}^m`$ are defined by $`|\stackrel{~}{\phi }_k=_{k=1}^md_j^k|\phi _j^{}.`$ Again, if we make all the same assumptions about $`|A_j^2`$ and $`|E_j^2`$ as were made for $`|A_j`$ and $`|E_j`$ in the previous subsection, and if we use the bi-partite factorization $`\{^S^{A1}^{A2},^{E1}^{E2}\}`$ in place of the bi-partite factorization $`\{^S^A,^E\}`$ in all the arguments found therein, the preferred decomposition of $`|\psi _{\text{final 2}}`$ is found to be
$$D_{\text{final 2}}=\left\{\left\{(c_kd_j^k,|\phi _j^{}|A_k^1|E_k^1|A_j^2|E_j^2)\right\}_{k=1}^m\right\}_{j=1}^m.$$
(37)
We can therefore conclude that there is a probability $`\left|c_kd_j^k\right|^2`$ that the first apparatus indicates outcome $`k`$ and the second apparatus indicates outcome $`j`$ after the second measurement. It follows that the probability of the second apparatus indicating outcome $`j`$ given that the first apparatus indicates outcome $`k`$ after the second measurement is $`\left|d_j^k\right|^2=\left|\phi _j|\stackrel{~}{\phi }_k\right|^2.`$
However, we have still not determined the probability for the second apparatus to indicate outcome $`j`$ after the second measurement given that the first apparatus indicates outcome $`k`$ after the first measurement, which is the quantity specified by the generalized Born rule. The problem is that it has not been shown that the outcome indicated by the first apparatus is stable over time. Whether it is or not depends on the dynamics of the property state vector, which is determined by Eqs. (29) and (30). Now although it may be reasonable to assume that the dynamical state vector after a measurement is such that theorem 4 applies, it is unlikely that this theorem applies during the entire interaction leading up to this outcome. Given this, we cannot at present determine the time-sequence of preferred decompositions nor the preferred paths through Hilbert space defined by this sequence. Since Eqs. (29) and (30) depend on the identity of these preferred paths, we cannot at present determine the dynamics of the property state vector.
Thus, for the moment we simply assume that within the minimal entropy proposal, the apparatus is never described as ‘jumping’ between macroscopically different readings. We call this assumption stability. Given stability, the minimal entropy proposal reproduces the predictions of the generalized Born rule.
### C The Faithfulness criterion revisited
We now reconsider the experiment of section 3.1 in the context of the minimal entropy proposal. Since this experiment involves a sequence of two measurements, we can make use of the model of measurement presented in the previous section. Consider the property ascription at the time $`t,`$ after the first measurement. The dynamical state vector is $`|\psi _{\text{final 1}}`$, defined in Eq.(34), and its preferred decomposition is $`D_{\text{final 1}},`$ defined in Eq.(35). If the first apparatus indicates the outcome $`k`$ at time $`t`$, the property state vector must be the $`k`$th element of $`D_{\text{final 1}},`$ that is, $`|\mathrm{\Phi }_{\text{final 1}}=|\stackrel{~}{\phi }_k|A_k^1|E_k^1|A_0^2|E_0^2.`$ The critical feature of the experiment of section 3.1 is that the vector $`|\stackrel{~}{\phi }_k`$ that is prepared when the first apparatus indicates outcome $`k,`$ is an eigenstate of $`V_{(k)}`$, the variable measured by the second apparatus. It follows that the variable $`V_{(k)}I`$ (where $`I`$ is the identity operator for $`^{A1}^{A2}^{E1}^{E2})`$ is determinate and has value $`v_{(k),1}`$ at time $`t.`$ Finally, it follows from the reductionist rule that $`V_{(k)}`$ is determinate and has value $`v_{(k),1}`$ at time $`t.`$ This is precisely what is required in order for the faithfulness criterion to be satisfied.
It should be noted that since the variable measured by the second apparatus depends upon the outcome of the first measurement, the initial state of the second apparatus may well be different for different outcomes of the first measurement. Thus, rather than the first measurement interaction being described by Eq.(31), it may be described by
$`|\phi _k|A_0^1|E_0^1|A_0^2|E_0^2`$ (38)
$``$ $`|\stackrel{~}{\phi }_k|A_k^1|E_k^1|A_{(k),0}^2|E_{(k),0}^2,`$ (39)
where $`\{|A_{(k),0}^2\}_{k=1}^m`$ and $`\{|E_{(k),0}^2\}_{k=1}^m`$ are orthonormal sets of vectors, and $`|A_{(k),0}^2`$ corresponds to the apparatus being ready to measure the variable $`V_{(k)}.`$ In any event, by making the same assumptions for $`|A_{(k),0}^2`$ and $`|E_{(k),0}^2`$ as were made for $`|A_0^2`$ and $`|E_0^2`$ in the last subsection, one can show that the minimal entropy proposal is in agreement with the predictions of operational quantum mechanics even when the nature of the second measurement depends on the outcome of the first.
We end this section with a discussion of the case wherein the second measurement is of a variable whose eigenstates are not all product states with respect to the distinguished factorization $`F`$. For such measurements, the faithfulness criterion cannot be satisfied within the minimal entropy proposal. The reason is as follows. Suppose $`|\phi `$ is an eigenstate of the measured variable that is entangled with respect to $`F.`$ If $`|\phi `$ is prepared by the first measurement and measured by the second, then the faithfulness criterion requires that the projector $`P_\phi `$ be determinate with value $`1`$ immediately prior to the second measurement. However, for this to occur the property state vector must be an eigenstate of $`P_\phi ,`$ and hence must be entangled with respect to $`F`$. But the property state vector is always a product state with respect to $`F`$ in the minimal entropy proposal.
The failure of the faithfulness criterion for such measurements in the context of the atomic modal interpretation of Bacciagaluppi and Dickson and Dieks has been discussed by Dieks, and also by Vermaas. These authors have suggested that an explanation of the outcomes of these measurements might be provided by dispositional properties or collective effects of the composite. This explanation can also be invoked in the context of the minimal entropy proposal.
## VI Conclusions
In modal interpretations, the properties of a system are given by a specification of the set of determinate variables (the ontology) and the value ascription to these variables, jointly referred to as the property ascription. Such interpretations also assume that the property ascription which obtains at a given time is just one of several possibilities. There is always a unique ‘smallest’ projector which receives the value $`1`$ in each of these possible property ascriptions, which we call the preferred projector for that property ascription.
We have shown that these preferred projectors must be non-orthogonal if one seeks to satisfy the faithfulness criterion, that is, if one seeks to explain the outcomes of certain perfectly predictable measurements in terms of pre-existing properties of the system under investigation. The possibility of such an explanation has historically been a strong motivation for the modal approach.
We have also shown that non-orthogonal preferred projectors are inconsistent with the assumption, common among previous modal interpretations, that at a given time there is only a single possible ontology. In order to consider non-orthogonal preferred projectors, we have developped a framework for modal interpretations wherein at a given time, the possible property ascriptions may differ with respect to ontology. As is required for any modal interpretation, the state vector appearing in the Schrödinger equation, which we call the dynamical state vector, does not uniquely fix the property ascription. Rather, a preferred decomposition of the dynamical state vector into a sum of orthogonal vectors must be specified at every time, and a single element of this decomposition, dubbed the property state vector, fixes the property ascription. The property state vector evolves stochastically according to a Markovian dynamics. Finally, subsystems receive only those properties they inherit from the total system by the reductionist rule.
It is of course possible to generalize this framework in many ways. One could consider non-Markovian dynamics, alternatives to the reductionist rule, and even non-orthogonal decompositions of the dynamical state vector. Nonetheless, we feel that the framework presented is a natural starting place for the interpretive program at hand.
Within the context of this framework, we have presented a novel proposal for the preferred decomposition. The proposal assumes that there is a distinguished set of subsystems of the universe, that is, a distinguished factorization of the total Hilbert space into a tensor product of Hilbert spaces<sup>§</sup><sup>§</sup>§Other modal interpreters have been led to this assumption by considerations of the correlations between the properties of a system and its subsystems.. It is also assumed that the preferred decomposition is a product decomposition with respect to this factorization. In the case of a distinguished factorization that is bi-partite, it is then natural to follow previous authors in identifying the bi-orthogonal decomposition as preferred. However, the obvious generalization of the bi-orthogonal decomposition to an $`n`$-partite distinguished factorization, namely the $`n`$-orthogonal decomposition, does not exist for all state vectors, as shown by Peres. The preferred decomposition in our proposal is the one that minimizes the IU entropy from among all product decompositions with respect to the distinguished factorization. This decomposition always exists and turns out to be equal to the $`n`$-orthogonal decomposition when the latter exists. It therefore can be thought of as a natural generalization of the bi-orthogonal decomposition to $`n`$-partite systems.
At present the strongest justification for the minimal entropy proposal is its success in dealing with the quantum measurement problem and in satisfying the faithfulness criterion. The measurement problem is resolved for a wide variety of measurements including certain types of non-ideal measurements, in particular, disturbing measurements, assuming particular microscopic models of the apparatus and environment. Within the same microscopic models, the faithfulness criterion is satisfied for sequences of disturbing measurements. It is this feature of the minimal entropy proposal that sets it apart from previous modal interpretations.
The solution of the measurement problem relies on the assumption of ontological correspondence, that the ontology of macroscopic systems corresponds to our everyday perceptions of them, and the assumption of stability, that the dynamics of the properties assigned to macroscopic systems are consistent with our stable perceptions of them. Ideally, these features would be demonstrated rather than assumed. However, the demonstration of ontological correspondence is likely to require a better specification of the distinguished factorization than has been provided in the present work, while the demonstration of stability must await progress in solving the entropy minimization problem in cases where theorem 4 does not apply. Progress on the minimization problem will also help to determine whether one can solve the quantum measurement problem for more general types of measurements than the ones considered here, for example, measurements of variables whose eigenstates are entangled with respect to the distinguished factorization. In addition, such progress is required to determine what the proposal has to say about more realistic models of measurements. Finally, it may indicate whether the IU entropy is the correct quantity to minimize in the rule for determining the preferred decomposition, or whether some other entropic quantity might be a better choice.
So we see that there remain many unanswered questions. In addition to these, there are difficulties with the minimal entropy proposal. For one, the product decomposition that minimizes the IU entropy may fail to be unique for certain dynamical state vectors. It may be that further technical work will show that this is not a problem after all. For instance, dynamical state vectors for which the preferred decomposition is non-unique for a finite interval of time may constitute a set of measure zero. Another difficulty is that the faithfulness criterion explicitly fails to be satisfied in measurements of variables whose eigenstates are entangled with respect to the distinguished factorization. Given that not all Hermitian operators necessarily correspond to variables that can be measured, it may happen that with a suitable choice of distinguished factorization, the measurements for which the faithfulness criterion fails to be satisfied are precisely those which are impossible to implement. On the other hand, it may be that this problem cannot be avoided within the minimal entropy proposal, but can be avoided if some other choice of preferred decomposition is made. As a third possibility, one might find that the faithfulness criterion for variables with entangled eigenstates, cannot be satisfied by any interpretation within the framework we have set out. Justifying any one of these answers would certainly be an interesting result, and motivates further investigation of these issues.
The use of a preferred decomposition, sometimes called an ‘interpretation basis’, has been viewed by some as necessary within interpretive strategies distinct from modal interpretations. This has been suggested by Deutsch in the context of the many-worlds interpretation and by Kent and McElwaine in the context of consistent histories. A preferred decomposition might also be useful in nonlinear modifications of quantum mechanics. Thus, the preferred decomposition of the minimal entropy proposal may well be of relevance to such interpretive strategies as well. In any event, a mathematically precise proposal, even though not without problems, can be useful in stimulating progress on interpretive issues, as is evidenced by the recent profusion of work on modal interpretations. We hope that the minimal entropy proposal will not be an exception in this respect.
## VII Acknowledgments
We wish to thank Rob Clifton for helpful comments on a draft of this paper. This work was supported by the National Sciences and Engineering Research Council of Canada.
## VIII Appendix: Proof of theorem 3.
It will be assumed throughout that the distinguished factorization is bi-partite, and the two factor spaces are denoted by $`^A`$ and $`^B.`$ All references to product decompositions are to be understood as product decompositions with respect to this factorization. We say that a product decomposition $`\{(c_k,|\chi _k^A|\varphi _k^B)\}_{k=1}^m`$ is $`A`$-orthogonal $`(B`$-orthogonal) if the set of vectors $`\{|\chi _k^A\}_{k=1}^m`$ $`\left(\{|\varphi _k^B\}_{k=1}^m\right)`$is orthogonal. A bi-orthogonal decomposition (also called a Schmidt decomposition) is one that is both $`A`$-orthogonal and $`B`$-orthogonal. We shall make use of several well-known properties of bi-orthogonal decompositions, an exposition of which can be found in Ref. . Finally, we remind the reader that $`S_{|\psi }(D)`$ denotes the IU entropy of $`|\psi `$ for the decomposition $`D,`$ which is defined by Eq.(28).
Theorem 3 follows from two lemmas:
For any vector $`|\psi ,`$ if $`D`$ is an arbitrary product decomposition of $`|\psi ,`$ then there always exists an $`A`$-orthogonal decomposition of $`|\psi ,`$ $`D_{\text{A-orth}}`$, such that
$`S_{|\psi }(D_{\text{A-orth}})S_{|\psi }(D).`$
For any vector $`|\psi ,`$ if $`D_{A\text{-orth}}`$ is any $`A`$-orthogonal decomposition of $`|\psi ,`$ and $`D_{\text{bi-orth}}`$ is any bi-orthogonal decomposition of $`|\psi ,`$ then
$`S_{|\psi }(D_{\text{bi-orth}})S_{|\psi }(D_{A\text{orth}}).`$
Together these imply that for any vector $`|\psi ,`$ if $`D`$ is an arbitrary product decomposition of $`|\psi `$, and $`D_{\text{bi-orth}}`$ is any bi-orthogonal decomposition of $`|\psi ,`$ then
$`S_{|\psi }(D_{\text{bi-orth}})S_{|\psi }(D),`$
which is simply theorem 3.
The task at hand is therefore to prove lemmas A.1 and A.2. We begin by reviewing a partial ordering relation among probability distributions, namely that of majorization, which has recently seen application in the study of entanglement purification. Suppose $`𝐩(p_1,p_2,\mathrm{},p_m)`$ and $`𝐪(q_1,q_2,\mathrm{},q_m)`$ are two $`m`$-element probability distributions. By definition, $`𝐩`$ majorizes $`𝐪`$ if for every $`l`$ in the range $`\{1,..,m\},`$
$`{\displaystyle \underset{k=1}{\overset{l}{}}}p_k^{}{\displaystyle \underset{k=1}{\overset{l}{}}}q_k^{},`$
where $`p_k^{}`$ indicates the $`k`$th largest element of $`𝐩,`$ so that $`p_1^{}p_2^{}\mathrm{}p_m^{}.`$
The notion of majorization is important in the present investigation because of the following well-known result (Theorem II.3.1 of Ref. ): The following two conditions are equivalent
$`(i)\text{ }𝐩\text{ majorizes }𝐪.`$
$`(ii)\text{ }{\displaystyle \underset{k=1}{\overset{m}{}}}f(p_k)`$ $``$ $`{\displaystyle \underset{k=1}{\overset{m}{}}}f(q_k)\text{ for all concave functions }f.`$
Since $`x\mathrm{log}x`$ is a concave function of $`x`$, it follows that $`H(𝐩)H(𝐪)`$ if and only if $`𝐩`$ majorizes $`𝐪,`$ where $`H(𝐩)`$ is the Shannon entropy of a probability distribution $`𝐩`$, defined in Eq.(27). Now consider two decompositions of a state vector, $`D=\{(c_k,|\varphi _k\}_{k=1}^m`$ and $`D^{}=\{(c_k^{},|\varphi _k^{}\}_{k=1}^m^{}.`$ Although these may have different cardinalities, they can be associated with probability distributions of equal cardinality by simply adding zeroes. Specifically, if $`mm^{},`$ then $`D`$ is associated with the distribution $`p_k=\left|c_k\right|^2`$ for $`k\{1,\mathrm{},m\}`$ and $`D^{}`$ is associated with the distribution $`q_k=`$ $`\left|c_k^{}\right|^2`$ for $`k\{1,\mathrm{},m^{}\}`$ and $`q_k=0`$ for $`k\{m^{}+1,\mathrm{},m\}.`$ Since the IU entropy of $`|\psi `$ for the decomposition $`D(D^{})`$ is simply the Shannon entropy of $`𝐩`$($`𝐪)`$, it follows that $`S_{|\psi }(D)S_{|\psi }(D^{})`$ if $`𝐩`$ majorizes $`𝐪.`$
In order to facilitate the proof of lemma A.1, we set out two minor lemmas.
Consider two probability distributions $`𝐩(p_1,p_2,\mathrm{},p_m)`$ and $`𝐪(q_1,q_2,\mathrm{},q_m).`$ If for every $`l`$ in the range $`\{1,\mathrm{},m\},`$
$`{\displaystyle \underset{k=1}{\overset{l}{}}}p__k{\displaystyle \underset{k=1}{\overset{l}{}}}q_k^{},`$
then $`𝐩`$ majorizes $`𝐪.`$
Proof. This result follows from the definition of majorization and the fact that
$`{\displaystyle \underset{k=1}{\overset{l}{}}}p_k^{}{\displaystyle \underset{k=1}{\overset{l}{}}}p_k`$
for every $`l`$ in the range $`\{1,\mathrm{},m\}.`$ This inequality is obviously true since the $`l`$-element subset of $`𝐩`$ with the largest sum must be the subset containing the $`l`$ largest elements of $`𝐩.`$ QED.
For the second minor lemma, we make use of some notational conventions introduced in the text: $`P_𝒮`$ denotes the projector onto the subspace $`𝒮,`$ and ‘$`P_𝒮<P_𝒮^{}`$’ denotes that $`𝒮`$ is a proper subspace of $`𝒮^{}.`$
If $`P_𝒮P_𝒮^{}`$ then $`\psi \left|P_𝒮\right|\psi \psi \left|P_𝒮^{}\right|\psi .`$
Proof. If $`P_𝒮=P_𝒮^{},`$ then the inequality is saturated. Otherwise, $`P_𝒮<P_𝒮^{},`$ and there exists a projector $`P_𝒯`$ such that $`P_𝒮+P_𝒯=P_𝒮^{}.`$ The desired inequality follows from the positivity of $`\psi \left|P_𝒯\right|\psi .`$ QED.
We are now in a position to prove lemma A.1.
Proof of lemma A.1. An arbitrary product decomposition has the form $`D=\{(d_k,|\varphi _k^A|\chi _k^B)\}_{k=1}^m,`$ where the lists of vectors $`\{|\varphi _k^A\}_{k=1}^m`$ and $`\{|\chi _k^A\}_{k=1}^m`$ are not necessarily orthogonal nor even linearly independent (although the list of vectors $`\{|\varphi _k^A|\chi _k^B\}_{k=1}^m`$ is orthogonal). The decomposition $`D`$ defines an $`m`$-element probability distribution $`𝐪=\{q_1,q_2,\mathrm{},q_m\},`$where $`q_k\left|d_k\right|^2.`$ As before, let $`q_k^{}`$ denote the $`k`$th largest element of $`𝐪,`$ and let $`|\varphi _k^A`$ and $`|\chi _k^A`$ denote the vectors associated with $`q_k^{}.`$
Now, identify every vector in the list $`\{|\varphi _k^A\}_{k=1}^m`$ that cannot be obtained as a linear combination of vectors with lower indices from this list. Suppose there a number $`m^{}`$ of such vectors, corresponding to a particular subset $`S`$ of the indices $`\{1,2,\mathrm{},m\},`$ so that the set of vectors is denoted by $`\{|\varphi _k^A\}_{kS}`$. By definition, this is a linearly independent set. The remaining vectors are denoted by $`\{|\varphi _k^A\}_{k\overline{S}},`$ where $`\overline{S}`$ is the set of indices that remain after removing the elements of $`S`$ from $`\{1,2,\mathrm{},m\}.`$ Obviously the elements of $`\{|\varphi _k^A\}_{k\overline{S}}`$ can all be written as linear combinations of the elements of $`\{|\varphi _k^{}^A\}_{k^{}S,\text{ }k^{}<k}.`$ Finally, for future reference, we define $`g(k)`$ as the number of indices $`k^{}`$ in $`S`$ such that $`k^{}k.`$ It is clear from the definition of $`S`$ that $`g(k)k.`$
Let $`\{|\mu _j^A\}_{j=1}^m^{}`$ be the ordered set of orthogonal vectors that are obtained by applying the Gram-Schmidt orthogonalization procedure to $`\{|\varphi _k^A\}_{kS},`$ in order of ascending $`k.`$ This new set yields an $`A`$-orthogonal decomposition of $`|\psi ,`$ $`D_{\text{A-orth}}=\{(c_j,|\mu _j^A|\nu _j^B)\}_{j=1}^m^{},`$ where $`|\nu _j^B=\mu _j^A|\psi /c_j`$ and $`c_j=\left|\mu _j^A|\psi \right|.`$ It also defines an $`m`$-element probability distribution $`𝐩=(p_1,p_2,\mathrm{},p_m)`$ where $`p_j\left|c_j\right|^2`$ for $`j`$ in the range $`\{1,\mathrm{},m^{}\},`$ and $`p_j0`$ for $`j`$ in the range $`\{m^{}+1,\mathrm{},m\}.`$
Let $`P_\varphi ^A`$ denote the projector onto the ray spanned by $`|\varphi ^A`$ and for convenience define $`P_{\mu _j}^AP_{\text{null}}`$ for $`j`$ in the range $`\{m^{},\mathrm{},m\}.`$ The nature of the Gram-Schmidt orthogonalization procedure ensures that for every $`kS,`$ $`|\varphi _k^A=_{j=1}^{g(k)}f_j|\mu _j^A`$ for some set of complex amplitudes {$`f_j\}_{j=1}^{g(k)}`$. Thus, $`(_{j=1}^{g(k)}P_{\mu _j}^A)|\varphi _k^A=|\varphi _k^A,`$ or equivalently, $`_{j=1}^{g(k)}P_{\mu _j}^AP_{\varphi _k}^A.`$ Since $`g(k)k,`$ this is trivially extended to $`_{j=1}^kP_{\mu _j}^AP_{\varphi _k}^A.`$ Moreover, if $`k\overline{S},`$ then $`|\varphi _k^A`$ can be written as a linear combination of the elements of $`\{|\varphi _k^{}^A\}_{k^{}S,\text{ }k^{}<k},`$ so that $`|\varphi _k^A=_{j=1}^{g((h(k))}\overline{f}_j|\mu _j^A`$ for some set of complex amplitudes {$`\overline{f}_j\}_{j=1}^{g(h(k))}`$, where $`h(k)=\mathrm{max}_{k^{}S,k^{}<k}k^{}.`$ It follows that $`_{j=1}^{g((h(k))}P_{\mu _j}^AP_{\varphi _k}^A`$ for all $`k\overline{S}`$. Since $`g(h(k))<k,`$ this is trivially extended to $`_{j=1}^kP_{\mu _j}^AP_{\varphi _k}^A.`$ It follows therefore that for every $`k`$ in the range $`\{1,\mathrm{},m\}`$ we have $`_{j=1}^kP_{\mu _j}^AP_{\varphi _k}^A.`$ Now, since $`I^BP_{\chi _l}^B,`$ we can infer that $`_{j=1}^kP_{\mu _j}^AI^BP_{\varphi _k}^A`$ $`P_{\chi _k}^B,`$ and by the orthogonality of the projectors in the set $`\{P_{\varphi _l}^AP_{\chi _l}^B\}_{l=1}^k,`$ we conclude that $`_{j=1}^kP_{\mu _j}^AI^B_{l=1}^kP_{\varphi l}^A`$ $`P_{\chi _l}^B`$ for every $`k`$ in the range $`\{1,\mathrm{},m\}.`$
Now we note that the probability distributions $`𝐪`$ and $`𝐩`$ are related to the projectors by $`q_l^{}=\psi \left|P_{\varphi _l}^AP_{\chi _l}^B\right|\psi `$ and $`p_j=\psi \left|P_{\mu _j}^AI^B\right|\psi .`$ From the inequality derived above together with lemma A.4, we find that $`_{j=1}^kp_j_{l=1}^kq_l^{}`$ for every $`k`$ in the range $`\{1,\mathrm{},m\}.`$ By lemma A.3, it follows that $`𝐩`$ majorizes $`𝐪.`$ QED.
Finally, we prove lemma A.2.
Proof of lemma A.2. An arbitrary $`A`$-orthogonal decomposition of $`|\psi `$ has the form $`D_{A\text{-orth}}=\{(c_j,|\mu _j^A|\nu _j^B)\}_{j=1}^m,`$ where the vectors $`\{|\mu _j^A\}_{j=1}^m`$ are orthogonal, but $`\{|\nu _j^B\}_{j=1}^m`$ need not be orthogonal nor even linearly independent. A bi-orthogonal decomposition of $`|\psi `$ has the form $`D_{\text{bi-orth}}=\{(\stackrel{~}{c}_j,|\stackrel{~}{\mu }_j^A|\stackrel{~}{\nu }_j^B)\}_{j=1}^{\stackrel{~}{m}},`$ where both the vectors $`\{|\stackrel{~}{\mu }_j^A\}_{j=1}^{\stackrel{~}{m}}`$ and $`\{|\stackrel{~}{\nu }_j^A\}_{j=1}^{\stackrel{~}{m}}`$ form orthogonal sets. The probability distributions associated with each decomposition are $`(\left|c_1\right|^2,\left|c_2\right|^2,\mathrm{},\left|c_m\right|^2)`$ and $`(\left|\stackrel{~}{c}_1\right|^2,\left|\stackrel{~}{c}_2\right|^2,\mathrm{},\left|\stackrel{~}{c}_{\stackrel{~}{m}}\right|^2)`$ respectively (even if there is more than one bi-orthogonal decomposition for a particular state vector, these do not differ in their coefficients). For ease of comparison of these distributions, we add zeroes until the number of elements in each is equal to the dimensionality, $`d,`$ of the Hilbert space $`^A`$. Denote the resulting distributions by $`𝐩`$ and $`\stackrel{~}{𝐩}`$ respectively. We establish that $`S_{|\psi }(D_{\text{bi-orth}})S_{|\psi }(D_{A\text{orth}})`$ by showing that $`\stackrel{~}{𝐩}`$ majorizes $`𝐩.`$
To begin, we express the probabilities as expectation values of projectors. We introduce an arbitrary orthogonal set of vectors $`\{|\mu _j^A\}_{j=m+1}^d`$ which together with $`\{|\mu _j^A\}_{j=1}^m`$ form an orthogonal basis for the Hilbert space $`^A`$, and similarly for $`\{|\stackrel{~}{\mu }_j^A\}_{j=1}^{\stackrel{~}{m}}.`$ Then, we have for all $`j`$ in the range $`\{1,\mathrm{},d\},`$
$`p_j`$ $`=`$ $`Tr_A(\rho ^AP_{\mu _j}^A),\text{ and}`$
$`\stackrel{~}{p}_j`$ $`=`$ $`Tr_A(\rho ^AP_{\stackrel{~}{\mu }_j}^A).`$
Let the unitary operator that transforms the elements of $`\{|\stackrel{~}{\mu }_j^A\}_{j=1}^d`$ to the elements of $`\{|\mu _j^A\}_{j=1}^d`$ be denoted by $`U^A,`$ so that
$`|\mu _j^A=U^A|\stackrel{~}{\mu }_j^A.`$
It follows that
$`p_j`$ $`=`$ $`Tr_A(\rho ^AU^AP_{\stackrel{~}{\mu }_j}^AU^A)`$
$`=`$ $`Tr_A(U^A\rho ^AU^AP_{\stackrel{~}{\mu }_j}^A),`$
where in the last step we have used the cyclic property of the trace. What distinguishes the bi-orthogonal decomposition from other $`A`$-orthogonal decompositions is that the projectors $`\{P_{\stackrel{~}{\mu }_k}^A\}_{k=1}^d`$ diagonalize $`\rho ^A`$,
$`\rho ^A={\displaystyle \underset{k=1}{\overset{d}{}}}\stackrel{~}{p}_kP_{\stackrel{~}{\mu }_k}^A.`$
Plugging this form of $`\rho ^A`$ into the expression for $`p_j,`$ we obtain
$`p_j={\displaystyle \underset{k=1}{\overset{d}{}}}\left|U_{jk}^A\right|^2\stackrel{~}{p}_k,`$
where $`U_{jk}^A=\stackrel{~}{\mu }_j\left|U^A\right|\stackrel{~}{\mu }_k.`$ By the unitarity of $`U^A`$, we find that $`_j\left|U_{jk}^A\right|^2=\stackrel{~}{\mu }_k\left|U^AU^A\right|\stackrel{~}{\mu }_k=1,`$ and $`_k\left|U_{jk}^A\right|^2=\stackrel{~}{\mu }_j\left|U^AU^A\right|\stackrel{~}{\mu }_j=1.`$ Thus, the transition matrix between the probability distributions $`\stackrel{~}{𝐩}`$ and $`𝐩`$ is doubly stochastic, from which it follows by a well-known result (theorem II.1.9 of Ref. ) that $`\stackrel{~}{𝐩}`$ majorizes $`𝐩.`$ QED.
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# On breaking the age-metallicity degeneracy in early-type galaxies: Outflows versus Star Formation Efficiency
## 1 Introduction
The process of galaxy formation and evolution can be explored in two complementary ways: A “forwards” approach takes into account the physics underlying the most basic processes of structure evolution and star formation and — after finding a suitable set of initial conditions — evolves the system forward so that the final output is compared with observations. This is the philosophy behind N-body simulations or semi-analytic modelling (e.g. Baugh et al. 1998; Kauffmann & Charlot 1998). On the other hand, a “backwards” approach simplifies the physics behind galaxy formation and evolution to a phenomenological problem comprising a reduced set of parameters, using local observations as constraints. Then the system is evolved backwards so that the predictions for a given set of parameters are compared with observations at moderate-to-high redshifts (e.g. Bouwens, Broadhurst & Silk 1998a,b). This “brute force” method allows one to search a reasonable volume of parameter space, throwing light on processes, such as star formation, which are otherwise too complicated to tackle from basic principles.
Cluster early-type galaxies are ideal candidates for a comparison between models and observations. Significant samples of these galaxies can be found over a large redshift range. Furthermore, the tight observed correlations such as the colour-magnitude relation or the fundamental plane can be used as powerful constraints on a phenomenological backwards approach. The list of observed clusters at moderate to high redshift is quite large (e.g. Dressler et al. 1999 ; Stanford, Eisenhardt & Dickinson 1998; Oke, Postman & Lubin 1998; Van Dokkum 1999) and ever increasing (e.g. Yee et al. 1999). Clusters observed at redshifts $`z11.5`$ yield valuable information about the epoch of star formation, pushing it to very high redshift ($`z_F3`$).
Unfortunately, the direct spectrophotometric determination of the star formation history is hampered by the age-metallicity degeneracy (Worthey 1994; Ferreras, Charlot & Silk 1999) which allows the variations of most of the spectrophotometric observables to be explained either by a range of ages or metallicities. Broadband photometry is strongly dependent both on age and metallicity but even spectral indices targeting single lines such as Balmer absorption or magnesium abundance can change both with age and metallicity. Not surprisingly, a similar age estimation technique based on spectral indices applied to similar sets of elliptical galaxies yields contradictory results: Kuntschner (2000) and Kuntscher & Davies (1998) find coeval stellar populations in Fornax cluster ellipticals so that the colour range is explained by a metallicity sequence. On the other hand, the sample of field and group ellipticals observed by González (1993) and further analysed by Trager et al. (2000) presents a relatively large spread in ages. Hence, the issue of the stellar age distribution in galaxies still requires the aid of modelling. We will show that incorporating chemical enrichment allows one to potentially solve the age-metallicity degeneracy problem. The next three sections describe our chemical enrichment model and the meaning of the reduced set of parameters used to trace the star formation history in cluster ellipticals. §5 deals with the comparison of predicted and observed mass-to-light ratios and its use at high redshift to discriminate between a mass sequence driven by age or metallicity. Finally in §6 we discuss the predictions and list the main conclusions.
## 2 Model description
The basic mechanisms describing chemical enrichment in galaxies can be reduced to infall of gas, metal-rich outflows triggered by supernovae, and a star formation prescription. Ferreras & Silk (2000, hereafter FS00) described one such simple model for an exponentially decaying infall rate of primordial gas in terms of five parameters: flow rate, timescale and delay of infall, ejected fraction in outflows, and star formation efficiency. Monotonically decreasing star formation rates (SFRs) always encounter the so-called G-dwarf problem. In FS00 the problem was avoided by assuming the metallicity effect was predominant with regard to the spectrophotometric output, thereby enabling us to convolve simple stellar populations with different metallicities given by the chemical enrichment equations but with an average common age.
The G-dwarf problem arises when comparing the predictions of closed box models with the observed paucity of low-metallicity low-mass stars in the halo of our galaxy (Van den Bergh 1962). However, this problem extends to all galaxies and morphologies (Worthey et al. 1996) since any model assuming a monotonically decreasing infall rate overproduces low-metallicity, low-mass stars. These models can be readily ruled out for bright cluster early-type galaxies because the convolution in age and metallicity of simple stellar populations with the chemical enrichment tracks predicted by these models yield $`UV`$ or $`VK`$ colours which are significantly bluer than the observed values, $`UV1.6`$, $`VK3.3`$ for Coma or Virgo (Bower, Lucey & Ellis 1992). The solution to this problem involves either pre-enrichment, i.e. shifting the zero-point of the metal abundance by assuming non-primordial infall (e.g. Sansom & Proctor 1998), or imposing an initial stage of moderate star formation that enriches the interstellar medium, followed by a second stage where most of the stars (with non-primordial abundances) are formed. Several authors (Elbaz, Arnaud & Vangioni-Flam 1995; Chiappini, Matteucci & Gratton 1997) have suggested two-stage processes to solve this problem, as well as a top-heavy initial mass function for the first star forming stage that overproduces high mass stars for a prompt enrichment of the interstellar medium, motivated by the fact that starbursts can be best explained by top-heavy mass functions (Charlot et al. 1993).
In this paper we present a model which avoids these assumptions and solves the G-dwarf problem by considering infall of primordial gas with a rate given by a gaussian distribution. This infall combined with a linear Schmidt law results in a gaussian profile for the star formation rate which is a suitable way of parametrizing a strongly peaked starburst. Furthermore, the extension of this profile to disk galaxies, and in particular to our own galaxy, accounts for the observed metallicity distribution of stars (Rocha-Pinto & Maciel 1996, Wyse & Gilmore 1995): the first stars are formed in a stage with a low SFR, but their contribution to the enrichment of the IGM, along with an increasing SFR with time yields a large population of stars with metallicities close to solar. Once the IGM has reached a very large metallicity — after the peak in the SFR — a low formation rate keeps the tally of stars with a high metal content as low as observed. In this scenario, the difference between early-type and late-type galaxies would amount to an extra component in the SFR for late-type systems, which continue to form stars at a slow rate (Ferreras & Silk, in preparation). Hence, the work in this paper uses the same chemical enrichment equations (described in FS00) based on the formalism described by Tinsley (1980), with just four parameters describing the star formation process, namely:
* Star Formation Efficiency ($`C_{\mathrm{eff}}`$): The star formation rate is modelled by a linear Schmidt law, whose proportionality constant determines the formation efficiency. This parameter is a phenomenological approximation to the complex physics underlying star formation. Its inverse represents the timescale in Gyr over which the SFR is extended for a sharply peaked infall of gas. In the Instantaneous Recycling Approximation (IRA), the SFR ($`\psi (t)`$) and the infall rate ($`f(t)`$) are related by:
$$\psi (t)=C_{\mathrm{eff}}_0^{\mathrm{}}𝑑sf(ts)e^{s/\tau _g},$$
(1)
$$\tau _g=\frac{1}{C_{\mathrm{eff}}\left[1(1B_{\mathrm{out}})R\right]};$$
(2)
$`B_{\mathrm{out}}`$ is the ejected fraction of gas in outflows — defined below — and $`R`$ is the returned fraction. Low values of $`C_{\mathrm{eff}}`$ extend the process of star formation, generating a larger age spread of the stellar populations.
* Gas Outflows ($`B_{\mathrm{out}}`$): Part of the metal-enriched gas in the IGM is heated by supernovae and ejected from the galaxy, decreasing the effective yield. This will be modulated by the gravitational potential well of the galaxy, so that bright and massive galaxies have lower ejected fractions and so, higher metal abundances (Larson 1974, Arimoto & Yoshii 1987). $`B_{\mathrm{out}}`$ represents the fraction of gas being ejected at each timestep.
* Formation redshift ($`z_F`$): This parameter describes the epoch — $`t(z_F)`$ — at which the infall rate (or roughly the SFR) was maximum.
* Infall timescale ($`\tau _f`$): Determines the duration of infall, so that the infall rate can be written:
$$f(t,\tau _f,z_F)=\frac{1}{\tau _f\sqrt{2\pi }}\mathrm{exp}\left[\frac{\left(tt(z_F)\right)^2}{2\tau _f^2}\right]$$
(3)
For a given set of these four parameters ($`C_{\mathrm{eff}}`$,$`B_{\mathrm{out}}`$, $`z_F`$,$`\tau _f`$) the equations give a chemical enrichment track which is used to convolve the simple stellar populations from the models of Bruzual & Charlot (2000) both in age and metallicity, to obtain a spectral energy distribution for a given redshift. Motivated by recent results of the angular power spectrum of the Cosmic Microwave Background (Melchiorri et al. 1999), we assume a flat cosmology with a cosmological constant ($`\mathrm{\Omega }_m=0.3`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$, H$`{}_{0}{}^{}=60`$ km s<sup>-1</sup>Mpc<sup>-1</sup>), although figure 1 shows the result for an open cosmology ($`\mathrm{\Omega }_\mathrm{\Lambda }=0`$) in the range of infall parameters ($`\tau _f`$ and $`z_F`$). The fiducial Initial Mass Function (IMF) chosen is a hybrid one between the Scalo (1986) and the Salpeter (1955) IMF with mass cutoffs $`0.1<M/M_{}<60`$. We have adopted the behaviour of the Scalo IMF for low stellar masses ($`M<2M_{}`$) and a Salpeter IMF for the high mass end. The shallow slope of the former for low masses better explains the observations, whereas the steeper slope at high masses of the Salpeter IMF seem to explain better the stellar populations in starbursts. For comparison purposes we show the difference between a Salpeter IMF and our mass function when calculating the correlation between mass-to-light ratio and mass (§5).
Broadband colours are computed from the spectral energy distribution obtained for a given star formation history described by these four parameters ($`C_{\mathrm{eff}}`$,$`B_{\mathrm{out}}`$,$`z_F`$,$`\tau _f`$). The colour-magnitude relation observed in local clusters is used as a constraint, so that a one-to-one mapping is assumed between colour and absolute luminosity. We have used the observations of Coma early-type galaxies (ellipticals and lenticulars) by Bower et al. (1992), whose colour-magnitude relation has a scatter ($`\pm 0.05`$ mag) as small as the uncertainties expected from population synthesis models (Charlot, Worthey & Bressan 1996). Hence, for a given $`UV`$ colour obtained from the spectral energy distribution we infer an absolute luminosity:
$$M_V=\frac{(UV)+0.3830}{0.0871}.$$
(4)
This correlation was obtained by a robust two-stage linear fit technique which computes a first estimate of slope and zero point using absolute deviations (e.g. Press et al. 1992) and then applies a least-squares fit to the resulting sample after culling points that deviate more than a given threshold. This method prevents outliers from contributing significantly to the final slope and zero point. In this case the number of outliers was 4 out of 36 galaxies. Once the $`V`$-band absolute luminosity is found, the total stellar mass is obtained using the predicted mass-to-light ratio.
## 3 Constraints from bright cluster galaxies
The colour of the brightest galaxies in clusters at several redshifts is one of the strongest and most robust constraint that can be imposed on the model parameters. Rest frame $`UV`$ colour is very sensitive to age (as well as metallicity). The red $`UV`$ colours of the brightest early-type systems observed in Coma ($`UV1.6`$) need a dominant population of old stars with a metallicity around solar or higher. Hence, large infall timescales ($`\tau _f`$) or low formation redshifts are readily ruled out since they predict bluer $`UV`$ colours. Furthermore, observations of clusters at moderate redshifts (e.g. Ellis et al. 1997; Stanford, Eisenhardt & Dickinson 1998; Dressler et al. 1999) have found a colour-magnitude relation consistent with passive evolution, imposing again an old age and high metal content. This motivates the need to use a low ejected fraction ($`B_{\mathrm{out}}0`$) as well as a high star formation efficiency ($`C_{\mathrm{eff}}5`$) in order to reproduce the spectrophotometric properties of these galaxies. Figure 1 shows the constraint on infall parameters for the brightest galaxies using a grid of models with $`B_{\mathrm{out}}=0`$ and $`C_{\mathrm{eff}}=10`$. The $`z0`$ constraint imposes a colour$``$colour correlation between $`UV`$ and $`VK`$ as observed in Coma and Virgo (Bower et al. 1992) with a scatter of $`\pm 0.05`$ mag. Notice that as long as the bursting stage is short enough ($`\tau _f1`$ Gyr), formation redshifts as low as $`z_F1`$ can be accomodated with the photometric data (Ferreras et al. 1999). In fact, a simple stellar population (which is the extreme case of the model described in this paper as $`\tau _f0`$) gives a $`z=0`$ color of $`UV1.5`$ for a stellar population with solar metallicity and formation redshift $`z_F=1`$. This corresponds to a bright ($`M_V21.6`$) cluster galaxy. The $`z0.5`$ constraint imposes a colour F555W$``$F814W (which maps into restframe $`UV`$) redder than $`2.4`$ for the brightest cluster galaxies, as observed in Cl0016+16 ($`z=0.55`$, Ellis et al. 1997 ). The high redshift constraint requires $`RK5.7`$ as observed in cluster RX J0848.9+4452 ($`z=1.26`$, Rosati et al. 1999). These constraints rule out models with infall timescales $`\tau _f1.5`$ Gyr, which reduces the main star formation episode to less than $`3`$ Gyr. The formation redshift is less restricted, allowing the possibility of $`z_F3`$, or even $`z_F2`$ with a flat cosmology with a cosmological constant ($`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$).
The constraint on the allowed volume of parameter space with respect to infall and its subsequent process of star formation shown in figure 1 is rather robust as it takes into account an observable with a very small uncertainty, namely the broadband colours of the most luminous galaxies. Any significant process of star formation at later stages will readily translate into bluer colours which are ruled out by observations. The modelling of star formation rate as a gaussian profile allows the extension of single burst models to processes where several bursts of star formation take place, as in a “gaseous” merger scenario. The spectrophotometric predictions of a multi-burst scenario can be mimicked by a single gaussian infall rate with a longer time scale (Ferreras & Silk 2000b). Hence, if we consider RX J0848.9+4452 as a representative cluster which will evolve into something like Coma or Virgo, then any hierarchical merging scenario should yield gaseous mergers only at epochs earlier than $`z_F3`$. Merging stages at lower redshifts should not trigger any significant star formation, for example involving gas-poor progenitors and hot gas, as hinted at by observations of red merging galaxies in cluster MS$`105403`$ ($`z=0.83`$ Van Dokkum et al. 1999).
It is worth emphasizing at this point that the brightest galaxies we are dealing with in this model represent the brightest “non-peculiar” systems. Hence, the very brightest cluster galaxies (BCGs) such as cDs should be excluded. In a recent paper, Aragón-Salamanca, Baugh & Kauffmann (1998) analysed the evolution with redshift of a sample of BCGs in 25 clusters in the redshift range $`0<z<1`$. They found no evolution, or even a fading of the absolute $`K`$-band luminosity which can be explained by a significant change (a factor of 2) in the stellar mass of these galaxies between redshifts $`z=1`$ and $`0`$. In the formalism of our model, the reddest — thereby brightest — galaxies are obtained in the corner of parameter space associated with high star formation efficiency ($`C_{\mathrm{eff}}510`$). This implies the star formation rate will be approximately given by the infall rate $`\psi (t)f(t)`$. For the allowed values of infall parameters ($`\tau _f`$, $`z_F`$) shown in figure 1, the change in stellar mass between redshifts $`z=1`$ and $`0`$ is very small, and so this model would not be compatible with BCGs. This type of galaxy requires a strong merger rate at late times ($`z1`$). Furthermore, the mergers should undergo no star formation, involving progenitors with either very little gas or hot gas. On the other hand, the evolution of the bulk of cluster early-type systems was explored in a sample of 38 clusters in a similar redshift range ($`0.1<z<1`$) by analysing the near-infrared luminosity function (De Propris et al. 1999). A significant positive luminosity evolution was found for the Schechter parameter $`K_{}`$ which traces the luminosity of the brightest galaxies. The evolution of $`K_{}`$ is found to be compatible with the passive evolution of a simple stellar population, which is consistent with the models presented in this paper for the range of infall parameters shown in figure 1.
## 4 Exploring parameter space
Once we have fixed a region in parameter space for the brightest “non-peculiar” cluster galaxies, we can extend the model to all early-type systems by exploring the parameters describing star formation efficiency and fraction of ejected gas in outflows. Figure 2 shows a contour of stellar masses for two realizations of infall ($`z_F=5`$, $`\tau _f=0.5`$ Gyr, dotted lines) and ($`z_F=5`$, $`\tau _f=1`$ Gyr, solid lines). The only constraint imposed is the colour-magnitude relation in Coma in order to relate colour — obtained from the star formation history described by a chosen set of parameters ($`B_{\mathrm{out}}`$, $`C_{\mathrm{eff}}`$, $`\tau _f`$, $`z_F`$) — and luminosity, or stellar mass. The contours change slightly with infall parameters but the qualitative behaviour is unchanged. We need both a high star formation efficiency and a low ejected fraction in outflows to account for the most massive galaxies. Lowering the SF efficiency will generate a significant spread in stellar ages that will yield colours that are too blue compared with the observations. On the other hand, a higher ejected fraction will decrease the effective yield, so that the average metallicity will be lower, blueing the predicted colours with respect to the observations. Hence, regardless of the choice of parameters, we find the brightest cluster galaxies to lie in a well-defined corner of the 4-dimensional parameter space explored in this paper. Lower mass galaxies are harder to identify because of a degeneracy between outflow rates and star formation efficiency. The blueing of galaxies as we progress down in luminosity along the colour-magnitude relation can be explained either by decreasing the average metallicity of the stellar populations (i.e. a range of ejected fractions, $`B_{\mathrm{out}}`$) or by increasing the fraction of younger stars (i.e. a range of SF efficiencies, $`C_{\mathrm{eff}}`$). The real picture will likely involve a range both in $`B_{\mathrm{out}}`$ and $`C_{\mathrm{eff}}`$. At low redshift it is hard to disentangle the contribution from these two parameters, mainly caused by the age-metallicity degeneracy (Worthey 1994). However, these two alternative mechanisms will predict different evolutions with lookback time: as we go to higher redshift, the sequence driven by SF efficiency (hereafter $`C_{\mathrm{eff}}`$ sequence) will show stronger evolutionary scars because of the larger scatter in age, whereas the sequence driven by outflows (hereafter $`B_{\mathrm{out}}`$ sequence) will be unaffected until the observations go back to the epoch of star formation.
Figure 3 shows the predicted distributions of stellar ages (left) and metallicities (right). The average (top) and standard deviation (bottom) of these distributions are shown for three sequences: $`B_{\mathrm{out}}`$,$`C_{\mathrm{eff}}`$ and $`B_{\mathrm{out}}+C_{\mathrm{eff}}`$ driven by outflows, star formation efficiency and a mixture of both, respectively. All of these models have $`\tau _f=1`$ Gyr, $`z_F=5`$; the $`B_{\mathrm{out}}`$ sequence fixes $`C_{\mathrm{eff}}=5`$; the $`C_{\mathrm{eff}}`$ sequence fixes $`B_{\mathrm{out}}=0`$; and the mixed $`B_{\mathrm{out}}+C_{\mathrm{eff}}`$ sequence imposes a correlation between these two parameters. We expect the star formation efficiency to be higher in more massive galaxies (where the ejected fraction should be lower), hence we have chosen: $`\mathrm{log}C_{\mathrm{eff}}=13B_{\mathrm{out}}`$, as an illustrative example. An outflow-driven sequence generates a pure metallicity sequence where the stellar populations of all galaxies are coeval. This is the model favored by the analysis of the colour-magnitude relation by Kodama et al. (1998) for a sample of 17 clusters in a wide redshift range ($`0.3<z<1.3`$). In a $`B_{\mathrm{out}}`$ sequence the colour-magnitude relation is due to a pure mass-metallicity correlation, so that the brightest galaxies have higher metal abundances. On the other hand, a $`C_{\mathrm{eff}}`$ or a $`B_{\mathrm{out}}+C_{\mathrm{eff}}`$ model represents mixed age+metallicity sequences as it generates a significant age spread of the stellar populations as well as a metallicity range. A pure age sequence explored in other papers (e.g. Kodama & Arimoto 1997) is ruled out when including chemical enrichment for any reasonable assumption about the star formation history of the galaxy. The low metallicity expected in the fainter galaxies can therefore be explained either by a lower effective yield caused by gas outflows, or by a low star formation efficiency that slows down the process of enrichment. Nagashima & Gouda (1999) mention the effect of the UV background radiation as a possible alternative to supernovae winds to suppress chemical enrichment. However, in this scenario it seems hard to motivate the correlation found with respect to luminosity (e.g. a stronger feedback should be expected from the UV background for the fainter galaxies). Figure 3 shows that a $`C_{\mathrm{eff}}`$ sequence implies a $`3`$ Gyr age difference between the brightest systems and $`M_V16`$ galaxies with a larger age spread for the fainter galaxies, as expected for a system with a lower star formation efficiency. This model would thus be consistent with the smaller age and metallicity scatter found in Coma ellipticals (Jørgensen 1999). This age difference is negligible for local clusters, since the predicted youngest stars are still too old ($`8`$ Gyr) to be detected spectrophotometrically. However, predictions between a coeval sequence and a $`C_{\mathrm{eff}}`$ sequence at lookback times $`8`$ Gyr (i.e. $`z1`$) differ significantly as shown below (§5). A $`B_{\mathrm{out}}`$ sequence predicts a small age and metallicity spread in faint ellipticals. In this case the observed photometric scaling relations are just driven by the average metallicity.
## 5 Mass-to-Light ratios: Stellar versus Observed
The evolution of the stellar mass-to-light ratio is a suitable probe for the analysis of the ages of the stellar populations in galaxies because of its weak dependence on metallicity (especially when measured in NIR passbands). However, observed $`M/L`$ ratios carry large uncertainties and systematic offsets as they are computed from central velocity dispersions ($`\sigma _0`$), sizes ($`r_e`$) and surface brightnesses ($`\mathrm{\Sigma }_e`$), requiring a prescription for the structure in order to infer masses and luminosities. From a simple dimensional analysis, we can write the mass and luminosity of a galaxy as:
$$M=\alpha \sigma _0^2r_e$$
(5)
$$L=\beta r_e^2\mathrm{\Sigma }_e$$
(6)
The proportionality constant $`\beta `$ only depends on the way the surface brightness is defined with respect to the effective size. However, $`\alpha `$ depends on the structure of the galaxy, which requires a dynamical model (e.g. King 1966; Bender, Burstein & Faber 1992), usually assumed to be invariant with regard to galaxy mass or size (i.e. homologous). Furthermore, a comparison between observations and model predictions depends on a prescription that relates stellar and total matter. Despite all these caveats, mass-to-light ratios are one of the best candidates for breaking the age-metallicity degeneracy.
Figure 4 shows the predicted $`M/L`$ ratios in two passbands: optical ($`V`$) and NIR ($`K`$) for a grid of $`B_{\mathrm{out}}`$ (thin) and $`C_{\mathrm{eff}}`$ (bold) sequences. A $`B_{\mathrm{out}}+C_{\mathrm{eff}}`$ mixed sequence would be represented by a trajectory between the grid spanned by these lines. The results are shown for three different infall timescales $`\tau _f=0.5,1,1.5`$ Gyr. The solid triangles are the observations of Coma cluster galaxies from Mobasher et al. (1999), whereas the dots represent ellipticals in 11 clusters from the sample of Pahre (1999). The slope of the correlation between mass-to-light ratios and masses is $`0.24`$ in the $`V`$ band and $`0.14`$ in $`K`$ band, whereas purely stellar $`M/L`$ ratios yield slopes around $`0.1`$ and $`0.0`$ in the $`V`$ and $`K`$ bands, respectively. This mismatch cannot be related to a different stellar population. Being a strongly age-dependent observable, the only way to generate the observed slopes in $`M/L`$ would require a stellar age spread with respect to galaxy mass whose restframe $`UV`$ colour would significantly disagree with the colours in moderate redshift cluster ellipticals. As a simple check to confirm this, we used simple stellar populations (SSP, i.e. no age or metallicity spread). If we relate the brightest galaxies ($`M_V22.5`$) to a SSP with an age of 12 Gyr and a super-solar metallicity $`Z=1.5Z_{}`$, the required age for the SSP for a $`M_V=19`$ galaxy which would explain the observed slope $`M/L_VM^{0.24}`$ should be $`t=4`$ Gyr ($`Z=Z_{}`$) or $`t=5`$ Gyr ($`Z=Z_{}/2`$). These two possibilities yield $`UV`$ colours of $`1.20`$ and $`1.07`$ respectively, so that the latter is barely consistent with the observed broadband photometry in nearby clusters (Bower et al. 1992). However, this hypothesis fails when considering moderate redshift clusters: for instance, cluster Cl0016+16 ($`z0.5`$) has a well-defined red-envelope in a large range of luminosities (Ellis et al. 1997), yet the lookback time ($`6.2`$ Gyr in our adopted cosmology) would be larger than the expected age for the stars in the fainter galaxies !
Hence, the only plausible way to solve this slope mismatch is by imposing a non-linear correlation between total mass (including dark matter) and stellar mass. A simple power law $`MM_{\mathrm{ST}}^{1.2}`$ (FS00) has been used in figure 4 to transform predicted stellar masses and mass-to-light ratios to total masses and $`M/L`$. The agreement is good in both bands, which means the observed slope difference in the fundamental plane between passbands can be explained by the stellar populations alone. This result agrees with the analysis of Jørgensen (1999) who refers to a higher fraction of dark matter in the most massive ellipticals in order to account for the mismatch between age estimates using $`M/L`$ ratios (which involves the total mass of the galaxies) or $`H\beta _G`$ indices (involving just the stellar component). However, one should bear in mind that there are alternative explanations for this slope mismatch, such as a systematic variation of the initial mass function or the breaking of the homology for the structural parameters of early-type systems (Jørgensen 1999; Pahre et al. 1998). Graham & Colless (1997) examined the effect of a non-homologous light profile in 26 early-type systems in the Virgo cluster and found no differences in the fundamental plane from estimates using a universal de Vaucouleurs profile.
Notice the scatter of the correlation is not the minimum one can obtain when computing the fundamental plane. The dependence of the fundamental plane on the observables (velocity dispersion, surface brightness and galaxy size) is not precisely the one inferred for the correlation between mass-to-light ratio and mass. This scatter has an observational component as well as an intrinsic contribution. The intrinsic part can be related to differences in either the stellar ages or the initial mass function of galaxies with the same mass. The former can be caused by slightly different formation redshifts or infall timescales. Notice the predictions of $`M/L`$ ratios for three different infall timescales in figure 4. A variation of 1 Gyr in infall timescale can account for changes in $`\mathrm{log}M/L`$ roughly of order $`0.1`$ dex. Furthermore, a non-universal IMF will add a variation around $`0.20.3`$ dex. The top-center panel displays the predictions for our fiducial “hybrid” IMF and a Salpeter IMF with the same mass cutoffs ($`0.1<M/M_{}<60`$). A different cutoff or slope at the high mass end will not change the result very much since these changes only have an effect on metallicity. However, changes in the low mass end of the IMF will significantly alter the mass-to-light ratio. Our “hybrid IMF” behaves like a Scalo function at low masses, whereas a Salpeter mass function produces a large fraction of low mass stars, thereby increasing the stellar $`M/L`$ ratio.
The different age distributions obtained for the $`B_{\mathrm{out}}`$ and $`C_{\mathrm{eff}}`$ sequences generate different predictions for the evolution of the slope of the correlation between stellar $`M/L`$ ratio and mass. We have shown above that a $`B_{\mathrm{out}}`$ sequence (driven by outflows) is equivalent to a metallicity sequence. Hence, a very small slope change is expected since mass-to-light ratios are not very sensitive to metal abundance. Furthermore, the decrease in $`M/L`$ will be greater at higher metallicities (i.e. for the brightest galaxies), so that the correlation will get flatter. In order to quantify the slope change, an index $`\eta _X`$ is defined as the slope change between redshift $`z>0`$ and $`z=0`$ using mass-to-light ratios in the $`X`$ passband, namely:
$$\eta _X(z)\frac{\mathrm{\Delta }\mathrm{log}M/L_X}{\mathrm{\Delta }\mathrm{log}M}|_z\frac{\mathrm{\Delta }\mathrm{log}M/L_X}{\mathrm{\Delta }\mathrm{log}M}|_{z=0}$$
(7)
Hence, a $`B_{\mathrm{out}}`$ sequence predicts small negative values for $`\eta _X(z)`$. On the other hand, a $`C_{\mathrm{eff}}`$ sequence is driven by age, so that the faintest galaxies — which will have a younger stellar population — will decrease their $`M/L`$ ratios with redshift faster than the brighter (and older) galaxies. This corresponds to a more significant and positive value for $`\eta _X(z)`$. Figure 5 shows the evolution of the slope (top panels) and zero point (bottom panels) of the correlation between mass-to-light ratio and mass for a range of infall timescales (left) or formation redshifts (right). The points show the observations of a few clusters in restframe $`B`$ band: Coma ($`z=0.02`$), Cl 1358+62 ($`z=0.33`$), Cl 0024+16 ($`z=0.39`$), MS 2053+03 ($`z=0.58`$) and MS1054-03 ($`z=0.83`$) from the compilation in Van Dokkum et al. (1998). Any sequence ($`B_{\mathrm{out}}`$, $`C_{\mathrm{eff}}`$, $`B_{\mathrm{out}}+C_{\mathrm{eff}}`$) will predict a decrease of the zero point as the stellar populations get younger at higher redshift, although the zero point decreases faster for a $`C_{\mathrm{eff}}`$ sequence, for which the average stellar population is younger. Current observations of moderate redshift cluster ellipticals allow an estimate of the zero point but not of the slope with enough accuracy. For instance, the study of Jørgensen et al. (1999) with clusters at low and moderate redshifts seem to indicate a steepening of the correlation between $`M/L`$ and mass, thereby favouring a $`C_{\mathrm{eff}}`$ sequence. However, selection effects could actually mimic such behaviour. Forthcoming spectroscopic observations of cluster ellipticals will enable us to estimate this slope, thereby allowing us to determine the importance of outflows and star formation efficiency in the evolution of cluster early-type galaxies.
## 6 Discussion
A simple phenomenological treatment is described in this paper, where the mechanisms underlying the evolution of galaxies are reduced to a set of four parameters. The star formation rate is assumed to follow a linear Schmidt law whose proportionality constant is used to describe a varying star formation efficiency ($`C_{\mathrm{eff}}`$). The supply of primordial gas fuelling star formation is controlled by gaussian infall characterised by the epoch at which the infall rate is maximum ($`z_F`$), and the width of the gaussian profile gives a characteristic infall timescale ($`\tau _f`$). The model is allowed to eject a fraction ($`B_{\mathrm{out}}`$) of the enriched gas, thereby lowering the effective yield. A first stage in this analysis involves finding a suitable pair of infall parameters ($`z_F`$,$`\tau _f`$) which are capable of generating the restframe $`UV`$ colours of the brightest cluster galaxies. The colour constraints imposed by the reddest (and brightest) galaxies in clusters at moderate and high redshift allow us to discard long infall timescales and recent star formation epochs. One could argue that the constraint on the star formation history of the brightest systems need not be the same as for less massive ellipticals. However, the current most plausible scenarios for galaxy formation assume either a simultaneous process of star formation regardless of galaxy mass, or a hierarchical structure where the most massive galaxies might have undergone the latest bursts of star formation. Furthermore, the possibility of an “inverted-hierarchical” scenario can still be accomodated in this model, as long as the constraints imposed ($`z_F3`$,$`\tau _f1`$ Gyr) are held even for the less massive galaxies.
Out of the four parameters considered in the model, we have found that — within the framework of this model — only the SF efficiency ($`C_{\mathrm{eff}}`$) and the ejected fraction in outflows ($`B_{\mathrm{out}}`$) help determine the mass sequence in cluster early-type galaxies. Infall parameters ($`z_F`$,$`\tau _f`$) do not appreciably change much the output unless very recent stages of star formation are included. However, this will result in restframe $`UV`$ colours that are in contradiction with observations. The efficiency parameter generates a significant spread in the age distribution of stars, although mixed with a metallicity range (any model with a reasonable IMF must include this range of abundances). Alternatively, a range of outflow rates result in a range of metallicities with no significant spread in ages. Both mechanisms are degenerate in local clusters because the average age of the stellar populations predicted for any model at $`z=0`$ are too old to be able to disentangle the effects of age and metallicity. However, the predictions of age-sensitive observables at high redshift differ noticeably for these two sequences. Unfortunately, present data is not capable of ruling out one model against the other, but the continuing flow of data from clusters at moderate-to-high redshift will eventually enable us to break this degeneracy. A sequence driven by efficiency ($`C_{\mathrm{eff}}`$ sequence) predicts a steepening of the slope of the fundamental plane, whereas a $`B_{\mathrm{out}}`$ sequence — driven by outflows — predicts no change or a slight decrease of this slope. High precision observations of the dynamical and spectrophotometric properties of cluster galaxies at high redshift will confirm the importance of either age or metallicity in the mass range of cluster ellipticals.
## Acknowledgments
The authors would like to thank the anonymous referee for useful comments and suggestions.
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# 1 Introduction and Summary
## 1 Introduction and Summary
The study of the AdS/CFT correspondence has been the major occupation of the hep-th-community since the initial breakthroughs of . Amongst the many aspects of the correspondence, one of the most intriguing is the possibility of formulating the field theoretical RG flow in terms of the classical dynamics of the gravitational theory in the bulk .
In studying the RG flow induced by certain operators from an ultra-violet (UV) fixed point, one needs to have a dictionary associating each operator to the appropriate field in the bulk. This requires resolving some possible ambiguity arising for a specific range of bulk masses near the stability bound , which amounts to making a choice between two theories, both in principle described by the same bulk fields. Usually, only one such theory is supersymmetric and the allowed values for the conformal weight $`\mathrm{\Delta }`$ can be read off from the representations of supersymmetry. After the association is made, one still needs to determine whether one is deforming the UV Lagrangian by adding the operator itself or simply going to a different vacuum where the operator acquires a non-zero vacuum expectation value (vev). (For a clear exposition of this point see ). In view of this physical interpretation, not all the solutions to the gravity theory are acceptable, for instance one should rule out flows in which a positive definite operator acquires a negative vev. We will shortly discuss these conditions.
In a related development, it was noticed in that a formulation of gravity in first order Hamiltonian formalism provides further insights into the RG <sup>6</sup><sup>6</sup>6That the RG equations can be given a Hamiltonian structure was noticed some years ago .. The double way of writing the equations – one as a second order Lagrangian system supplemented by the zero energy constraint and the other as a first order system written in terms of a “superpotential” – is the origin of some confusion in the implementation of the holographic RG. The name superpotential in this context is somewhat of a misnomer because, whereas the Hamilton-Jacobi equations have a continuum set of solutions parameterized by the constants of motion, only one such solution can be regarded as the superpotential arising in a supersymmetric theory of gravity. We shall reserve the name superpotential for the truly supersymmetric one and call all the solutions to the Hamilton-Jacobi equation generating functions.
In many circumstances, the few particular generating functions that can be found explicitly are precisely those that can be thought of as true superpotentials. If we are interested in flows between two fixed points of which only one (typically the UV one) appears as an extremum of the known superpotential, we cannot use the first order equations for this purpose and must revert to the Lagrangian system. In this case, it is impossible to obtain an analytical solution and it is impossible to resolve the vev/deformation ambiguity by asymptotically expanding it near the UV fixed point. How can one decide then which of these two cases is realized? Also, does the fact that the known generating function cannot be used to connect two fixed points mean that it is useless in connection with the RG flow? Or, if it can be used, how can the same boundary operator induce different flows?
These questions have been addressed in above cited literature and a coherent picture has emerged. So far, most of the attention has been focused on the case of various gauged supergravities in $`d=5`$, which are, of course, the ones most relevant to four dimensional field theories. We will address and solve these problems very explicitly for a particularly simple example – $`𝒩=1`$ $`d=7`$ gauged supergravity . This example has all the features we want to study and its simple field content allows for a clear-cut solution<sup>7</sup><sup>7</sup>7The case of $`𝒩=1`$ $`d=7`$ gauged supergravity has been recently discussed in , where some comments in the direction of the results of this paper have been made.. By presenting a thorough and explicit solution of the problem, we hope to contribute in clarifying some points that might have remained obscure from the previous discussions.
The status of $`𝒩=1`$ $`d=7`$ gauged supergravity can be summarized as follows: First of all, there is only one scalar field $`\varphi `$ in the bulk and its potential has two extrema (see figure 1), a maximum at $`\varphi =0`$ corresponding to a supersymmetric UV theory and a minimum at<sup>8</sup><sup>8</sup>8In appropriate units to be specified later. $`\varphi =\mathrm{log}2/\sqrt{5}`$, corresponding to a non-supersymmetric but nevertheless stable IR theory. The “tachyonic” excitation near the UV point has a mass $`m`$ given, in units of the AdS radius $`r`$, by $`m^2r^2=8`$. The boundary operator corresponding to $`\varphi `$ is $`𝒪_\varphi =\mathrm{\Phi }^2`$, where $`\mathrm{\Phi }`$ is a scalar in the tensor multiplet of the $`d=6`$ CFT or, better, its still unknown non-Abelian generalization. The conformal dimension of $`𝒪_\varphi `$ is $`\mathrm{\Delta }=4`$. The other possibility ($`\mathrm{\Delta }=2`$) is ruled out by looking at the table 1 of for the multiplets of extended ($`𝒩=2`$) supersymmetry and figure 2 in . In fact, $`\mathrm{\Delta }=2`$ corresponds to the singleton field $`\mathrm{\Phi }`$ itself.
Deforming the UV theory ($`\varphi =0`$) by the addition of $`\varphi 𝒪_\varphi `$ to the fixed point Lagrangian induces an RG flow that ends at the non-supersymmetric IR conformal fixed point. In this case the generating function cannot be obtained explicitly but it can be computed numerically and shown to have the correct behavior at both ends – it corresponds to a particular one among the 1-parameter family of solutions of the Hamilton-Jacobi equation. This is the only solution in which the field $`\varphi `$ is allowed to acquire negative values. All other such solutions correspond to negative vev’s for $`𝒪_\varphi `$ and should be ruled out, in tune with the fact that, when evaluated on solutions that are running away to $`\mathrm{}`$, the potential is not bounded from above and the metric singularity at the runaway point is that of a naked time-like nature .
It turns out that the only solution to the Hamilton-Jacobi equation that has an extremum and is analytic at $`\varphi =0`$ is the superpotential. The superpotential can be used to study new supersymmetric vacua of the theory, for which $`𝒪_\varphi >0`$, by studying runaway solutions in which $`\varphi +\mathrm{}`$. There is also a continuum set of solutions, still with $`\varphi +\mathrm{}`$, describing what we believe are consistent non-supersymmetric vacua.
Towards the end of the paper we will turn to the more complicated case of compactification of $`d=11`$ supergravity on “squashed” manifolds ($`\stackrel{~}{S}^7`$ and $`\stackrel{~}{N}(1,1)`$) and comment on some particular features of the models, complementing the discussion in .
In , the interesting question was raised of whether there exist trajectories connecting the squashed solutions with the corresponding unsquashed manifolds. The situation is similar to the well-studied case of gauged $`d=5`$ supergravity, but there is a crucial difference: In $`d=5`$ supergravity, the analog particle rolls from a saddle point to a minimum of the (inverted) potential, whereas here it should roll from a maximum to a saddle point, clearly a more unstable situation. If the RG equations where truly first order, one could argue from general theorems that there must still be a critical line connecting the points. However, the equations expressed in terms of the potential are second order and there is no guarantee that such a solution will survive. In fact, we have reasons to believe that such a flow does not exist, although more work is required to fully establish or refute this belief.
As far as the squashed solutions are concerned, one is able to find an explicit solution<sup>9</sup><sup>9</sup>9The explicit generating function has no fixed point at the unsquashed vacuum (or else a solution connecting the two would exist). to the Hamilton-Jacobi equation. It turns out that this corresponds to giving a vev to the squashing operator, thus breaking conformal invariance. Given the simple form of the generating function and the collected experience with similar models, it is tempting to conjecture that such a solution is in fact supersymmetric, although here we are working beyond the gauged supergravity truncation and considering fields from the higher levels of the Kaluza-Klein spectrum.
## 2 $`𝒩=1`$ $`d=7`$ gauged supergravity
The field content, Lagrangian and supersymmetry transformations for $`𝒩=1`$ $`d=7`$ gauged supergravity can be found in . For our purposes, we set all fields to zero except for the metric and the scalar $`\varphi `$. The action is
$$S=d^7x\sqrt{g}\left(\frac{1}{2}R\frac{1}{2}(\varphi )^2V(\varphi )\right).$$
(1)
The scalar potential is chosen to be<sup>10</sup><sup>10</sup>10In its full generality, the potential depends on two arbitrary constants $`h`$ and $`g`$ and it displays two minima as long as $`h/g>0`$ (c.f.r. ). One combination of $`h`$ and $`g`$ is eliminated by shifting $`\varphi `$ and the remaining one is an irrelevant overall multiplicative constant in front of the potential.
$$V(\varphi )=\frac{1}{4}e^{8\varphi /\sqrt{5}}2e^{3\varphi /\sqrt{5}}2e^{+2\varphi /\sqrt{5}}.$$
(2)
A plot of $`V(\varphi )`$ is shown in figure 1. There is a supersymmetric UV fixed point at $`\varphi =0`$ and a stable non-supersymmetric IR one at $`\varphi =\mathrm{log}2/\sqrt{5}`$.
The Lagrangian equations of motion following from (1), with the standard domain-wall ansatz
$$ds^2=dy^2+e^{2A(y)}\eta _{\mu \nu }dx^\mu dx^\nu ,\text{and}\varphi =\varphi (y)$$
(3)
can be derived by the following action for a mechanical system
$$S=𝑑ye^{6A}\left(15\dot{A}^2\frac{1}{2}\dot{\varphi }^2V(\varphi )\right).$$
(4)
When supplemented by the zero energy constraint, the equations read <sup>11</sup><sup>11</sup>11The primes denote the derivative with respect to $`\varphi `$ and the dots the derivative with respect to $`y`$.
$`\ddot{\varphi }+6\dot{A}\dot{\varphi }`$ $`=`$ $`V^{}(\varphi )`$ (5)
$`5\ddot{A}+15\dot{A}^2+{\displaystyle \frac{1}{2}}\dot{\varphi }^2`$ $`=`$ $`V(\varphi )`$ (6)
$`15\dot{A}^2{\displaystyle \frac{1}{2}}\dot{\varphi }^2`$ $`=`$ $`V(\varphi ).`$ (7)
Equation (6) can be easily shown to follow from (5) and (7).
Equivalently, one can consider the equation for Hamilton’s characteristic function $`F(A,\varphi ,c)`$ generating the canonical transformations to the cyclic coordinates<sup>12</sup><sup>12</sup>12There is only one constant, $`c`$, because the other conjugate variable is set to zero by (7).
$$\frac{1}{60}\left(\frac{F}{A}\right)^2\frac{1}{2}\left(\frac{F}{\varphi }\right)^2+e^{12A}V=0.$$
(8)
By substituting the ansatz $`F(A,\varphi ,c)=e^{6A}W(\varphi ,c)`$ into (8) the equation becomes the same as the defining equation for the superpotential. Altogether, expressing the canonical transformation in terms of $`W`$ we end up with the first order system of Hamilton-Jacobi equations
$`\dot{\varphi }`$ $`=`$ $`W^{}`$ (9)
$`\dot{A}`$ $`=`$ $`{\displaystyle \frac{1}{5}}W`$ (10)
$`V`$ $`=`$ $`{\displaystyle \frac{1}{2}}W_{}^{}{}_{}{}^{2}{\displaystyle \frac{3}{5}}W^2.`$ (11)
Equation (11) is obeyed by the superpotential of the theory but it also admits a continuum of solutions, parameterized by $`c`$, that have nothing to do with supersymmetry. If one wants to recover all the solutions to the Lagrangian equations this way, one needs to consider all possible solutions to (11).
Particularly confusing is the fact that there are different solutions to (11) that have an extremum at $`\varphi =0`$. One solution, $`W_{\mathrm{susy}}`$, can be easily found by inspection and identified with the superpotential <sup>13</sup><sup>13</sup>13This is defined up to an overall unimportant sign.:
$$W_{\mathrm{susy}}=2e^{\varphi /\sqrt{5}}\frac{1}{2}e^{4\varphi /\sqrt{5}}.$$
(12)
The flow between the two fixed points is generated by another solution, $`W_{\mathrm{ir}}`$, not supersymmetric and not analytic at $`\varphi =0`$ that can only be found numerically. The two functions are plotted for comparison in figure 2.
The function $`W_{\mathrm{ir}}`$ is rather tricky to find directly from (11) but it can be constructed a posteriori once the solution to the Lagrangian system (5), (6), (7) has been found numerically. Such a solution for $`\varphi _{\mathrm{ir}}`$ is presented in figure 3 and can be easily seen to interpolate between the UV and IR fixed points.
In fact, once the solution $`\varphi _{\mathrm{ir}}`$ is found, $`W_{\mathrm{ir}}`$ can be defined as
$`W_{\mathrm{ir}}(z)`$ $`=`$ $`{\displaystyle _{\mathrm{log}2/\sqrt{5}}^z}𝑑w\dot{\varphi }_{\mathrm{ir}}\left(\varphi _{\mathrm{ir}}^1(w)\right){\displaystyle \frac{5}{2^{1/5}\sqrt{3}}}`$ (13)
$`=`$ $`{\displaystyle _{\mathrm{}}^{\varphi _{\mathrm{ir}}^1(z)}}𝑑y\dot{\varphi }_{\mathrm{ir}}(y)^2{\displaystyle \frac{5}{2^{1/5}\sqrt{3}}}.`$
The constant in (13) is chosen to agree with (11) at the IR point. It is interesting to analyze the behaviors of $`W_{\mathrm{susy}}`$ and $`W_{\mathrm{ir}}`$ near the origin. Obviously, $`W_{\mathrm{susy}}`$ is analytic and
$$W_{\mathrm{susy}}(0)=\frac{5}{2},W_{\mathrm{susy}}^{}(0)=0,W_{\mathrm{susy}}^{\prime \prime }(0)=2,W_{\mathrm{susy}}^{\prime \prime \prime }(0)=\frac{6}{\sqrt{5}},$$
(14)
whereas $`W_{\mathrm{ir}}`$ is not analytic, since
$$W_{\mathrm{ir}}(0)=\frac{5}{2},W_{\mathrm{ir}}^{}(0)=0,W_{\mathrm{ir}}^{\prime \prime }(0)=1,W_{\mathrm{ir}}^{\prime \prime \prime }(0)=\mathrm{}.$$
(15)
The two solutions $`W_{\mathrm{susy}}`$ and $`W_{\mathrm{ir}}`$ act as boundaries for a continuum set of solutions that lay between them, all of which have the same behavior as (15) <sup>14</sup><sup>14</sup>14There are also solutions with only an extremum at the IR point which we do not consider ..
Since the second derivative of $`W`$ determines whether the behavior of $`\varphi `$ at $`y+\mathrm{}`$ is square-integrable or not, we see that $`W_{\mathrm{ir}}`$ gives rise to a non-square-integrable behavior, thus corresponding to deforming the fixed point Lagrangian by $`𝒪_\varphi `$.
In fact none of the other solutions (including the superpotential) is physically acceptable in the region $`\varphi <0`$ because they would correspond to giving a negative vev to $`𝒪_\varphi `$, a manifestly positive operator. If we write the asymptotics of $`\varphi `$ as<sup>15</sup><sup>15</sup>15For simplicity we do not write the polynomial corrections. Also recall that $`r^2=15/V(0)=4`$
$$\varphi Ae^{2y/r}+Be^{4y/r}$$
(16)
the above analysis shows that $`A=0`$ for $`W_{\mathrm{susy}}`$ and non-zero for the others. This is shown in figure 4 for the particularly interesting case where the generating function is $`W_{\mathrm{ir}}`$. If we take the case of $`W_{\mathrm{susy}}`$, so that the vev becomes the leading term, we get $`B<0`$, since we are studying the region $`\varphi <0`$. The term $`B`$ should still remain negative by continuity as we use generating functions laying between $`W_{\mathrm{susy}}`$ and $`W_{\mathrm{ir}}`$ and it will reach zero at $`W_{\mathrm{ir}}`$, precisely as $`A`$ reaches zero at the opposite end ($`W_{\mathrm{susy}}`$). $`B`$ corresponds to a vev for $`𝒪_\varphi `$ and therefore a negative value must be excluded. It would be nice to check numerically that our picture is correct but it is rather difficult to isolate the sub-leading term when $`A0`$.
Since the leading exponential behavior for all generating functions is known explicitly, we can be more precise and analyze the differences between positive and negative $`\varphi `$. From (9) and (10), following and doing the asymptotics for large $`|\varphi |`$, one can easily see that the metric has the following behavior, (shifting the singularity to $`y=0`$)
$`\varphi >0:ds^2`$ $`=`$ $`y^2\eta _{\mu \nu }dx^\mu dx^\nu +dy^2`$ (17)
$`\varphi <0:ds^2`$ $`=`$ $`y^{1/8}\eta _{\mu \nu }dx^\mu dx^\nu +dy^2.`$ (18)
Solution (18) corresponds to a naked time-like singularity and our analysis says that it should be excluded. On the other hand, the runaway solution for $`\varphi >0`$ is acceptable and plotted in the neighborhood of the UV point in figure 5.
This solution corresponds to going to new non-supersymmetric vacua where $`𝒪_\varphi >0`$, which for the special case of the superpotential corresponds to a supersymmetric vacuum. The reason we have many acceptable generating functions for $`\varphi >0`$ is that they simply correspond to different vev’s (or vacua), in contrast to the $`\varphi <0`$ case.
## 3 Squashing deformations of $`d=4`$ supergravity
The final issue we discuss is the behavior of $`𝒩=1`$, $`d=3`$ superconformal field theories obtained from solutions of $`d=11`$ supergravity on squashed seven-manifolds. These examples are of interest because they involve fields from higher levels of the Kaluza-Klein tower. They were recently considered in – we would like to make a few remarks complementing that analysis.
The two main examples of manifolds allowing for a squashed solution are the seven-sphere and the manifold $`N(1,1)`$ . Squashed metrics are Einstein metrics, obtained by stretching the original one in some directions. $`N(1,1)`$ is a particular instance of a class of seven-dimensional Einstein manifolds named $`N(p,q)`$ , which has the peculiarity of preserving $`𝒩=3`$ supersymmetries, whereas its squashed version has $`𝒩=0`$.<sup>16</sup><sup>16</sup>16However, the orientation-reversed or “skew-whiffed” solution is supersymmetric, with $`𝒩=1`$ . By the AdS/CFT correspondence all these solutions ought to correspond to some conformal limits of three-dimensional QFT representing the degrees of freedom living on M2-branes placed at the apex of the cone over the compactification manifolds .
By the standard procedure, one reads off the global symmetries (“flavor” and R-symmetry) from the isometries of the corresponding solution. Once an appropriate guess for the gauge group is made, it is possible to get a detailed mapping between the operators and the fields of the KK spectrum, based on matching supersymmetry representations. The theory dual to $`S^7`$ is a 3 dimensional $`𝒩=8`$ CFT with $`SO(8)`$ R-symmetry group, while $`N(1,1)`$ gives rise to a CFT with $`SU(3)\times SU(2)`$ global symmetries and supercharges transforming in the 3 of $`SU(2)`$. Moreover, the situation is more difficult for squashed solutions. Having $`𝒩=1`$ in three dimensions there is no R-symmetry and the usual procedure does not apply straightforwardly. The case of $`N(1,1)`$ is particularly puzzling as squashed and unsquashed solutions share the same global symmetries.
To study the possibility of having domain-wall solutions interpolating between such theories one considers a truncation of the Kaluza-Klein spectrum and derives an effective four-dimensional action for the non-zero fields. The potential for the sphere is known from the work of in terms of two scalars $`u`$ and $`v`$ appearing in the eleven dimensional metric as
$$ds^2=e^{7u}ds^2(AdS_4)+e^{2u+3v}ds^2(\mathrm{base})+e^{2u4v}ds^2(\mathrm{fibre}),$$
(19)
where the seven-sphere is thought of as a $`S^3`$ fibration over the base $`S^4`$. The potential for the squashed $`N(1,1)`$ has been given in terms of four scalars in but for our purposes it is sufficient to repeat the computation of using (19) where now the base manifold is $`CP^2`$ and the fiber is $`RP^3`$, thus obtaining a potential also dependent only on two scalars. In both cases the potential can be written as
$$V(u,v)=\lambda e^{9u}\left(\alpha e^{4v}e^{3v}\frac{1}{32\alpha }e^{10v}\right)+2Q^2e^{21u},$$
(20)
where $`Q`$ is the Page charge and, for the sphere, $`\alpha =1/8`$ and $`\lambda =48`$, whereas, for $`N(1,1)`$, we have $`\alpha =1/16`$ and $`\lambda =24`$. Amusingly, all the physical quantities, such as the conformal dimensions for the operators, turn out to be independent of $`\alpha `$ and $`\lambda `$.
The potential (20) has two fixed points but the field $`u`$ always describes a non-renormalizable (irrelevant) operator. From the equivalent mechanical problem, the flow between these two points would have to connect a maximum of $`V`$ to a saddle point of $`V`$, clearly an unstable situation – contrary to the situation occurring for some flows in $`d=5`$ gauged supergravity, where the “particle” rolls along a valley from a saddle point to a minimum.
As mentioned in the introduction, if the RG equations where truly first order, one could argue from general theorems that there must still be a critical line connecting the points. However, the equations expressed in terms of the potential are second order and the existence of this solution is not guaranteed in this case. After some numerical tests we now believe that there is no such flow.
The potential (20) has another peculiar property: It is possible to find explicitly one generating function $`W`$ that has one critical point at the squashed solution. The function is<sup>17</sup><sup>17</sup>17We take the case of the sphere for concreteness. For the $`N(1,1)`$ case the factor 3 in front of $`e^{2v}`$ is changed to 3/2.
$$W(u,v)=\frac{1}{\sqrt{8}}e^{\frac{9}{2}u}\left(3e^{2v}+6e^{5v}|Q|e^{6u}\right),$$
(21)
which is a solution to
$$V=\frac{16}{63}(_uW)^2+\frac{8}{21}(_vW)^212W^2.$$
(22)
At first, it seems rather counterintuitive that the point with less supersymmetries should appear as an extremum, but we must remember that we are not dealing with a gauged supergravity, where only low-lying KK excitations are included. Still, it is tempting to believe that the solution associated with $`W`$ describes different supersymmetric vacua of the theory. As a check one can show, expanding $`W`$ near its critical point, that the solution corresponds to the operator associated to $`v`$ getting a vev – more specifically $`v\mathrm{exp}(5y/3r)`$. There is a choice between a theory in which the conformal dimension of the operator is $`\mathrm{\Delta }=5/3`$ or $`\mathrm{\Delta }=4/3`$, which is also allowed . Finally, one finds that the runaway solution also satisfies the criterion of boundness from above of the potential. Hopefully, further investigations will reveal if even these models consistently describe holographic RG flows.
## 4 Acknowledgments
We wish thank U. Boscain, U. Gran, J. Kalkkinen and A. Tomasiello for useful discussions. D.M. wishes to thank Chalmers University of Göteborg for financial support and kind hospitality when this work was initiated, and acknowledges partial support from EU TMR program CT960045.
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# 1 Introduction
## 1 Introduction
This review contains four different topics on heavy baryon decays. First we discuss some recent theoretical determinations of the quasielastic $`\mathrm{\Lambda }_b\mathrm{\Lambda }_c`$ form factor in HQET where there are many different results in the literature. The quasielastic $`\mathrm{\Lambda }_b\mathrm{\Lambda }_c`$ form factor has also been calculated in the Relativistic Three Quark Model (RTQM) which is an all-encompassing tool for the description of exclusive heavy baryon decays. We briefly describe the RTQM model and discuss various applications of the RTQM. As concerns QCD sum rules we describe some recent three-loop results on the finite mass baryon current-current correlator at $`𝒪(\alpha _s)`$ which is a new result important for QCD sum rule calculations. Finally we discuss the semi-inclusive decays $`\mathrm{\Lambda }_bX_c+D_s^{()}`$ where $`𝒪(\alpha _s)`$ and $`O(1/m_b^2)`$ have been recently calculated. This is an important mode for $`\mathrm{\Lambda }_b`$-decays with an expected branching ratio of $`10\%`$.
## 2 The quasielastic $`\mathrm{\Lambda }_b\mathrm{\Lambda }_c`$ form factor
There exist many different results on the the quasielastic $`\mathrm{\Lambda }_b\mathrm{\Lambda }_c`$ form factor in the literature. The predicted slope values of the form factor range from $`\rho ^2=0.33`$ to $`\rho ^2=2.35`$. In the simplest approach one takes a heavy quark – light diquark model and describes the transition by a one-loop Feynman diagram . One takes $`M_Q=m_Q`$ and local point coupling factors $`g_1`$ and $`g_2`$ for the quark-diquark-baryon vertices whose strengths are fixed by the compositeness condition. The compositeness condition is nothing but the field theoretic equivalent of the familiar quantum mechanical concept of wave function normalization. In the heavy quark limit the result of such a calculation is given by the form factor $`\mathrm{\Phi }(\omega )=(\omega ^21)^{1/2}\mathrm{ln}(\omega +\sqrt{\omega ^21})`$ which is familiar from the $`\omega `$-dependent renormalization of the heavy quark current. An expansion in terms of powers of $`(\omega 1)`$ shows that the form factor is correctly normalized at the zero recoil point $`\omega =1`$, has a slope of $`\rho ^2=1/3`$ and a convexity of $`c=2/15`$. The form factor $`\mathrm{\Phi }(\omega )`$ is rather flat when e.g. compared to the heavy meson form factor where experimentally one finds slope values of $`1`$. $`\mathrm{\Phi }(\omega )`$ lies within the inclusive HQET sum rule bounds derived by Chiang but must nevertheless be discarded since it oversaturates the semileptonic inclusive rate $`\mathrm{\Lambda }_bX_c+l^{}+\overline{\nu }_l`$ as recently shown in .
Improvements on this simplest approach lead one to the Relativistic Three Quark Model (RTQM) . In the improvements one incorporates binding effects by replacing $`M_Q=m_Q`$ by $`M_Q=m_Q+\overline{\mathrm{\Lambda }}`$. The vertex is softened through introduction of a nonlocal vertex and one introduces a true three-quark structure by replacing the $`(ud)`$-diquark by single $`u`$, $`d`$ quarks.
## 3 The Relativistic Three Quark Model
According to the changes mentioned above the RTQM treats the decay $`\mathrm{\Lambda }_b\mathrm{\Lambda }_c+W_{\mathrm{off}\mathrm{shell}}^{}`$ in terms of a two-loop Feynman diagram with nonlocal vertices including binding effects. The result is that the $`\mathrm{\Lambda }_b\mathrm{\Lambda }_c`$ form factor becomes steeper. Depending on the choice of spin vertex structure for the heavy baryons the slope increases from the aforementioned $`\rho ^2=1/3`$ to $`\rho ^2=0.75÷1.35`$ depending on the choice of spin vertex structure to be discussed later on.
The RTQM is an all-encompassing and versatile tool for the description of heavy baryon decays in terms of a Feynman diagram description. The number of parameters associated with the nonlocality of the vertices, the binding effects and the values of the constituent quark masses is reasonably small and their values lie within common expectations. Many of the parameters are already fixed from light baryon decays where the RTQM also applies. For the loop integrations one uses the $`\alpha `$-parametrization in its exponential form. This introduces $`n`$ $`\alpha `$-parameters $`\alpha _1,\mathrm{},\alpha _n`$ for $`n`$ propagators and consequently $`n`$ integrations. One introduces a Laplace transform to facilitate the vertex form factor integration which is left to the very end. The exponential $`\alpha `$-parametrization allows one to do the tensor loop integrals directly through differentiation, i.e. without use of the Passarino-Veltman expansion. One transforms to spherical type variables which leaves one with one radial type integration and $`(n1)`$ angular type integrations. All $`(n+1)`$ numerical integrations including the Laplace transform can be done with ease. In fact the spherical integrations can also be done analytically but the ease of the numerical integration does not warrant this effort. We shall now discuss several applications of the RTQM to heavy baryon decays.
The dependence of the Isgur-Wise function on the choice of the vertex spin structure for the heavy baryons was investigated in . For the $`\mathrm{\Lambda }_Q`$-type baryons both the effective couplings
$$J_{\mathrm{\Lambda }_Q}^1=\overline{\psi }_Q\psi _{[u}^TC\gamma _5\psi _{d]}J_{\mathrm{\Lambda }_Q}^2=\psi _Q\psi _{[u}^TC\gamma _5v/\psi _{d]}$$
(1)
correctly describe the coupling of the $`\mathrm{\Lambda }_Q`$ to a heavy on-shell quark and two light off-shell quarks in the limit of Heavy Quark Symmetry (HQS). When inserted into the relevant two-loop diagram both coupling structures reproduce the required leading order HQET form factor structure including the unit normalization at zero recoil. However, one finds that the slope of the Isgur-Wise function depends on the choice of vertex structure. For the three choices $`J_{\mathrm{\Lambda }_Q}^1`$, $`\frac{1}{2}(J_{\mathrm{\Lambda }_Q}^1+J_{\mathrm{\Lambda }_Q}^2)`$ and $`J_{\mathrm{\Lambda }_Q}^2`$ one finds slope values of $`\rho ^2=1.35`$, $`1.05`$ and $`0.75`$. The fact that $`\rho ^2(J_{\mathrm{\Lambda }_Q}^1)>\rho ^2(J_{\mathrm{\Lambda }_Q}^2)`$ agrees with the sum rule analysis of although the difference in slope values in is not as large. Note that taking the geometric mean of the two currents leads to a constituent type vertex structure where the projector $`(v/+1)/2`$ projects onto the large components of the light quark fields.
Finite mass effects in heavy $`\mathrm{\Lambda }_{Q_1}\mathrm{\Lambda }_{Q_2}`$ transitions were analyzed in by replacing the heavy quark propagators in the Feynman diagrams by the full propagator. This was effected by the replacement
$$\frac{i}{lv+\overline{\mathrm{\Lambda }}}\frac{v/+1}{2}\frac{i(P/+/l+m_Q)}{(P+l)^2m_Q^2}$$
(2)
where $`(P+l)`$ is the momentum of the heavy quark and $`l`$ is a loop momentum. For $`\mathrm{\Lambda }_b\mathrm{\Lambda }_c`$ the rate is decreased by $`9.3\%`$ in qualitative agreement with the findings of . The decrease was found to be even larger for $`\mathrm{\Lambda }_c\mathrm{\Lambda }_s`$ where the strange quark was treated as a heavy quark in the reference rate for the sake of comparison. It is clear that an expansion of the full propagator in terms of powers of $`1/m_Q`$ would allow one to systematically explore higher order $`1/m_Q`$-effects in these transitions as e.g. in the zero recoil normalization of the relevant zero recoil $`\mathrm{\Lambda }_{Q_1}\mathrm{\Lambda }_{Q_2}`$ form factor.
In the RTQM was used to calculate exclusive nonleptonic decays of heavy baryons. There are so-called factorizing and nonfactorizing contributions to these decays. The nonfactorizing contributions had never been calculated before. In the Feynman diagram approach they involve a genuine three-loop calculation which was done in . As a sample result one finds that in the nonleptonic decays $`\mathrm{\Lambda }_b\mathrm{\Lambda }_c+\pi ^{}`$ the nonfactorizing contributions amount to $`20\%`$ and $`28\%`$ in the parity violating and parity conserving amplitudes, repectively, with an ensuing reduction in rate of $`40\%`$. The nonfactorizing contributions are therefore not negligible. A multitude of exclusive nonleptonic decays have been calculated within the RTQM model involving $`(\overline{s}c)(\overline{u}d)`$, $`(\overline{b}c)(\overline{u}d)`$ and $`(\overline{b}c)(\overline{c}s)`$ transitions .
The RTQM model is also well suited for heavy flavour-conserving one-pion and one-photon transitions between heavy baryons. The one-pion transitions are described by two-loop diagrams where the pion couples to a single light quark line. The $`1/f_\pi `$ coupling of chiral perturbation theory effectively appears through the quark level Goldberger-Treiman relation $`g_\pi =2m_q/f_\pi `$. Many one-pion transitions have been calculated including transitions from excited states . The results are remarkably close to the results of using the constituent quark model for the light quarks even though the light quarks are fully off-shell in the RTQM model.
For one-photon transitions the transverse on-shell photon couples only to the light quarks in the leading order of the heavy quark expansion with a coupling strength given by the light quark charge. In addition one has to include contact graphs to assure gauge invariance of the one-photon transitions. These are generated according to the path integral formalism of Mandelstam. Again the results of the RTQM model are remarkably close to the constituent quark model calculation . The relation of the RTQM description of one-photon transitions to the chiral approach remains to be explored, in particular to the recent calculation of which contains also chiral loops.
What we have discussed so far are some basic applications of the RTQM in the heavy baryon sector. Further work is in progress on the decays of double heavy baryons, on magnetic moments of heavy baryons and on heavy flavour-conserving nonleptonic charm and bottom baryon decays.
## 4 Finite mass baryonic current-current correlator at $`𝒪(\alpha _s)`$
The calculation of the spectral density associated with the baryonic current-current correlator is important for QCD sum rule applications. We want to report on some advances we have made in the calculation of the $`𝒪(\alpha _s)`$ radiative corrections to the spectral density with one finite mass quark mass and two zero quark masses . The calculation involves the evaluation of two-scale three-loop Feynman diagrams which only became possible due to recent technical advances in three-loop technology. Taking the appropiate limits we recover previous results derived for the zero mass case and for the infinite mass case . In the mesonic case the corresponding calculation has been done some time ago showing that radiative corrections to the spectral density can become quite important .
The basic object of study is the vacuum expectation value of the time-ordered product of two baryonic currents $`<TJ(x)J(0)>`$. In spinor space its Fourier transform is expanded along the spinor matrix structures $`q/`$ and $`m`$ with coefficients $`\pi _q(q^2)`$ and $`\pi _m(q^2)`$. We concentrate on the invariant $`\pi _m(q^2)`$ which has associated with it a spectral density $`\rho _m(s)`$ for which we shall present two- and three-loop results. Using the simplest possible current $`J=\mathrm{\Psi }(u^TCd)`$ and writing
$$\rho _m(s)=\frac{1}{128\pi ^4}s^2\left\{\rho _0(s)\left(1+\frac{\alpha _s}{\pi }\mathrm{ln}\left(\frac{\mu ^2}{m^2}\right)\right)+\frac{\alpha _s}{\pi }\rho _1(s)\right\}$$
(3)
we obtain the Born term two-loop contribution $`\rho _0(q^2)=1+9z9z^2z^3+6z(1+z)\mathrm{ln}z`$ where $`z=m^2/q^2`$. In the $`\overline{\text{MS}}`$ scheme the radiative three-loop contribution is given by
$`\rho _1(s)=9+{\displaystyle \frac{665}{9}}z{\displaystyle \frac{665}{9}}z^29z^3\left({\displaystyle \frac{58}{9}}+42z42z^2{\displaystyle \frac{58}{9}}z^3\right)\mathrm{ln}(1z)`$ (4)
$`+\left(2+{\displaystyle \frac{154}{3}}z{\displaystyle \frac{22}{3}}z^2{\displaystyle \frac{58}{9}}z^3\right)\mathrm{ln}z+4\left({\displaystyle \frac{1}{3}}+3z3z^2{\displaystyle \frac{1}{3}}z^3\right)\mathrm{ln}(1z)\mathrm{ln}z`$
$`+12z\left(2+3z+{\displaystyle \frac{1}{9}}z^2\right)\left({\displaystyle \frac{1}{2}}\mathrm{ln}^2z\zeta (2)\right)+4\left({\displaystyle \frac{2}{3}}+12z+3z^2{\displaystyle \frac{1}{3}}z^3\right)\mathrm{Li}_2(z)`$
$`+24z(1+z)\left(\mathrm{Li}_3(z)\zeta (3){\displaystyle \frac{1}{3}}\mathrm{Li}_2(z)\mathrm{ln}z\right)`$
Writing $`q^2=(m+E)^2`$ we can perform a threshold expansion of the spectral density in terms of powers of $`E/m`$. We write the leading order result in a factorized form in order to facilitate comparison with HQET. One has
$`m\rho _m`$ $`\stackrel{m\mathrm{}}{=}`$ $`{\displaystyle \frac{1}{128\pi ^4}}E^5\{1+{\displaystyle \frac{\alpha _s}{\pi }}({\displaystyle \frac{1}{2}}\mathrm{ln}\left({\displaystyle \frac{m^2}{\mu ^2}}\right){\displaystyle \frac{2}{3}})\}^2\times `$ (5)
$`\times \left\{1+{\displaystyle \frac{\alpha _s}{\pi }}\left(4\mathrm{ln}\left({\displaystyle \frac{\mu }{2E}}\right)+{\displaystyle \frac{2}{45}}\left(10\pi ^2+273\right)\right)\right\}.`$
The first bracket is the square of the appropiate HQET matching coefficient $`C(m/\mu ,\alpha _s)`$ first derived in and the second bracket is the appropiate result for the leading order HQET spectral density $`\rho ^{\mathrm{HQET}}(E,\mu )`$ first derived in . We have checked that the zero mass limit of the general spectral density reproduces the result of . Work is in progress on the momentum spectral density $`\rho _q(s)`$ and on correlators of baryonic currents with arbitrary spin structure and on sum rule applications of the spectral densities.
## 5 The semi-inclusive decay $`\mathrm{\Lambda }_bX_c+D_s^{()}`$
Following the analysis of the semi-inclusive B-meson decays $`\overline{B}X_c+D_s^{()}`$ in we looked at the corresponding semi-inclusive $`\mathrm{\Lambda }_b`$-decays. The $`\mathrm{\Lambda }_b`$ decays are potentially more interesting because of the possibility to observe $`\mathrm{\Lambda }_b`$ polarization effects in this decay. At the leading order of $`\alpha _s`$ and the heavy quark mass expansion one expects branching ratios of $`BR(\mathrm{\Lambda }_bX_c+D_s^{})=3.2\%`$ and $`BR(\mathrm{\Lambda }_bX_c+D_s^{})=(4.4(L)+2.4(T))\%`$ where we have separately listed the longitudinal ($`L`$) and transverse component ($`T`$) of the spin 1 $`D_s^{}`$. The two components can be separately measured by an angular analysis of the subsequent decays $`D_s^{}D_s^{}+\gamma `$ and $`D_s^{}D_s^{}+\pi ^0`$. The same holds true for the measurement of polarization effects which will not be discussed here.
In we calculated the perturbative $`𝒪(\alpha _s)`$ and the nonperturbative corrections to these decays using the factorization hypothesis. Numerically one finds:
$`\mathrm{\Lambda }_bX_c+D_s^{}:\widehat{\mathrm{\Gamma }}_S`$ $`=`$ $`(10.0960.013)`$
$`\mathrm{\Lambda }_bX_c+D_s^{}:\widehat{\mathrm{\Gamma }}_L`$ $`=`$ $`0.65(10.1100.034)`$
$`\widehat{\mathrm{\Gamma }}_T`$ $`=`$ $`0.35(10.108+0.026)`$
$`\widehat{\mathrm{\Gamma }}_{L+T}`$ $`=`$ $`(10.0960.009).`$ (6)
In order to clearly exhibit the percentage changes the rates have been normalized to their respective Born term rates. The second and third figures in the round brackets of Eq.(6) refer to the perturbative $`𝒪(\alpha _s)`$ corrections and the nonperturbative kinetic energy correction, respectively. The perturbative corrections are negative and quite uniform. They amount to $`10\%`$. The nonperturbative corrections range from 0.9% to 3.4% with differing signs. The longitudinal mode dominates the rate into $`D_s^{}`$’s. The $`L/T`$ rate ratio $`\mathrm{\Gamma }_L/\mathrm{\Gamma }_T`$ decreases by 6.8% from 1.86 to 1.73 after applying the perturbative and nonperturbative corrections.
The corresponding semi-inclusive $`bu`$ decays $`\mathrm{\Lambda }_bX_u+D_s^{()}`$ are suppressed due to the smallness of $`V_{bu}`$. They are nevertheless of interest for the analysis of so-called wrong sign $`D_s^{()}`$’s . Numerically one finds
$`\mathrm{\Lambda }_bX_u+D_s^{}:\widehat{\mathrm{\Gamma }}_S`$ $`=`$ $`(10.169\mathrm{0..013})`$
$`\mathrm{\Lambda }_bX_u+D_s^{}:\widehat{\mathrm{\Gamma }}_L`$ $`=`$ $`0.73(10.1780.029)`$
$`\widehat{\mathrm{\Gamma }}_T`$ $`=`$ $`0.27(10.115+0.030)`$
$`\widehat{\mathrm{\Gamma }}_{L+T}`$ $`=`$ $`(10.1610.001).`$ (7)
The dominance of the longitudinal mode in the decay $`\mathrm{\Lambda }_bX_u+D_s^{}`$ is now more pronounced. Also the radiative corrections are no longer uniform leading to a substantial 13.3% change in the $`L/T`$ rate ratio due to the perturbative and nonperturbative corrections. It would be interesting to study these semi-inclusive $`\mathrm{\Lambda }_b`$ decay modes including the $`L/T`$ composition of the $`D_s^{}`$’s at future colliders.
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# Metal-insulator transitions in systems with electron-phonon and Coulomb interactions
## I Introduction
The alkali-doped fullerenes, A<sub>n</sub>C<sub>60</sub> (A= K, Rb), raise many interesting problems due to the important role played by both the electron-electron and the electron-phonon interaction. For $`n=3`$ and $`n=4`$ each alkali atom is assumed to donate about one electron into a six-fold degenerate $`t_{1u}`$ band. A<sub>3</sub>C<sub>60</sub> are metals and superconductors, while A<sub>4</sub>C<sub>60</sub> are insulators. The $`t_{1u}`$ orbital is only partly filled for both systems, and band structure calculations predict both systems to be metallic. A<sub>4</sub>C<sub>60</sub> must then be an insulator due to interaction effects neglected in band structure calculations. Under pressure A<sub>4</sub>C<sub>60</sub> becomes metallic, while some fullerenes at normal pressure and with the doping $`n=3`$ (NH<sub>3</sub>K<sub>3</sub>C<sub>60</sub>, Cs<sub>3</sub>C<sub>60</sub>) are not superconductors and probably insulators, but become superconductors under pressure. This suggest that these systems are relatively close to a metal-insulator transition. A metal-insulator transition is usually discussed in terms of the ratio between the Coulomb interaction $`U`$ between two electrons on the same molecule and the one-particle band width $`W`$. For A<sub>3</sub>C<sub>60</sub> and A<sub>4</sub>C<sub>60</sub> this ratio is, however, almost identical. It is then interesting to ask which factors determine whether a system is on the metallic or insulating side.
These systems have been studied in a Hubbard-like model, and it was found that the three-fold degeneracy of the $`t_{1u}`$ orbital plays an important role by increasing the ratio $`U/W`$ where the metal-insulator transition takes place. It was furthermore found that the lattice structure is important for the difference between A<sub>3</sub>C<sub>60</sub> (fcc) and A<sub>4</sub>C<sub>60</sub> (bct). Hubbard-like models predict, however, that A<sub>4</sub>C<sub>60</sub> is an anti-ferromagnetic insulator, while it is known experimentally that there are no moments in A<sub>4</sub>C<sub>60</sub>. The electrons in A<sub>n</sub>C<sub>60</sub> have a relatively strong interaction with Jahn-Teller intramolecular phonons with H<sub>g</sub> symmetry and a weaker interaction with A<sub>g</sub> phonons. The interaction with the Jahn-Teller phonons favor a low spin state and might lead to a nonmagnetic insulator. At the same time there is, however, a Hund’s rule coupling, which favors a high spin state. This leads to an interesting competition between the Jahn-Teller effect and the Hund’s rule coupling. The purpose of this paper is therefore to study the influence of the Jahn-Teller effect and the Hund’s rule coupling on the metal-insulator transition. We do this in the context of the Fullerenes, but similar effects also occur in many other systems, e.g., transition metal compounds.
In this paper we study a model of A<sub>n</sub>C<sub>60</sub> which includes the electron-phonon coupling and the electron-electron interaction. In A<sub>n</sub>C<sub>60</sub> partly occupied orbital $`t_{1u}`$ couples to two nondegenerate phonons of A<sub>g</sub> symmetry and to eight five-fold degenerate (Jahn-Teller) phonons of H<sub>g</sub> symmetry. For A<sub>n</sub>C<sub>60</sub> the important coupling is believed to be to the H<sub>g</sub> phonons, and we therefore study a model with a coupling to one H<sub>g</sub> phonon. To see the effect of having a Jahn-Teller phonon, we also compare with a model containing an A<sub>g</sub> phonon. The models are solved in the dynamical mean-field theory (DMFT) or by using exact diagonalization. To interprete the results we also study analytically a single molecule and a simple two-site model. A brief summary of some aspects of this work are published elsewhere.
In Sec. II we present results for a free molecule and in Sec. III we discuss the typical parameter range. We present results for a two-site model in Sec. IV and describe the DMFT calculation in Sec. V. The results are given in Sec. VI and discussed in Sec. VII.
## II Isolated molecule
We first study an isolated molecule. We consider the case of the coupling to an A<sub>g</sub> phonon and two cases of coupling to Jahn-Teller phonons, namely the $`E\times E`$ case, where a two-fold degenerate level couples to a two-fold degenerate phonon, and the $`T\times H`$ case, where a three-fold degenerate level couples to a five-fold degenerate phonon. The $`E\times E`$ case is the simplest model with a Jahn-Teller effect, and the $`T\times H`$ case is a simple model of a C<sub>60</sub> molecule. Thus we consider the Hamiltonian
$`H_{\mathrm{el}\mathrm{ph}}`$ $`=\epsilon _0{\displaystyle \underset{m}{}}{\displaystyle \underset{\sigma }{}}c_{m\sigma }^{}c_{m\sigma }+\omega _0{\displaystyle \underset{\nu }{}}b_\nu ^{}b_\nu `$ (2)
$`+g{\displaystyle \underset{\nu }{}}{\displaystyle \underset{\sigma }{}}{\displaystyle \underset{m}{}}{\displaystyle \underset{m^{^{}}}{}}V_{mm^{^{}}}^{(\nu )}c_{m\sigma }^{}c_{m^{^{}}\sigma }(b_\nu +b_\nu ^{}),`$
where the first term describes the electronic level, the second the phonon and the third the electron-phonon coupling. The coupling matrices $`V^{(\nu )}`$ are determined by symmetry and they are given in, e.g., Ref. . For the coupling to an A<sub>g</sub> phonon $`V^{(1)}`$ is diagonal with all the diagonal elements equal to unity, $`V_{mm^{^{}}}^{(1)}=\delta _{mm^{^{}}}`$. For the $`E\times E`$ case
$$V^{(1)}=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)\mathrm{and}V^{(2)}=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)$$
(3)
and for the $`T\times H`$ case
$`V^{(1)}={\displaystyle \frac{1}{2}}\left(\begin{array}{ccc}1& 0& 0\\ 0& 1& 0\\ 0& 0& 2\end{array}\right)V^{(2)}={\displaystyle \frac{1}{2}}\left(\begin{array}{ccc}\sqrt{3}& 0& 0\\ 0& \sqrt{3}& 0\\ 0& 0& 0\end{array}\right)`$ (10)
$`V^{(3)}={\displaystyle \frac{\sqrt{3}}{2}}\left(\begin{array}{ccc}0& 1& 0\\ 1& 0& 0\\ 0& 0& 0\end{array}\right)V^{(4)}={\displaystyle \frac{\sqrt{3}}{2}}\left(\begin{array}{ccc}0& 0& 1\\ 0& 0& 0\\ 1& 0& 0\end{array}\right)`$ (17)
$`V^{(5)}={\displaystyle \frac{\sqrt{3}}{2}}\left(\begin{array}{ccc}0& 0& 0\\ 0& 0& 1\\ 0& 1& 0\end{array}\right)`$ (21)
The overall coupling strength is determined by $`g`$. In a metallic molecular solid with narrow bands, there is a simple relation between $`g`$ and the dimensionless coupling $`\lambda `$
$$\lambda _H=\frac{5}{3}N(0)\frac{g^2}{\omega _0}$$
(22)
for the $`T\times H`$ case and
$$\lambda _E=2N(0)\frac{g^2}{\omega _0}$$
(23)
for the $`E\times E`$ case. Here $`N(0)`$ is the electron density of states per spin. We also define a Jahn-Teller energy
$$E_{JT}=\frac{g^2}{\omega _0}.$$
(24)
In addition we include the Coulomb interaction
$`H_\mathrm{U}=`$ $`U_{xx}{\displaystyle \underset{m}{}}n_mn_m+U_{xy}{\displaystyle \underset{\sigma \sigma ^{^{}}}{}}{\displaystyle \underset{m<m^{^{}}}{}}n_{\sigma m}n_{\sigma ^{^{}}m^{^{}}}`$ (25)
$`+`$ $`{\displaystyle \frac{1}{2}}K{\displaystyle \underset{\sigma \sigma ^{^{}}}{}}{\displaystyle \underset{mm^{^{}}}{}}\psi _{m\sigma }^{}\psi _{m^{^{}}\sigma ^{^{}}}^{}\psi _{m\sigma ^{^{}}}\psi _{m^{^{}}\sigma }`$ (26)
$`+`$ $`{\displaystyle \frac{1}{2}}K{\displaystyle \underset{\sigma }{}}{\displaystyle \underset{mm^{^{}}}{}}\psi _{m\sigma }^{}\psi _{m\sigma }^{}\psi _{m^{^{}}\sigma }\psi _{m^{^{}}\sigma },`$ (27)
where $`n_{m\sigma }=\psi _{m\sigma }^{}\psi _{m\sigma }`$ is an occupation number, $`U_{xx}`$ and $`U_{xy}`$ are the Coulomb interactions between equal and unequal orbitals, respectively and $`K`$ is an exchange integral. The Coulomb integrals are related via
$$U_{xy}=U_{xx}2K.$$
(28)
For A<sub>g</sub> phonons and a free molecule, the electron-phonon problem is reduced to a displaced oscillator, and we can easily obtain the electron-phonon contribution to the ground-state energy
$$E^{A_g}(N)=N^2E_{JT}.$$
(29)
where $`E(N)`$ is the energy of a free molecule with $`N`$ electrons. The Jahn-Teller case is more complicated and cannot be solved exactly. We therefore focus on the lowest order contribution to $`E(N)`$. For the Coulomb interaction this contribution is of first order in $`K`$ and for the electron-phonon interaction it is of second order in $`g`$. The results are shown in Table I for the $`E\times E`$ case and in Table II for the $`T\times H`$ case. We have subtracted an average Coulomb energy,
$$\stackrel{~}{E}(N)E(N)\frac{1}{2}N(N1)U_{av}(k),$$
(30)
where $`U_{av}(k)`$ is the average interaction for the ground-state in a full shell with $`2k`$ electrons and the orbital degeneracy $`k`$. By definition, $`\stackrel{~}{E}(N)`$ is then zero for a full shell. Results for the low spin state (singlet or doublet) and the high spin state (triplet or quartet) are shown in Table I.
Neglecting the Coulomb interaction ($`U`$ and $`K`$) the results for the $`T\times H`$ case can be written as
$$E_H^{WC}(N)=\frac{5}{2}C(N)E_{JT},$$
(31)
where
$$C(N)=\{\begin{array}{ccc}1& \mathrm{for}& N=1,5\\ 4& \mathrm{for}& N=2,4\\ 3& \mathrm{for}& N=3\end{array}$$
(32)
This is a weak-coupling result for the case when $`\lambda `$ (or $`g`$) is small. For the strong-coupling case
$$E_H^{SC}(N)=C(N)E_{JT}.$$
(33)
In the strong-coupling case the prefactor in front of $`C(N)E_{JT}`$ has been reduced by a factor of $`5/2`$.
From tables I and II it follows that the low spin state is lower in energy for $`K<2E_{JT}`$ in the $`E\times E`$ case and for $`K<\frac{3}{2}E_{JT}`$ in the $`T\times H`$ case. The results illustrate the competition between the Jahn-Teller effect and the Hund’s rule coupling. According to Hund’s first rule, it is favorable to form a high spin state to minimize the Coulomb energy. On the other hand, the Jahn-Teller effect favors a low spin state. For instance, if the first phonon ($`\nu =1`$) of $`E`$ symmetry (the left matrix in Eq. (3)) is excited, it is favorable to put two electrons (with opposite spins) in the level $`m=2`$, resulting in a singlet. The tables show that in the low spin state the energy $`\stackrel{~}{E}(N)`$ is lowered by the electron-phonon coupling but is increased by the exchange integral $`K`$. In the high spin state, on the other hand, the exchange coupling lowers the energy, while the electron-phonon interaction lowers the energy less than in the low spin state or not at all.
This can be further illustrated by considering the case of two electrons in the $`E\times E`$ system. For $`K<2E_{JT}`$, the important states are
$$|1=\psi _1^{}\psi _1^{}|\mathrm{vac}|2=\psi _2^{}\psi _2^{}|\mathrm{vac}.$$
(34)
These states couple directly via the exchange interaction and indirectly via the electron-phonon interaction. Thus we also include states where one phonon has been excited. In the corresponding Hamiltonian matrix the part corresponding to the states with one phonon are “folded” into the part corresponding to $`|i`$, $`i=1`$ or 2, and we obtain the matrix
$$\left(\begin{array}{cc}U_{xx}6E_{JT}& 2E_{JT}+K\\ 2E_{JT}+K& U_{xx}6E_{JT}\end{array}\right)$$
(35)
If $`K<2E_{JT}`$ the corresponding energy is
$$U_{xx}6E_{JT}+(2E_{JT}+K)=U_{av}(2)8E_{JT}+\frac{8}{3}K,$$
(36)
which clearly shows the competition between the Jahn-Teller coupling and the Hund’s rule coupling. For $`K>2E_{JT}`$ a triplet state of the type $`\psi _1^{}\psi _2^{}|\mathrm{vac}`$ with the energy $`U_{av}(2)4K/3`$ becomes the lowest state.
## III Parameter range
It is now important to establish the parameter range appropriate for A<sub>n</sub>C<sub>60</sub>. In particular, the relative size of $`g`$ and $`K`$ is important, since this determines whether the Jahn-Teller effect or the Hund’s rule coupling wins. The exchange integral $`K`$ has been estimated from an ab initio SCF calculation. From the calculated multiplet splitting the result $`K=0.11`$ eV was obtained. Correlation effects are expected to reduce this number, since correlation effects in particular lower the energies of the low spin states. The difference in Coulomb energy between the low spin states and the high spin states is then reduced and the effective $`K`$ becomes correspondingly smaller. For instance, for some atomic multiplets, correlation effects were found to reduce the multiplet splitting by about 25 $`\%`$. The electron-phonon coupling constants have been estimated from photoemission experiments for a free molecule. Here we replace the eight H<sub>g</sub> modes by one effective mode. The frequency of this mode is chosen as the logarithmically averaged frequency
$`\lambda ={\displaystyle \underset{\nu =1}{\overset{8}{}}}\lambda _\nu ;\lambda \mathrm{ln}\omega _0={\displaystyle \underset{\nu =1}{\overset{8}{}}}\lambda _\nu \mathrm{ln}\omega _\nu `$ (37)
where $`\lambda _\nu `$ and $`\omega _\nu `$ are the electron-phonon couplings and frequencies, respectively. We have calculated the energies of the lowest singlet and triplet states for a free C<sub>60</sub> molecule. The results are shown in Fig. 1. Experimentally it is found that the low spin state wins for both A<sub>3</sub>C<sub>60</sub> and A<sub>4</sub>C<sub>60</sub>, A<sub>4</sub>C<sub>60</sub> being a nonmagnetic insulator and NH<sub>3</sub>K<sub>3</sub>C<sub>60</sub> being antiferromagnetic with a moment (0.7 $`\mu _B`$ per molecule) which corresponds to a spin 1/2 system. For A<sub>4</sub>C<sub>60</sub> the triplet-singlet splitting is estimated to be 0.1 eV. In Fig. 1 this splitting is obtained for $`K0.07`$ eV. As expected this value is smaller than the value $`K=0.11`$ eV deduced from the HF calculation. We observe that a change in the estimate of $`g`$ of course would also lead to a change in this empirical number of $`K`$. Since, however, the present estimate of $`K`$ seems reasonable relative the the HF value, this calculation also gives some support for the value of $`g`$ deduced from PES.
## IV A two-site model
In this section we study a two-site model to gain understanding of a) the different influences of an A<sub>g</sub> and a Jahn-Teller phonon on the metal-insulator transition and b) the competition between the Jahn-Teller effect and the Hund’s rule coupling. We therefore compare the case of an A<sub>g</sub> phonon coupling to a two-fold degenerate level ($`A\times E`$ problem) with the Jahn-Teller $`E\times E`$ problem. Thus we study the model
$$H=\underset{i=1}{\overset{2}{}}[H_{\mathrm{el}\mathrm{phon}}(i)+H_\mathrm{U}(i)]+H_{\mathrm{hop}},$$
(38)
where $`H_{\mathrm{el}\mathrm{phon}}(i)`$ (Eq. (2)) describes the electron-phonon interaction on site $`i`$, $`H_\mathrm{U}(i)`$ (Eq. (25)) describes the Coulomb interaction on site $`i`$, and
$$H_{\mathrm{hop}}=t\underset{m=1}{\overset{2}{}}\underset{\sigma }{}(\psi _{1m\sigma }^{}\psi _{2m\sigma }+\psi _{2m\sigma }^{}\psi _{1m\sigma }),$$
(39)
describes the hopping between the two sites. We have assumed that there is only hopping between equal $`m`$-quantum numbers.
The band gap $`E_g(n)`$ for the filling $`n`$ is given by
$$E_g(n)=E(2n1)+E(2n+1)2E(2n)$$
(40)
for a two-site model, where $`E(N)`$ is the ground-state energy of a system with $`N`$ electrons. It is useful to consider $`E_g`$ in the limit when $`U`$ is very large, and then to extrapolate to intermediate values of $`U`$. The limit of a large $`U`$ is particular transparent, allowing for simple calculations and a qualitative understanding, but it is still relevant for the intermediate values of $`U`$ where the metal-insulator transition takes place. The the metal-insulator transition happens when $`E_g=0`$, i.e., we find the critical value of $`U`$ for which this condition is satisfied. The two-site system is much too small for obtaining reliable quantitative estimates of the critical $`U`$ and the extrapolation of the large $`U`$ results is questionable. Nevertheless, this approach can give a qualitative understanding of more realistic calculations. We consider the limit
$`g<<\omega _0<<W<<U`$ (41)
$`K{\displaystyle \frac{g^2}{\omega _0}}E_{JT}<<W,`$ (42)
where the electron-phonon coupling and the exchange integral are just weak perturbations to the hopping and to the Coulomb interaction.
We first study the case of an A<sub>g</sub> phonon and calculate the gap for the filling $`n=1`$. It is then convenient to transform the coupling term to the form
$$g\underset{i=1}{\overset{2}{}}(n_i1)(b_i+b_i^{}),$$
(43)
where $`n_i=_{m\sigma }\psi _{im\sigma }^{}\psi _{im\sigma }`$ and we have neglected an irrelevant term $`(g^2/\omega _0)_i(2n_i1)`$. For a system with two electrons and in the limit of a large $`U`$, hopping is almost completely suppressed, and each site has almost exactly one electron. The term (43) then has essentially zero coupling to this state, and we obtain
$$E^{A_g}(2)0.$$
(44)
Since there is only one electron per site, there are no multiplet effects. To obtain the energy of a system with one electron, it is convenient to transform to bonding and antibonding operators, e.g., $`\psi _{\pm m\sigma }=(\psi _{1m\sigma }\pm \psi _{2m\sigma })`$. This gives the interaction term
$`{\displaystyle \frac{g}{\sqrt{2}}}[`$ $`(\psi _{+m\sigma }^{}\psi _{+m\sigma }+\psi _{m\sigma }^{}\psi _{m\sigma })(b_++b_+^{})`$ (45)
$`+`$ $`(\psi _{+m\sigma }^{}\psi _{m\sigma }+\psi _{m\sigma }^{}\psi _{+m\sigma })(b_{}+b_{}^{})].`$ (46)
The bonding level is then occupied by one electron, giving the kinetic energy $`t`$. The leading contribution to the energy due to the electron-phonon interaction is $`\frac{1}{2}E_{JT}`$, giving
$$E^{A_g}(1)=t\frac{1}{2}E_{JT}.$$
(47)
In the limit studied here ($`\omega _0<<W`$), the phonons react to the average electronic charge, which gives rise to the factor $`\frac{1}{2}`$.
We next consider the system with three electrons. The wave function for $`g=0`$ and $`K=0`$ is
$`|0={\displaystyle \frac{1}{\sqrt{6}}}(\psi _{11}^{}\psi _{11}^{}\psi _{22}^{}+\psi _{11}^{}\psi _{21}^{}\psi _{12}^{}+\psi _{21}^{}\psi _{11}^{}\psi _{12}^{}`$ (48)
$`+`$ $`\psi _{21}^{}\psi _{21}^{}\psi _{12}^{}+\psi _{21}^{}\psi _{11}^{}\psi _{22}^{}+\psi _{11}^{}\psi _{21}^{}\psi _{22}^{})|\mathrm{vac}.`$ (49)
The corresponding energy is $`0|H|0=2t`$. For a general lattice and for a state which is half-filled apart from one electron or one hole, we expect the hopping to be enhanced by roughly a factor $`\sqrt{N_{\mathrm{deg}}}`$, where $`N_{\mathrm{deg}}`$ is the orbital degeneracy. However, for the two-site system this enhancement is $`N_{\mathrm{deg}}`$. The effects of the electron-phonon coupling are now included in perturbation theory. For the parameter range considered here ($`\omega _0<<W`$), there is only coupling to the states $`b_i^{}|0`$. Furthermore, the Coulomb energy is given by $`0|H_\mathrm{U}|0`$, since there are no states with the kinetic energy -$`2t`$ which couple to $`|0`$ via $`H_\mathrm{U}`$. The energy is therefore
$$E^{A_g}(3)=U_{xx}2t\frac{1}{2}E_{JT}\frac{5}{3}K.$$
(50)
The corresponding band gap is given in Table III ($`A\times E`$), which also shows the result for the filling $`n=2`$. The table shows that the band gap is reduced by the electron-phonon coupling. Extrapolating to intermediate values of $`U`$ then suggests that the electron-phonon coupling increases the critical value of $`U`$ where the metal-insulator transition takes place. The reason is that the electron-phonon interaction in the model (43) is not effective for integer filling in the large $`U`$ case, since the charge fluctuations are then suppressed. For the states with an extra electron or hole, on the other hand, there is a fluctuating charge due to the hopping of the electron or hole and a corresponding lowering of the energy due to the electron-phonon coupling.
The multiplet effects trivially reduce the gap for $`n=1`$, since they are effective for the state with three electrons but not for the other two states. For $`n=2`$, however, they increase the gap. The reason is that in the integer occupation case ($`N=4`$), hopping plays no role for a large $`U`$ and the state can adjust to use the multiplet effects optimally. This is not possible for the case of an extra hole ($`N=3`$) or an extra electron ($`N=5`$), since then the states first of all adjusts to optimize hopping, and only in the second place adjust to optimize the multiplet effects. The increase of the gap due to the exchange integral has been observed earlier.
We next consider the Jahn-Teller case of two-fold degenerate phonons. We first calculate the energy of the state with one electron per site in the large $`U`$ limit. Perturbation theory shows that each phonon contributes an energy $`E_{JT}`$. Since there are two phonons per site and two sites, the energy is
$$E^{E_g}(2)=4E_{JT},$$
(51)
where E<sub>g</sub> labels the Jahn-Teller case. Although charge fluctuations are completely suppressed in this case, the system can still gain energy by introducing a (dynamical) Jahn-Teller distortion. This is in strong contrast to the A<sub>g</sub> phonons, which only couple to the net charge on a given site. In perturbation theory, the energies of the states with an extra electron or hole are
$`E^{E_g}(1)=tE_{JT}`$ (52)
$`E^{E_g}(3)=U_{xx}2tE_{JT}{\displaystyle \frac{5}{3}}K.`$ (53)
The electron-phonon interaction lowers the energy twice as much as in the A<sub>g</sub> case (Eqs. (47) and (50)), simply because the phonon is two-fold degenerate and the coupling therefore is to twice as many phonons as in the A<sub>g</sub> case. The corresponding gap is shown in Table III, which also shows the gap for the filling $`n=2`$.
For $`n=2`$ and $`K<E_{JT}`$ the Jahn-Teller effect dominates over the Hund’s rule coupling. In this case as well as for $`n=1`$ the gap is then increased by the electron-phonon interaction, while in the A<sub>g</sub> case it is decreased. The reason is similar as for the multiplet effects in the A<sub>g</sub> case, discussed in the second paragraph below Eq. (50). For the parameter range (Eq. (41)) we are considering and for integer filling, the hopping is very efficiently suppressed, and the Jahn-Teller system can therefore adjust efficiently to the electron-phonon interaction. With an extra electron or hole, however, the system cannot take advantage of the electron-phonon interaction to the same extent, since the wave function primarily optimizes the hopping of the electron or hole. According to Eq. (40), this leads to an increase of the gap. In the A<sub>g</sub> case, on the other hand, the system cannot couple to the phonons in the integer filling case, and therefore even the reduced coupling to the phonons in the case of an extra electron or hole is sufficient to reduce the gap.
Table III furthermore illustrates the competition between the Jahn-Teller effect and the Hund’s rule coupling. For $`K<2E_{JT}`$ and $`n=2`$ an increase of $`K`$ leads to a reduction of the gap. As $`K`$ is further increased, the system instead tends to go into high spin states, and the gap is increased as K is increased.
It is interesting to compare the results with a different approach. We can calculate an effective on-site $`U_{eff}(n)`$ for the filling $`n`$ as
$$U_{eff}(n)=E(n+1)+E(n1)2E(n),$$
(54)
where $`E(N)`$ now refers to the energy of a free molecule, calculated in Sec. II. For $`K<2E_{JT}`$ and for the $`E\times E`$ case, we then obtain
$`U_{eff}(1)=U_{av}(2)4E_{JT}+{\displaystyle \frac{8}{3}}K`$ (55)
$`U_{eff}(2)=U_{av}+12E_{JT}{\displaystyle \frac{16}{3}}K.`$ (56)
It is then tempting to assume that for the filling $`n=1`$ and $`n=2`$ the gap is given by $`U_{eff}(1)3t`$ and $`U_{eff}(2)4t`$, respectively. Comparison with Table III shows that this is incorrect in the limit studied here (Eq. (41)). In particular, for $`n=1`$, this approach would even predict the wrong sign for the the electron-phonon contribution to the gap. The reason is that Eq. (54) assumes that the electron-phonon and Hund’s rule couplings can adjust to the instantaneous occupation of a given site. As we have seen above, this is not possible in the limit (Eq. (41)) for the states with an extra electron or hole. $`U_{eff}(2)`$ may instead become relevant in the limit when the exchange coupling $`K`$ and the Jahn-Teller energy $`E_{JT}`$ are much larger than the hopping energy.
In a similar way we have calculated the gap for the $`T\times H`$ two-site problem. The results are shown in Table IV. These results also illustrate the competition between the Jahn-Teller effect and the Hund’s rule coupling.
## V Dynamical mean-field calculations
We now consider the $`T\times H`$ problem in the dynamical mean-field theory (DMFT). We formulate the infinite dimensional limit on the Bethe lattice, where the nearest-neighbor hopping integral $`t_{im,jm^{}}`$ is rescaled as $`t^{}/\sqrt{z}`$ with the connectivity $`z`$ going to infinity. We further simplify the hopping by setting $`t_{im,jm^{}}\delta _{mm^{}}`$ only allowing the diagonal hopping. The unit of energy is chosen such that the bandwidth, $`W`$, is set to 2. The Jahn-Teller phonons embedded at each lattice sites of the bath can be easily incorporated into the effective medium of the impurity Anderson model, since they are Einstein phonons without direct intersite couplings between them. Therefore the lattice contribution of phonons is implicitly included in the medium electron Green’s function.
We here focus on the problem where the exchange integral $`K=0`$. The effective impurity Anderson model is solved using the quantum Monte Carlo (QMC) technique with the Fye-Hirsch algorithm, which has only mild “fermion sign-problems” with our Hamiltonian. Here we treat fully quantum mechanically the phonon fields which are updated together with the fermion auxiliary fields in each Monte Carlo step. Details of the implementation of the QMC technique in the DMFT can be found elsewhere in the literature. We have used the discretization step for the Trotter breakup, $`\mathrm{\Delta }\tau =1/3`$, throughout this paper unless mentioned otherwise and more than one million Monte Carlo sweeps are taken for each iteration of the self-consistency loop. For the case where the exchange integral $`K>0`$ the QMC method has a serious “sign-problem”, and this case is treated in Sec. VID using exact diagonalization.
## VI Results
We study how the metal-insulator transition depends on the parameters, e.g., the ratio $`U/W`$ for different strengths of the electron-phonon coupling $`\lambda `$ and the exchange integral. When the system becomes insulating a gap is opened up in the electron spectral function $`A(\omega )`$. This shows up in the electron Green’s function $`G(\tau )`$ calculated for imaginary times $`\tau `$. For instance if
$$A(\omega )=0\mathrm{f}\mathrm{o}\mathrm{r}|\omega |<\mathrm{\Delta },$$
(57)
we obtain
$$G(\tau =\beta /2)e^{\mathrm{\Delta }\beta /2},$$
(58)
where $`\beta =1/T`$. Thus $`G(\beta /2)`$ decays exponentially with $`\beta `$ for an insulator. We therefore use the behavior of $`G(\beta /2)`$ as a measure whether the system is a metal or an insulator. It is also interesting to study the charge fluctuation $`(nn_0)^2`$, which is an average of $`(n_in_0)^2`$ and where $`n_0`$ is the average occupancy per site. This quantity is expected to become small but nonzero at the metal-insulator transition.
### A Jahn-Teller H<sub>g</sub> phonons
We first study the Jahn-Teller H<sub>g</sub> phonons and consider the case of half-filling, i.e., three electrons per site. Fig. 2 shows $`G(\beta /2)`$ as a function of $`U/W`$ for different values of $`\lambda `$. The figure illustrates that for a given $`\lambda `$, $`G(\beta /2)`$ is reduced as $`U`$ is increased and at some critical value $`U_c`$, $`G(\beta /2)`$ becomes very close to zero (not exactly equal to zero due to the finite temperature), where a metal insulator transition takes place. The figure also illustrates that the charge fluctuations $`(n3)^2`$ are strongly reduced in the insulating state. The critical $`U_c/W`$ is reduced as $`\lambda `$ grows. The reason for this was discussed extensively in Sec. IV.
It is interesting that $`U_c`$ as a function of $`\lambda `$ initially is reduced very strongly as $`\lambda `$ is increased, while for larger values of $`\lambda `$ the decrease is slower. This can be understood in terms of the results in Eq. (31, 33) for the Jahn-Teller energy in the weak- and strong-coupling limits. This shows that the electron-phonon energy increases much faster with $`\lambda `$ (by a factor of $`5/2`$) in the weak-coupling limit than in the strong-coupling limit. This is particularly relevant for the large $`U`$ integer filling case, where the electron-phonon interaction has a similar effect as in the free molecule, and it gives a qualitative explanation for the dependence of $`U_c`$ on $`\lambda `$.
### B A<sub>g</sub> versus Jahn-Teller H<sub>g</sub> phonons
We next compare the Jahn-Teller H<sub>g</sub> phonons with A<sub>g</sub> phonons. Fig. 3 shows $`G(\beta /2)`$ in a system with coupling to A<sub>g</sub> phonons as a function of $`U`$ for $`\lambda =0`$ and for $`\lambda =0.8`$. Comparing the results for $`\lambda =0`$ and $`\lambda =0.8`$, we can see that the coupling to A<sub>g</sub> phonons increases the critical $`U_c`$ where the metal-insulator transition takes place, while in the case of H<sub>g</sub> phonons this value is decreased (see Fig. 2). The reason for this change was discussed extensively in Sec. IV.
It is interesting to observe that both the electron-phonon interaction and the electron-electron interaction by themselves would tend to reduce $`G(\beta /2)`$. In a weak-coupling theory where we simply add the lowest order contribution to the self-energy from each interaction, we would then predict that the two interactions work together in reducing $`G(\beta /2)`$. This is indeed what is found in Fig. 3 for $`U/W=0.5`$. Not too surprisingly, the arguments presented below Eq. (52) and assuming the large $`U`$ limit are incorrect for $`U/W=0.5`$. It is not surprising that these arguments become correct for large values of $`U/W`$, but it is important that they are qualitatively correct already for $`U/W1.5`$, which is on the metallic side of the metal-insulator transition.
For a system with only a Coulomb interaction $`U`$, the metal-insulator transition can be thought of as resulting from the competition between kinetic and potential energy. When the coupling to the phonons is introduced the hopping of the electrons is reduced, since an electron tends to drag a cloud of phonons along its path. It is then tempting to assume that the coupling to phonons will move the metal-insulator transition towards smaller values of $`U`$. Our results show that for the parameter range considered here, the effect is the opposite for the case of coupling to A<sub>g</sub> phonons.
The coupling to the A<sub>g</sub> phonons in A<sub>3</sub>C<sub>60</sub> is weak, and the A<sub>g</sub> phonons should not play an important role. There is, however, a strong coupling to a charge carrying plasmon derived from the $`t_{1u}`$ electrons. This plasmon couples in the same way as the A<sub>g</sub> phonons. Its energy (0.5eV) is larger than the phonon energies. However, whether we consider the plasmon energy to be small or large we arrive at the conclusion that it increases $`U_c`$. For small values of the plasmon energy this follows from the study of the A<sub>g</sub> phonons, and for large energies of the plasmon, we can introduce an effective $`U_{eff}`$ as in Eq. (54). The coupling to the plasmon reduces $`U_{eff}`$ and therefore the metal-insulator transition happens for a larger bare $`U`$.
### C A<sub>3</sub>C<sub>60</sub> versus A<sub>4</sub>C<sub>60</sub>
Fig. 4 shows $`G(\beta /2)`$ for A<sub>4</sub>C<sub>60</sub>. The figure illustrates that $`U_c`$ is smaller at filling four than filling three (see Fig. 4). This is not surprising in view of Eq. (31-33) and Table IV, since the energy is lowered more by the electron-phonon coupling for filling four than for filling three. According to Eq. (40) this increases the gap of the insulating state more for filling four than for filling three, i.e., it makes $`U_c`$ smaller for filling four. As discussed in the next section, this result is, however, modified by the Jahn-Teller coupling.
### D Competition between Jahn-Teller and Hund’s rule coupling
We finally discuss the competition between the Jahn-Teller effect and the Hund’s rule coupling. Since there is a sign-problem in the dynamical mean-filed calculations if we use the full multiplet coupling in Eq. (25), we use an exact diagonalization technique. This requires that the system size is small. Therefore we consider a system with just four sites and we consider the $`E\times E`$ problem (Eq. (3)). We furthermore limit the Hilbert space by not allowing more than two phonons per site. To reduce the discreteness of the one-particle spectrum for such a small system, we pick the hopping integrals randomly. We calculate the gap according to Eq. (40) and add some finite size corrections
$$E_g^{red}(U_{xx})=E_g\frac{U_{av}(2)}{N_{site}}E_g^{U=0},$$
(59)
where $`U_{av}(2)/N_{site}`$ is a contribution to the gap from the electrostatic energy of a finite system, $`N_{site}`$ is the number of sites and $`E_g^{U=0}`$ is the gap of a system without any Coulomb interaction. Both these corrections go to zero for a large system. Assuming that that $`E_g^{red}`$ grows linearly with $`U_{xx}`$, we can then use
$$U_{xx}E_g^{red}(U_{xx})$$
(60)
as a crude estimate of the critical $`U_{xx}`$. This quantity is shown in Fig. 5. The results can be qualitatively understood from the results in Table III for a two-site system, although the parameter range considered here is outside the range where the results in the table are valid. For $`\lambda =0`$, the critical value of $`U_{xx}`$ is reduced as $`K`$ is increased, as expected. However, as $`\lambda `$ is increased, the competition between the Jahn-Teller and Hund’s rule effects leads to an increase of the critical $`U_{xx}`$. This is in agreement with Table III, although the increase is faster than for the parameter range of this table. For a particular value of $`\lambda `$, the critical $`U_{xx}`$ becomes comparable to the value for $`K=\lambda =0`$. As $`\lambda `$ is further increased, the Jahn-Teller effect becomes dominating and the critical value of $`U_{xx}`$ decreases again. This critical value is, however, larger than for $`K=0`$, due to the competition between the Jahn-Teller and Hund’s rule coupling.
## VII Discussion
It is interesting to discuss these results in the context of A<sub>n</sub>C<sub>60</sub>. Considering just the Hubbard $`U`$ ($`g=K=0`$), it is found that the metal-insulator transition in A<sub>3</sub>C<sub>60</sub> takes place at the upper range of what is believed to be physical values of $`U/W`$, while for A<sub>4</sub>C<sub>60</sub> this happens at the lower range of these parameters. This agrees nicely with the fact that A<sub>3</sub>C<sub>60</sub> is a metal but A<sub>4</sub>C<sub>60</sub> is an insulator. However, to explain why A<sub>4</sub>C<sub>60</sub> is not antiferromagnetic, we need to include the coupling to the Jahn-Teller H<sub>g</sub> phonons. This substantially lowers the critical $`U_c`$ for both systems, and it puts $`U_c/W`$ of A<sub>3</sub>C<sub>60</sub> at the lower range of the physical parameters. This puts our understanding of A<sub>3</sub>C<sub>60</sub> being a metal into question. We find, however, that the competition between the Jahn-Teller effect and the Hund’s rule coupling increases $`U_c`$ again. The coupling to the $`t_{1u}`$ plasmons in A<sub>3</sub>C<sub>60</sub> should lead to an additional increase of $`U_c`$ in this system. This makes it understandable that A<sub>3</sub>C<sub>60</sub> can be a metal.
## VIII Acknowledgements
This work has been supported by the Max-Planck-Forschungspreis.
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# 𝑆𝑂(10) Unification of Color Superconductivity and Chiral Symmetry Breaking ?
## Acknowledgements
We are indebted to D. H. Friedan for extremely helpful advice and to D. Kaplan and K. Rajagopal for clarifying and stimulating discussions. We also like to thank M. Alford and S.-C. Zhang for interesting discussions about color and high-temperature superconductivity. We gratefully acknowledge the hospitality of the INT in Seattle where this work was done. Our work is supported in part by funds provided by the U.S. Department of Energy (D.O.E.) under cooperative research agreements DE-FC02-94ER40818 and DE-FG02-96ER40945. U.-J. W. also acknowledges the support of the A. P. Sloan foundation.
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# 1 Preparations
## 1 Preparations
Let $`X`$ be a compact Kähler manifold. Then $`𝒦_XH^{1,1}(X,)`$ denotes the Kähler cone, i.e. the open set of all Kähler classes on $`X`$. For a class $`\alpha H^{1,1}(X,)`$ we usually denote by $`\stackrel{~}{\alpha }𝒜^{1,1}(X)_{}`$ a closed real $`(1,1)`$-form representing $`\alpha `$. Let us recall the following version of the Aubin-Calabi-Yau theorem
###### Theorem 1.1
— Let $`X`$ be an $`N`$-dimensional compact Kähler manifold with a given volume form $`\mathrm{vol}𝒜^{N,N}(X)_{}`$. For any Kähler class $`\omega 𝒦_X`$ there exists a unique Kähler form $`\stackrel{~}{\omega }𝒜^{1,1}(X)_{}`$ representing $`\omega `$, such that $`\stackrel{~}{\omega }^N=c\mathrm{vol}`$, with $`c`$.
Since $`\stackrel{~}{\omega }^N`$ is harmonic with respect to $`\stackrel{~}{\omega }`$, this can be equivalently expressed by saying that any Kähler class $`\omega `$ can uniquely be represented by a Kähler form $`\stackrel{~}{\omega }`$ with respect to which the given volume form is harmonic. Note that the constant $`c`$ can be computed as $`c=_X\omega ^N/\mathrm{vol}(X)`$.
###### Definition 1.2
— For a given volume form $`\mathrm{vol}𝒜^{N.N}(X)_{}`$ we let $`\stackrel{~}{𝒦}_X𝒜^{1,1}(X)_{}`$ be the set of Kähler forms $`\stackrel{~}{\omega }`$ with respect to which $`\mathrm{vol}`$ is harmonic.
By the Aubin-Calabi-Yau theorem the natural projection $`\stackrel{~}{𝒦}_X𝒦_X`$ is bijective. But, in general, $`\stackrel{~}{𝒦}_X`$ is not an open subset of a linear subspace of $`𝒜^{1,1}(X)`$ (cf. 2.1). Let $`\stackrel{~}{\omega }\stackrel{~}{𝒦}_X`$. The tangent space of $`\stackrel{~}{𝒦}_X`$ at $`\stackrel{~}{\omega }`$ can be computed as follows. Firstly, we may write $`\stackrel{~}{𝒦}_X=_+\times \stackrel{~}{𝒦}_X^c`$, where $`\stackrel{~}{𝒦}_X^c=\{\stackrel{~}{\omega }𝒦_X|\stackrel{~}{\omega }^N=c\mathrm{vol}\}`$. Secondly, the infinitesimal deformations of $`\stackrel{~}{\omega }`$ in the direction of $`\stackrel{~}{𝒦}_X^c`$ are of the form $`\stackrel{~}{\omega }+\epsilon \stackrel{~}{v}`$, where $`\stackrel{~}{v}`$ is a closed real $`(1,1)`$-form and such that $`(\stackrel{~}{\omega }+\epsilon \stackrel{~}{v})^N=\stackrel{~}{\omega }^N`$. The latter condition gives $`\stackrel{~}{\omega }^N+\left(\genfrac{}{}{0pt}{}{N}{2}\right)\epsilon \stackrel{~}{\omega }^{N1}\stackrel{~}{v}=\stackrel{~}{\omega }^N`$, i.e. $`\stackrel{~}{v}`$ is primitive. As any closed primitive $`(1,1)`$-form is harmonic, this shows that the tangent space of $`\stackrel{~}{𝒦}_X^c`$ at $`\stackrel{~}{\omega }`$ is the space $`^{1,1}(\omega )_{,prim}`$ of real $`\stackrel{~}{\omega }`$-primitive $`\stackrel{~}{\omega }`$-harmonic $`(1,1)`$-forms. Thirdly, the $`_+`$-direction corresponds to the scaling of $`\stackrel{~}{\omega }`$ and this tangent direction is therefore canonically identified with $`\stackrel{~}{\omega }`$. Altogether, one obtains that $`T_{\stackrel{~}{\omega }}\stackrel{~}{𝒦}_X=^{1,1}(\stackrel{~}{\omega })_{}`$ is the space of real $`\stackrel{~}{\omega }`$-harmonic $`(1,1)`$-forms. In particular, $`\stackrel{~}{𝒦}_X`$ is a smooth connected subset of $`𝒜^{1,1}(X)_{}`$. To make this approach rigorous, one completes $`𝒜^{1,1}(X)`$ in the $`L^2`$-topology. The projection of the closed forms to cohomology is a differential map (use e.g. Hodge theory, cf. ). The lifted Kähler cone $`\stackrel{~}{𝒦}_X`$ is the intersection of the space of closed $`L^2`$-forms with the space of sections of the submanifold of the bundle of $`(1,1)`$-forms that consists of those positive forms whose top exterior power equals (a scalar multiple of) the given $`(N,N)`$-form at every point.
###### Definition 1.3
— Let $`X`$ be a compact Kähler manifold with a given volume form. Then one associates to a given Kähler class $`\omega 𝒦_X`$ the space $`^{p.q}(\omega ):=^{p,q}(\stackrel{~}{\omega })`$ of $`(p,q)`$-forms that are harmonic with respect to the unique $`\stackrel{~}{\omega }\stackrel{~}{𝒦}_X`$ representing $`\omega `$.
Note that two different Kähler forms $`\stackrel{~}{\omega }_1`$ and $`\stackrel{~}{\omega }_2`$ representing the same Kähler class $`\omega _1=\omega _2`$ always have different spaces of harmonic $`(1,1)`$-forms. Indeed, $`\stackrel{~}{\omega }_1`$ and $`\stackrel{~}{\omega }_2`$ are $`\stackrel{~}{\omega }_1`$-harmonic respectively $`\stackrel{~}{\omega }_2`$-harmonic. Since any class, in particular $`\omega _1=\omega _2`$, is represented by a unique harmonic form and $`\stackrel{~}{\omega }_1\stackrel{~}{\omega }_2`$, this yields $`^{1,1}(\stackrel{~}{\omega }_1)^{1,1}(\stackrel{~}{\omega }_2)`$. One might ask more generally what the relation is between the spaces of harmonic forms with respect to different Kähler forms not representing the same Kähler class. It is quite interesting to observe that the dependence of $`^{1,1}(\stackrel{~}{\omega })`$ on the Kähler class $`\omega `$ is related to the problem discussed in the introduction. We will try to make this more explicit in the next section.
## 2 How ‘harmonic’ depends on the Kähler form
Let us begin with the following fact which relates the shape of $`\stackrel{~}{𝒦}_X`$ to the dependence of $`^{1,1}(\omega )`$ on $`\omega `$.
###### Proposition 2.1
— The subspace $`^{1,1}(\omega )𝒜^{1,1}(X)`$ is independent of $`\omega `$ if and only if $`\stackrel{~}{𝒦}_X`$ spans an $``$-linear subspace of dimension $`h^{1,1}(X)`$.
Proof. Let $`^{1,1}(\omega )𝒜^{1,1}(X)`$ be independent of $`\omega 𝒦_X`$. Since for any $`\omega 𝒦_X`$ the unique $`\stackrel{~}{\omega }\stackrel{~}{𝒦}_X`$ representing it is $`\stackrel{~}{\omega }`$-harmonic, the assumption immediately yields $`\stackrel{~}{𝒦}_X^{1,1}(\omega )_{}`$ for any $`\omega 𝒦_X`$.
Conversely, if $`\stackrel{~}{𝒦}_X`$ spans an $``$-linear subspace of dimension $`h^{1,1}(X)`$, then this subspace coincides with the tangent space of $`\stackrel{~}{𝒦}_X`$ at every point $`\stackrel{~}{\omega }\stackrel{~}{𝒦}_X`$. But the latter was identified with $`^{1,1}(\omega )_{}`$. Hence, the linear subspace equals $`^{1,1}(\omega )_{}`$ for any $`\omega 𝒦_X`$ and $`^{1,1}(\omega )`$, therefore, does not depend on $`\omega `$. $`\mathrm{}`$
###### Remark 2.2
— The assertion might be rephrased from a slightly different point of view as follows. The bijective map $`\stackrel{~}{𝒦}_X𝒦_X`$ can be used to define a differentiable map $`𝒦_X𝒜^2(X)`$ (in the $`L^2`$-topology). The proposition then just says that this map is linear if and only if the Gauss map is constant. It might be instructive to rephrase some of the results later on in this spirit, e.g. Prop. 3.2.
The next proposition states that the ‘global’ change of $`^{1,1}(\omega )`$ for $`\omega 𝒦_X`$ is determined by the ‘harmonic’ behaviour with respect to a single $`\omega 𝒦_X`$.
###### Proposition 2.3
— Let $`X`$ be a compact Kähler manifold of dimension $`N`$ with a fixed Kähler form $`\stackrel{~}{\omega }_0`$ and volume form $`\stackrel{~}{\omega }_0^N/N!`$. Then the following statements are equivalent:
i) The linear subspace $`^{1,1}(\omega )𝒜^{1,1}(X)_{}`$ does not depend on $`\omega 𝒦_X`$.
ii) For all $`\alpha ^{1,1}(\omega _0)`$ one has $`\alpha ^N^{N,N}(\omega _0)`$.
Proof. Let us assume i). By the previous lemma the lifted Kähler cone $`\stackrel{~}{𝒦}_X`$ spans the $``$-linear subspace $`^{1,1}(\omega _0)`$. Since $`\stackrel{~}{𝒦}_X`$ is open in $`^{1,1}(\omega _0)_{}`$ and all $`\alpha \stackrel{~}{𝒦}_X`$ satisfy the $``$-linear equation
$$\alpha ^N=(_X\alpha ^N/_X\omega _0^n)\omega _0^N$$
(1)
which is an algebraic condition, in fact all $`\alpha ^{1,1}(\omega _0)`$ satisfy (1). Hence, for all $`\alpha ^{1,1}(\omega _0)`$ the top exterior power $`\alpha ^N`$ is harmonic, i.e. ii) holds true.
Let us now assume ii). If $`\alpha ^{1,1}(\omega _0)`$, such that its cohomology class $`\omega :=[\alpha ]`$ is a Kähler class, let $`\stackrel{~}{\omega }\stackrel{~}{𝒦}_X`$ denote the distinguished representing Kähler form of $`\omega `$. If $`\alpha `$ itself is strictly positive definite, then the unicity of $`\stackrel{~}{\omega }`$ and ii) imply $`\alpha =\stackrel{~}{\omega }`$. Thus, the intersection of the closed subset $`^{1,1}(\omega _0)_{}`$ with the open cone of strictly positive definite real $`(1,1)`$-forms is contained in $`\stackrel{~}{𝒦}_X`$. This intersection is non-empty, as it contains $`\stackrel{~}{\omega }_0`$. Since $`\stackrel{~}{𝒦}_X`$ is a closed connected subset of this open cone of the same dimension as $`^{1,1}(\omega _0)_{}`$ this yields $`\stackrel{~}{𝒦}_X^{1,1}(\omega _0)_{}`$. By Prop. 2.1 one concludes that $`^{1,1}(\omega )`$ does not depend on $`\omega 𝒦_X`$. $`\mathrm{}`$
## 3 The positive cone
The next proposition is a first step towards a geometric understanding of the failure of harmonicity of $`\alpha ^N`$ for a harmonic form $`\alpha `$. To state it we recall the following notation.
###### Definition 3.1
— For a compact Kähler manifold $`X`$ the positive cone $`𝒞_XH^{1,1}(X,)`$ is the connected component of $`\{\alpha H^{1,1}(X,)|_X\alpha ^N>0\}`$ that contains the Kähler cone.
Note that by definition $`𝒦_X𝒞_X`$.
###### Proposition 3.2
— If $`X`$ is a compact Kähler cone such that $`𝒦_X`$ is strictly smaller than $`𝒞_X`$, then for any Kähler form $`\stackrel{~}{\omega }`$ there exists a $`\stackrel{~}{\omega }`$-harmonic $`(1,1)`$-form $`\alpha `$ such that $`\alpha ^N`$ is not $`\stackrel{~}{\omega }`$-harmonic.
Proof. Assume that there exists a Kähler form $`\stackrel{~}{\omega }_0`$ such that for all $`\alpha ^{1,1}(\stackrel{~}{\omega }_0)`$ also $`\alpha ^N`$ is $`\stackrel{~}{\omega }_0`$-harmonic. We endow $`X`$ with the volume form $`\stackrel{~}{\omega }_0^N/N!`$. By Prop. 2.3 the lifted Kähler cone $`\stackrel{~}{𝒦}_X`$ is contained in $`^{1,1}(\stackrel{~}{\omega }_0)`$. Since $`𝒦_X`$ is strictly smaller than $`𝒞_X`$ there exists a sequence $`\omega _t𝒦_X`$ converging towards a $`\omega 𝒞_X𝒦_X`$. As $`\stackrel{~}{𝒦}_X`$ is contained in the finite-dimensional space $`^{1,1}(\stackrel{~}{\omega }_0)`$ the lifted Kähler forms $`\stackrel{~}{\omega }_t\stackrel{~}{𝒦}_X`$ will converge towards a form (!) and not only a current $`\stackrel{~}{\omega }^{1,1}(\stackrel{~}{\omega }_0)\stackrel{~}{𝒦}_X`$. As a limit of strictly positive definite forms $`\stackrel{~}{\omega }`$ is still semi-positive definite. Moreover, $`\stackrel{~}{\omega }`$ is strictly positive definite at $`xX`$ if and only if $`\stackrel{~}{\omega }^N`$ does not vanish at $`x`$. By assumption $`\stackrel{~}{\omega }^N=c\stackrel{~}{\omega }_0^N`$ with $`c=_X\omega ^N/_X\omega _0^N`$. Since $`\omega 𝒞_X`$, the scalar factor $`c`$ is strictly positive. Hence, $`\stackrel{~}{\omega }^N`$ is everywhere non-trivial. Thus $`\stackrel{~}{\omega }`$ is strictly positive definite. This yields the contradiction.$`\mathrm{}`$
The interesting thing here is that the proposition in particular can be used to determine the positivity of a class with positive top exterior power just by studying the space of harmonic forms with respect to a single given, often very special Kähler form:
###### Corollary 3.3
— Let $`X`$ be a compact Kähler manifold with a given Kähler form $`\stackrel{~}{\omega }_0`$. If for all $`\stackrel{~}{\omega }_0`$-harmonic $`(1,1)`$-forms $`\alpha `$ the top exterior power $`\alpha ^N`$ is also $`\stackrel{~}{\omega }_0`$-harmonic, then any class $`\omega 𝒞_X`$ is a Kähler class.$`\mathrm{}`$
We conclude this section with a few examples, where the assumption of the corollary is met a priori. In the later sections we will discuss examples where $`𝒦_X`$ is strictly smaller than $`𝒞_X`$ and where Prop. 3.2 can be used to conclude the ‘failure’ of harmonicity.
###### Examples 3.4
i) If $`X`$ is a complex torus and $`\omega `$ is a flat Kähler form, then harmonic forms are constant forms and their products are again constant, hence harmonic. In particular, one recovers the fact that on a torus the Kähler cone and the positive cone coincide.
ii) If for two Kähler manifolds $`(X,\stackrel{~}{\omega })`$ and $`(X^{},\stackrel{~}{\omega }^{})`$ with $`b_1(X)b_1(X^{})=0`$ the top exterior power of any harmonic $`(1,1)`$-forms on $`X`$ or on $`X^{}`$ is again harmonic, then the same holds for the product $`(X\times X^{},\stackrel{~}{\omega }\times \stackrel{~}{\omega }^{})`$. The additional assumption on the Betti-numbers is necessary as the product of two curves shows. Indeed, any $`\phi H^{1,0}(X)`$, for a curve $`X`$, is harmonic, but $`\phi \overline{\phi }`$ is not. Hence, $`\alpha =\phi \times \overline{\phi }+\overline{\phi }\times \phi `$ is a harmonic $`(1,1)`$-form on $`X\times X^{}`$ with non-harmonic $`\alpha ^2`$.
iii) If $`X`$ is a Kähler manifold, such that $`^{1,1}(\omega )`$ does not depend on $`\omega `$, then the same holds for any smooth finite quotient of $`X`$.
iv) For hermitian symmetric spaces of compact type it is known that the space of harmonic forms equals the space of forms invariant under the real form. As the latter space is invariant under products, the Kähler cone of an irreducible hermitian symmetric space coincides with the positive cone.
## 4 K3 surfaces
As indicated earlier the behaviour of the Kähler cone is closely related to the geometry of the manifold. We shall study this in more detail for K3 surfaces. The next proposition follows directly from the description of the Kähler cone of a K3 surface.
###### Proposition 4.1
— Let $`X`$ be a K3 surface containing a smooth rational curve. Then for any Kähler form $`\stackrel{~}{\omega }`$ there exists an $`\stackrel{~}{\omega }`$-harmonic form $`(1,1)`$-form $`\alpha `$ such that $`\alpha ^2`$ is not harmonic.
Proof. If $`X`$ contains a smooth rational curve, then $`𝒦_X`$ is strictly smaller than $`𝒞_X`$ and we apply Prop. 3.2. Indeed, a smooth rational curve $`CX`$ determines a $`(2)`$-class $`[C]`$, whose perpendicular hyperplane $`[C]^{}`$ cuts $`𝒞_X`$ into two parts and $`𝒦_X`$ is contained in the part that is positive on $`C`$.$`\mathrm{}`$
If the harmonicity of the top exterior powers fails for a Kähler manifold with a given Kähler form $`(X,\stackrel{~}{\omega })`$ then it should do so for any small deformation of $`(X,\stackrel{~}{\omega })`$. For a Ricci-flat Kähler structure on a K3 surface the argument can be reversed and one can use the existence of rational curves on arbitrarily near deformations to prove the above proposition on any K3 surface.
###### Corollary 4.2
— Let $`X`$ be an arbitrary K3 surface. If $`\stackrel{~}{\omega }`$ is any hyperkähler form on $`X`$, then there exists a $`\stackrel{~}{\omega }`$-harmonic $`(1,1)`$-form $`\alpha `$ such that $`\alpha ^2`$ is not $`\stackrel{~}{\omega }`$-harmonic.
Proof. Let $`H^0(X,\mathrm{\Omega }_X^2)=\sigma `$. Then
$$\begin{array}{ccc}\hfill ^2(\stackrel{~}{\omega })& =& ^{1,1}(\stackrel{~}{\omega })^{2,0}(\stackrel{~}{\omega })^{0,2}(\stackrel{~}{\omega })\hfill \\ & =& ^{1,1}(\stackrel{~}{\omega })\sigma \overline{\sigma }\hfill \end{array}$$
As the space of harmonic forms only depends on the underlying hyperkähler metric $`g`$, $`^{1,1}(\stackrel{~}{\omega })\sigma \overline{\sigma }`$ contains $`^{1,1}(\stackrel{~}{\omega }_{aI+bJ+cK})`$ for all $`(a,b,c)S^2`$. Here, $`I,J,`$ and $`K`$ are the three complex structures associated with the hyperkähler metric $`g`$.
Assume $`\alpha ^2`$ is $`g`$-harmonic for all $`\alpha ^{1,1}(\stackrel{~}{\omega })`$. Since $`\sigma =\stackrel{~}{\omega }_J+i\stackrel{~}{\omega }_K`$ (up to a scalar factor) and since the product of a harmonic form with the Kähler form is again harmonic, also $`\sigma \overline{\sigma }`$ is harmonic. This implies that $`\alpha ^2`$ is harmonic for all $`\alpha ^2(\omega )`$, as $`\sigma ^2=\overline{\sigma }^2=\alpha \sigma =\alpha \overline{\sigma }=0`$ for $`\alpha ^{1,1}(\stackrel{~}{\omega })`$. Thus, $`\alpha ^2`$ is $`g`$-harmonic for all $`\alpha ^{1,1}(\stackrel{~}{\omega }_{aI+bJ+cK})`$ and all $`(a,b,c)S^2`$. On the other hand, it is well-known that for a non-empty (dense) subset of $`S^2`$ the K3 surface $`(X,aI+bJ+cK)`$ contains a smooth rational curve. Indeed, if $`eH^2(X,)`$ is any $`(2)`$-class, then the subset of the moduli space of marked K3 surfaces for which $`e`$ is of type $`(1,1)`$ is a hyperplane section. This hyperplane section, necessarily, cuts the complete curve given by the base $`^1=S^2`$ of the twistor family. Hence, on one of the K3 surfaces $`(X,aI+bJ+cK)`$ the class $`e`$ represents a smooth rational curve. Contradiction.$`\mathrm{}`$
###### Remark 4.3
— What are the bad harmonic $`(1,1)`$-forms? Certainly $`\stackrel{~}{\omega }^2`$ is harmonic and for any harmonic form $`\alpha `$ also $`\stackrel{~}{\omega }\alpha `$ is harmonic. So, if there is any bad harmonic $`(1,1)`$-form there must be also one that is $`\stackrel{~}{\omega }`$-primitive. Most likely, it is even true that the square of any primitive harmonic form is not harmonic. The proof of it should closely follow the arguments in the proof of Prop. 3.2, but there is a slight subtlety concerning the existence of sufficiently many $`(2)`$-classes, that I cannot overcome for the moment. We sketch the rough idea: Assume there exists a $`\stackrel{~}{\omega }`$-harmonic $`\stackrel{~}{\omega }`$-primitive real $`(1,1)`$-form $`\alpha `$ such that $`\alpha ^2`$ is $`\stackrel{~}{\omega }`$-harmonic. As a $`\stackrel{~}{\omega }`$-harmonic $`\stackrel{~}{\omega }`$-primitive $`(1,1)`$-form $`\alpha `$ is also of type $`(1,1)`$ with respect to any complex structure $`\lambda =aI+bJ+cK`$ induced by the hyperkähler metric corresponding to $`\stackrel{~}{\omega }`$ (see Prop. 7.5 ). Moreover, $`\alpha `$ is also primitive with respect to all Kähler forms $`\stackrel{~}{\omega }_\lambda `$. Assume that there exists a complex structure $`\lambda S^2`$, such that $`𝒞_X[\alpha ]\omega _\lambda `$ is not contained in $`𝒦_X`$. This condition can be easily rephrased in terms of $`(2)`$-classes and thus becomes a question on the lattice $`3U2(E_8)`$. It looks rather harmless, but for the time being I do not know a complete proof of it. Under this assumption, we may even assume that in fact $`\lambda =I`$. Since $`\alpha ^2`$ is harmonic, in fact $`\beta ^2`$ is harmonic for all $`\beta \alpha \stackrel{~}{\omega }^{1,1}(\omega )`$. Going back to the proof of Prop. 3.2, we see that the second part of it can be adapted to this situation and shows that $`\psi ^1(𝒦_X[\alpha ]\omega )\alpha \stackrel{~}{\omega }`$, where $`\psi :\stackrel{~}{𝒦}_X𝒦_X`$. The space $`\psi ^1(𝒦_X[\alpha ]\omega )`$ is the space of the distinguished Kähler forms whose classes are linear combinations of $`[\alpha ]`$ and $`\omega `$. Therefore, all these forms are harmonic and linear combinations of $`\alpha `$ and $`\stackrel{~}{\omega }`$ themselves. To conclude, we imitate the proof of Prop. 3.2 and choose a sequence $`\omega _t𝒦_X[\alpha ]\omega `$ converging towards $`\omega ^{}𝒞_X𝒦_X`$. The corresponding sequence $`\stackrel{~}{\omega }_t\stackrel{~}{𝒦}_X`$ is contained in $`\alpha \stackrel{~}{\omega }`$ and converges towards a form(!) $`\stackrel{~}{\omega }^{}`$. As in the proof of Prop. 3.2 this leads to a contradiction.
## 5 Hyperkähler manifolds
We will try to improve upon Prop. 3.2 in the case of hyperkähler manifolds. In particular, we will replace the question whether the top exterior power $`\alpha ^N`$ of an harmonic form $`\alpha `$ is harmonic by the corresponding question for the square of $`\alpha `$. The motivation for doing so stems from the general philosophy that hyperkähler manifolds should be treated in almost complete analogy to K3 surfaces and that instead of the top intersection pairing one should consider the Beauville-Bogomolov form as the higher dimensional analogue of the intersection pairing for K3 surfaces.
Let us begin by recalling some notations and basic facts. By a compact hyperkähler manifold $`X`$ we understand a simply-connected compact Kähler manifold, such that $`H^0(X,\mathrm{\Omega }^2)=\sigma `$, where $`\sigma `$ is an everywhere non-degenerate holomorphic two-form. A Ricci-flat Kähler form $`\stackrel{~}{\omega }`$ turns out to be a hyperkähler form (cf. ), i.e. there exists a metric $`g`$ and three complex structures $`I`$, $`J`$, and $`K:=IJ`$, such that the corresponding Kähler forms $`\stackrel{~}{\omega }_{aI+bJ+cK}`$ are closed for all $`(a,b,c)S^2`$, such that $`I`$ is the complex structure defining $`X`$, and such that $`\stackrel{~}{\omega }=\stackrel{~}{\omega }_I`$. One may renormalize $`\sigma `$, such that $`\sigma =\stackrel{~}{\omega }_J+i\stackrel{~}{\omega }_K`$. In particular, multiplying with $`\sigma `$ maps harmonic forms to harmonic forms, for this holds true for the Kähler forms $`\stackrel{~}{\omega }_J`$ and $`\stackrel{~}{\omega }_K`$.
The positive cone $`𝒞_XH^{1,1}(X,)`$ is a connected component of $`\{\alpha H^{1,1}(X,)|q_X(\alpha )>0\}`$, where $`q_X`$ is the Beauville-Bogomolov form (cf. ).
###### Proposition 5.1
— Let $`X`$ be a $`2n`$-dimensional compact hyperkähler manifold with a fixed hyperkähler form $`\omega _0`$ and the unique holomorphic two-form $`\sigma `$. Then $`\alpha ^2(\sigma \overline{\sigma })^{n1}`$ is harmonic for all $`\alpha ^{1,1}(\omega _0)`$ if and only if the linear subspace $`^{1,1}(\omega )𝒜^{1,1}(X)`$ does not depend on $`\omega 𝒦_X`$.
Proof. Assume that for all $`\alpha ^{1,1}(\omega _0)`$ also $`\alpha ^2(\sigma \overline{\sigma })^{n1}`$ is harmonic. If $`\alpha `$ is in addition strictly positive definite and $`\stackrel{~}{\omega }\stackrel{~}{𝒦}_X`$ with $`[\alpha ]=\omega `$, then $`\alpha ^2(\sigma \overline{\sigma })^{n1}=\stackrel{~}{\omega }^2(\sigma \overline{\sigma })^{n1}`$. We adapt Calabi’s classical argument to deduce that in this case $`\alpha =\stackrel{~}{\omega }`$: If $`\alpha ^2(\sigma \overline{\sigma })^{n1}=\stackrel{~}{\omega }^2(\sigma \overline{\sigma })^{n1}`$, then $`(\alpha \stackrel{~}{\omega })(\alpha +\stackrel{~}{\omega })(\sigma \overline{\sigma })^{n1}=0`$. Since $`\alpha `$ and $`\stackrel{~}{\omega }`$ are strictly positive definite, also $`(\alpha +\stackrel{~}{\omega })`$ is strictly positive definite. By Lemma 6.1 of also $`(\alpha +\stackrel{~}{\omega })(\sigma \overline{\sigma })^{n1}`$ is strictly positive. As $`[\alpha ]=\omega =[\stackrel{~}{\omega }]`$, the difference $`\alpha \stackrel{~}{\omega }`$ can be written as $`dd^c\phi `$ for some real function $`\phi `$. But by the maximum principle the equation $`(\alpha +\stackrel{~}{\omega })(\sigma \overline{\sigma })^{n1}dd^c\phi =0`$ implies $`\phi const`$. Hence, $`\alpha =\stackrel{~}{\omega }`$.
As in the proof of Prop. 3.2 this shows that the intersection of the closed subset $`^{1,1}(\omega _0)_{}`$ with the open cone of strictly positive definite forms in $`𝒜^{1,1}(X)_{}`$ is contained in $`\stackrel{~}{𝒦}_X`$ and one concludes that $`\stackrel{~}{𝒦}_X^{1,1}(\omega _0)_{}`$.
Hence, $`𝒦_X`$ spans a linear subspace of the same dimension and, by Lemma 2.1 this shows that $`^{1,1}(\omega )`$ is independent of $`\omega 𝒦_X`$.
Conversely, let $`^{1,1}(\omega )`$ be independent of $`\omega 𝒦_X`$. Then $`\stackrel{~}{𝒦}_X^{1,1}(\omega )_{}`$ for any $`\omega 𝒦_X`$. Therefore, $`\alpha ^2(\sigma \overline{\sigma })^{n1}=c(\sigma \overline{\sigma })^n`$ with $`c`$ for $`\alpha `$ in the Zariski-dense open subset $`\stackrel{~}{𝒦}_X^{1,1}(\omega )_{}`$. Hence, $`\alpha ^2(\sigma \overline{\sigma })^{n1}`$ is harmonic for any $`\alpha ^{1,1}(\omega )`$ (cf. proof of Prop. 2.3). $`\mathrm{}`$
Of course, as for K3 surfaces one expects that $`^{1,1}(\omega )`$ does in fact depend on $`\omega `$. This would again follow from the existence of rational curves in every nearby hyperkähler manifold. But it would actually be more interesting to reverse the argument: Assume that $`X`$ is a hyperkähler manifold, such that for any small deformation $`X^{}`$ of $`X`$ the Kähler cone $`𝒦_X^{}`$ equals $`𝒞_X^{}`$. I expect that this is equivalent to saying that $`^{1,1}(\omega )`$ does not depend on $`\omega `$. If for some other reason than the existence of rational curves as used in the K3 surface case this can be excluded, then one could conclude that there always is a nearby deformation $`X^{}`$ for which $`𝒦_X^{}`$ is strictly smaller than $`𝒞_X^{}`$. The latter is expected to imply the existence of rational curves on $`X^{}`$. Along these lines one could try to attack the Kobayashi conjecture, as the existence of rational curves on nearby deformations would say that $`X`$ itself cannot be hyperbolic. Unfortunately, I cannot carry this through even for K3 surface.
## 6 Various other examples
Here we collect a few examples where algebraic geometry predicts the failure of harmonicity of the top exterior power of harmonic two-forms. In all examples this is linked to the existence of rational curves.
Varieties of general type. Let $`X`$ be a non-minimal smooth variety of general type. As I learned from Keiji Oguiso this immediately implies that the Kähler cone is strictly smaller than the positive cone. His proof goes as follows: By definition the canonical divisor $`K_X`$ is big and by the Kodaira Lemma (cf. ) it can therefore be written as the sum $`K_X=H+E`$ of an ample divisor $`H`$ and an effective divisor $`E`$ (with rational coefficients). Consider the segment $`H_t:=H+tE`$ with $`t[0,1)`$. If all $`H_t`$ were contained in the positive cone $`𝒞_X`$, then $`K_X`$ would be in the closure of $`𝒞_X`$. If the Kähler cone coincided with the positive cone $`𝒞_X`$, then $`K_X`$ would be nef, contradicting the hypothesis that $`X`$ is not minimal. Hence $`t_0:=sup\{t|H_t𝒞_X\}(0,1)`$. If $`H_{t_0}`$ is not nef, then $`𝒦_X`$ is strictly smaller than $`𝒞_X`$. Thus, it suffices to show that $`H_{t_0}`$ is not nef. If $`H_{t_0}`$ were nef then all expressions of the form $`H_{t_0}^{Ni}.H^{i1}.E`$ would be non-negative. Then $`0=H_{t_0}^N=H_{t_0}^{N1}(H+t_0E)=H_{t_0}^{N1}.H+t_0H_{t_0}^{N1}.E`$, so both summands must vanish. In particular, $`0=H_{t_0}^{N1}.H=H^2.H_{t_0}^{N2}+t_0H.H_{t_0}^{N2}.E`$ Again this yields the vanishing of both terms and in particular $`0=H^2.H_{t_0}^{N2}`$. By induction we eventually obtain $`0=H^{N1}.H_{t_0}`$ and, furthermore, $`0=H^{N1}.H_{t_0}=H^N+t_0H^{N1}.E`$. But this time $`H^N>0`$ yields the contradiction. Therefore, for a non-minimal variety of general type one has $`𝒦_X𝒞_X`$ and hence there exist harmonic (with respect to an arbitrary Kähler metric) two-forms with non-harmonic top exterior power. Note that a non-minimal variety contains rational curves. As the reader will notice, the above proof goes through on an arbitrary manifold $`X`$ that admits a big, but not nef line bundle $`L`$ (replacing the canonical divisor). Also in this case the positive cone and the Kähler cone differ.
Birational Calabi-Yau. Let $`X`$ and $`X^{}`$ be birational Calabi-Yau manifolds, i.e. $`K_X`$ and $`K_X^{}`$ are trivial, then the birational map extends to an isomorphism or there exist harmonic $`(1,1)`$-forms on $`X`$, such that their top exterior power is not harmonic. Again, a non-trivial birational correspondence produces rational curves. As one expects for hyperkähler manifolds that $`^{1,1}(\omega )`$ does depend on the hyperkähler form even when $`X`$ does not contain a rational curve, e.g. for K3 surfaces, it would be interesting to see an example of a simply-connected Calabi-Yau manifold (in particular not a torus), where it does not.
The same argument could be applied to the case of different birational minimal models (minimal models are not unique!). This shows that in the previous example the Kähler cone could be strictly smaller than the positive cone, even when $`K_X`$ is nef or ample.
Blow-ups. This example is very much in the spirit of the previous two. Let $`X`$ be a non-trivial blow-up of a projective variety $`Y`$. Then $`𝒦_X`$ is strictly smaller than $`𝒞_X`$ and, therefore, for any Kähler structure on $`X`$ there exist harmonic $`(1,1)`$-forms with non-harmonic maximal exterior power. Indeed, if $`L`$ is an ample line bundle on $`Y`$ then $`f^{}(L)`$ is nef, but not ample, and it is contained in the positive cone. Hence, $`f^{}(L)𝒞_X𝒦_X`$. Note that also the first example could be proved along these lines. By evoking the contraction theorem one shows that any non-minimal projective variety $`X`$ admits a non-trivial contraction to a projective variety $`Y`$. The above argument then yields that $`𝒦_X`$ and $`𝒞_X`$ are different.
## 7 Chern forms
Let $`X`$ be a compact Kähler manifold with a Ricci-flat Kähler form $`\stackrel{~}{\omega }`$. If $`F`$ denotes the curvature of the Levi-Cevita connection $``$, then the Bianchi identity reads $`F=0`$. The Kähler-Einstein condition implies $`\mathrm{\Lambda }_\omega F=0`$. The last equation can be expressed by saying that $`F`$ is $`\stackrel{~}{\omega }`$-primitive. Analogously to the fact that any closed primitive $`(1,1)`$-form is in fact harmonic, one has that for $`F`$ with $`F=0`$ the primitivity condition $`\mathrm{\Lambda }_{\stackrel{~}{\omega }}F=0`$ is equivalent to the harmonicity condition $`F=0`$. As for untwisted harmonic $`(1,1)`$-forms one might ask for the harmonicity of the product $`F^m`$. Slightly less ambitious, one could ask whether the trace of this expression, an honnest differential form, is harmonic. This trace is, in fact, a scalar multiple of the Chern character $`ch_m(X,\stackrel{~}{\omega })`$.
Question. — Let $`(X,\stackrel{~}{\omega })`$ be a Ricci-flat Kähler manifold. Are the Chern forms $`ch_m(X,\stackrel{~}{\omega })`$ harmonic with respect to $`\stackrel{~}{\omega }`$ ?
By what was said about K3 surface we shall expect a negative answer to this question at least in this case:
Problem. — Let $`X`$ be a K3 surface with a hyperkähler form $`\stackrel{~}{\omega }`$. Let $`c_2A^2(X)`$ be the associated Chern form. Show that $`c_2`$ is not harmonic with respect to $`\stackrel{~}{\omega }`$ !
So, this should be seen in analogy to the fact that $`\alpha ^2`$ is not harmonic for any primitive harmonic $`(1,1)`$-form $`\alpha `$. Here, $`\alpha `$ is replaced by the curvatue $`F`$ and $`\alpha ^2`$ by $`trF^2`$. It is likely that the non-harmonicity of $`c_2`$ can be shown by standard methods in differential geometry, in particular by using the fact that $`c_2`$ is essentially $`F\stackrel{~}{\omega }^2`$ (see ), but I do not know how to do this.
Furthermore, it is not clear to me what the relation between the above question and the one treated in the previous sections is. I could imagine that the non-harmonicity of $`ch_m`$ in fact implies the existence of harmonic $`(1,1)`$-forms with non-harmonic top exterior power.
Acknowledgements. I wish to thank U. Semmelmann for his interest in this work and M. Lehn and D. Kaledin for making valuable comments on a first version of it. I am most grateful to Keiji Oguiso for its enthusiastic help with several arguments.
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# The east–west effect for atmospheric neutrinos
## 1 Introduction
The precise calculation of the fluxes of atmospheric neutrinos has become a much more interesting task since the measurements of Super–Kamiokande and other detectors have given convincing evidence that $`\nu `$ flavor transitions exist and are modifying the intensity and distorting the angular distributions of the $`\nu `$ fluxes. As the field of atmospheric neutrino is maturing from the “discovery era” to the era of “precision measurements” the need for detailed, accurate predictions for the the atmospheric neutrino fluxes has clearly grown in importance, since comparison of data and predictions is now used to extract from the data the parameters that describe exciting new physics. At the same time the task of calculating the expected fluxes is now more demanding, since errors in the calculation, or the use of flawed data as input can result in the introduction of biases in the estimate of these parameters.
The study of the zenith angle distributions of the neutrino fluxes is obviously of central importance in the determination of the oscillation parameters. The neutrino pathlength $`L`$ (and the density profile encountered by a neutrino during its path) are strictly correlated with the zenith angle $`\theta _\nu `$ (the correlation is not perfect because the neutrino creation point along the straight line defined by its momentum at detection is not exactly known). The neutrino flavor transition probabilities depend on the pathlength $`L`$ (or more in general on the neutrino path), and therefore leave their imprints on the zenith angle distributions.
The azimuth angle distributions are much less interesting. The only significant source of non–flatness for these distributions are the effects of the geomagnetic field. In fact, neglecting the existence of the magnetic field, the geometry of the neutrino source volume (a sherical shell of air) and of the volume where neutrinos propagate, with very good approximation have cylindrical symmetry for rotations around the vertical axis<sup>1</sup><sup>1</sup>1 Other sources of asymmetry can in principle be the presence of mountains (especially above the detector) and the existence of different profiles of the atmosphere in different geographycal locations. The first effect is easily calculable. Both effects are expected to be negligibly small. . Neutrino oscillations (or in fact any other forms of transition or disappearance proposed as an explanation for the atmospheric $`\nu `$ data) do not disturb this symmetry<sup>2</sup><sup>2</sup>2 If an east–weast asymmetry is generated by some other source (for example geomagnetic effects), neutrino oscillations can in principle affect the size of the asymmetry, since oscillations can distort the $`\nu `$ energy spectrum. This effect is of second order, and a small correction to the calculated east–west asymmetry.. The fact that the azimuth angle distributions are independent from the new physics that is investigated, has however a positive consequence, their study can provide an important cross check for the theoretical predictions and the experiments. The calculations must be able to predict the shape of these distributions, without the need of additional (unknown) physics beyond the standard model, and experiments must be able to measure angular effects that are unambiguously predicted.
### 1.1 The hint of a discrepancy between data and prediction
The fluxes of (positively charged) primary cosmic rays that reach the Earth’s atmosphere contain more particles traveling from west toward east than in the opposite direction. The discovery of this effect in 1933 allowed to determine that the dominant component of the cosmic rays is positively charged. This “primary asymmetry” is also reflected in the intensity of the fluxes of secondary particles that are generated in the showers of the primary particles and have a direction correlated with the primary one.
The Super–Kamiokande (SK) detector has measured the azimuth angle distributions of atmospheric neutrinos, finding as expected that the distributions are not flat, and that there in an excess of particles traveling toward east. In the SK analysis the east–west asymmetry is defined as:
$$A=\frac{\left(N_EN_W\right)}{\left(N_E+N_W\right)}$$
(1)
where $`N_E`$ ($`N_W`$) is the number of events with the detected charged lepton traveling toward east (west). After selecting events with the charged lepton ($`\mathrm{}=e,\mu `$) in the momentum interval $`p_{\mathrm{}}\left(\mathrm{GeV}\right)=[0.4,3]`$ and the zenith angle region $`\left|\mathrm{cos}\theta _{\mathrm{}}\right|<0.5`$, the measured asymmetries are:
$$A_e^{\mathrm{SK}}=0.21\pm 0.04,A_\mu ^{\mathrm{SK}}=0.08\pm 0.04$$
(2)
This has to be compared with predictions from Honda, Kasahara, Kajita and Midorikawa (HKKM) and the Bartol group :
$$A_e^{\mathrm{HKKM}}=0.13A_\mu ^{\mathrm{HKKM}}=0.11$$
(3)
$$A_e^{\mathrm{Bartol}}=0.17A_\mu ^{\mathrm{Bartol}}=0.15$$
(4)
Considering the statistical errors and the spread in the theoretical prediction, the SK collaboration concludes that both measured asymmetries are compatible with the expectations. However, if we consider the ratio $`A_e/A_\mu `$, we can see that there is a hint (at the level of 2–2.5 standard deviations) of a discrepancy between the calculations that predict asymmetries that are approximately equal for $`e`$–like and $`\mu `$–like events (with only a small excess for $`e`$–like events) and the data where $`A_e`$ is events is more than twice larger that the asummetry for $`\mu `$–like events. We also note that there is an indication that the measured asymmetry for $`e`$–like events is larger than the prediction while on the contrary the prediction for $`\mu `$–like events is smaller.
A discrepancy between data and calculation for the azimuthal distributions can have a non negligible importance a for the interpretation of the data. The study of the east–west aymmetry probes the atmospheric neutrino fluxes in the angular region around the horizontal plane that is actually the most important one for the determination of $`\mathrm{\Delta }m^2`$. The fact that the results for $`e`$–like and $`\mu `$–like events differ from the expectations, and in opposite directions, in a situation where $`\nu `$ oscillations cannot play a role, is clearly important for the studies of neutrino oscillations, where the comparison of the events rates for the two lepton types is a cornerstone of the analysis.
In this work we predict that with increasing statistics the Super–Kamiokande measurement of the east–west asymmetry will clearly show to be in disagreements with the predictions of and , and indeed with any other prediction based on a calculation that does not take into account the effects the bending in the geomagnetic field of secondary charged particles in the cosmic ray showers (even if the calculation is three dimensional ). The existing calculations predict an east–west asymmetry that is simply the reflection of the asymmetry of the primary cosmic ray. In this case the asymmetries for $`e`$–like and $`\mu `$–like events (and indeed for all for neutrino types $`\nu _e`$, $`\overline{\nu }_e`$, $`\nu _\mu `$, $`\overline{\nu }_\mu `$). As we will discuss in the following the effect of the magnetic field in the development of the shower in the atmosphere results in different asymmetries for the different neutrino types, in good agreement with the Super–Kamiokande data.
This work is organized as follows. In the next section we rapidly summarize the effects of the geomagnetic field on the primary cosmic rays. Section 3 contains a discussion of the effects of the geomagnetic field on the shower development. Section 4 discusses the effects of the bending of the muon trajectories in in the geomagnetic field. This is the key mechanism that enhances or suppresses the asymmetry for the different neutrino flavors. In section 5 we show the results of a complete calculation that includes the bending in the geomagnetic fields of all secondary charged particles. Section 6 gives a summary and some conclusions. An appendix contains a brief discussion on the difficulty of performing a detailed and accurate “three–dimensional” calculation of the atmospheric neutrino fluxes and on the possible strategies that are under study.
## 2 The East–West effect for the primary cosmic rays
The flux of primary cosmic rays protons and nuclei that arrive in the vicinity of the Earth surface exhibits the well known east–west effect, that is there are more particles traveling from west toward east than in the opposite direction. This is due to the effects of the geomagnetic field that forbids the lowest rigidity particles from reaching the Earth. The effect is strongest for west–ward going positively charged particles.
As an illustration let us approximate the geomagnetic field as a dipole, and study the trajectories of charged particles in the equatorial plane. The magnetic field is orthogonal to the plane, pointing “up” toward the magnetic pole in the northern hemisphere. Near the surface of the Earth the field has an approximately constant value $`B0.31`$ Gauss. A positively (negatively) charged particle with unit charge and momentum $`p59`$ GeV can have a a clockwise (counterclockwise) circular trajectory that grazes the Earth along the equator going from east to west (from west to east)<sup>3</sup><sup>3</sup>3The trajectory is unstable even in the case of an exactly dipolar field. It is clear that an observer at the equator looking near the horizon toward magnetic east (west) can only see unit charge primary particles with a momentum larger than 59 GeV, all lower momentum particles have “forbidden trajectories”. If the equatorial observer turns around by 180 degrees, he/she can see primary particles of much lower momenta. This is qualitatively the origin of the east–west effect.
The problem of calculating the geomagnetic effects on the primary cosmic rays (before they reach the atmosphere) can be reduced to the problem of calculating “allowed” and “forbidden trajectories”. The definition of allowed and forbidden trajectories is the following. Let us consider a cosmic ray of charge $`q`$, momentum $`\stackrel{}{p}`$ near the Earth surface in position $`\stackrel{}{x}`$, and let us study the past trajectory of the particle. This study can have three results:
* the trajectory originates from the Earth’s surface;
* the trajectory remains confined in the volume $`R_{}<r<\mathrm{}`$ without ever reaching “infinity”;
* the particle in the past was at very large distances from the Earth.
Trajectories belonging to the classes (a) and (b) are considered as “forbidden”, because no primary cosmic ray particle can reach the Earth from a large distance traveling along one of these trajectories. All other trajectories are allowed.
The algorithm (used in all calculations of the atmospheric neutrino fluxes) to obtain the primary cosmic ray fluxes as a function of position and angle is based on two steps:
1. It is assumed that the primary cosmic rays at a distance of one astronomical unit from the sun are isotropic, unless they are disturbed by the presence of the Earth.
2. The flux arriving at an imaginary surface in the vicinity of the Earth (the “top of the atmosphere”) is equal to the isotropic primary flux after the subtraction of all “forbidden trajectories”.
The theoretical basis for the “subtraction agorithm” is solidly motivated , and is based on two fundamental assumptions, that the cosmic ray spectrum “in the absence of the Earth” is isotropic, and that the field around the Earth is well described by a static purely magnetic field. The Liouville theorem states that the density of points in phase space volume is constant. The flux $`\varphi (\stackrel{}{p},\stackrel{}{x})`$ is in fact proportional to the phase space density for relativistic particles, and since the momentum of a charged particle in a magnetic field is constant, then the differential flux $`\varphi (\stackrel{}{p},\stackrel{}{x}\left(t\right))`$ along a particle trajectory is constant. It follows that if the flux is isotropic at large distances from the Earth, then the flux in a small cone around any trajectory is constant and independent from the position. If a particle can reach the Earth, then the momentum spectrum, and the angular distribution around it are not deformed.
For the explicit calculation of allowed and forbidden trajectories, the most accurate method is the so called “back–tracking method” that is a straightforward direct application of the definition outlined above. Given a detailed map of the magnetic field around the Earth (for example see ) it is a straightforward exercise to compute numerically the past trajectories of a cosmic ray from a point “just above” the atmosphere, and test if it corresponds to an allowed or forbidden trajectory.
It is well known that the problem of calculating allowed and forbidden trajectories can be solved analytically for the special case of a volume that is entirely filled with an exactly dipolar magnetic field (there are no class (a) trajectories in this model). In this case the geomagnetic effects result in a sharp rigidity cutoff”<sup>4</sup><sup>4</sup>4In a numerical calculation for a non exactly dipolar field, or even in the case of a dipole, if the trajectories that have a segment inside the Earth are also considered as forbidden, there is not anymore a sharp cutoff, but in a narrow interval allowed and forbidden bands of rigidity alternate.. For a given position $`\stackrel{}{x}`$ and a given direction $`\widehat{n}`$ the trajectories of all positively (negatively) particles with rigidity $`R>R_S^+`$ ($`R<R_S^{}`$) are allowed and the trajectories of all particles with $`R<R_S^+`$ ($`R>R_S^{}`$) are forbidden because they remain confined to a finite distance from the dipole center. The quantity $`R_S^\pm (\stackrel{}{x},\widehat{n})`$ is the Störmer rigidity cutoff :
$$R_S^+(r,\lambda _M,\theta ,\phi )=\left(\frac{M}{2r^2}\right)\left\{\frac{\mathrm{cos}^4\lambda _M}{\left[1+\left(1\mathrm{cos}^3\lambda _M\mathrm{sin}\theta \mathrm{sin}\phi \right)^{1/2}\right]^2}\right\}$$
(5)
where we have made use of the cylindrical symmetry of the problem, $`M`$ is the magnetic dipole moment of field, $`r`$ is the distance from the dipole, $`\lambda _M`$ the magnetic latitude, $`\theta `$ the zenith angle and $`\phi `$ an azimuth angle measured clockwise from magnetic north. The Störmer cutoff has been of considerable historical importance. It is not sufficiently accurate for a modern detailed numerical study. but it remains a good approximation that contains all qualitative features of an “exact” numerical calculation, and remains a very valuable tool to gain physical insight and qualitative understanding. We can use it to ilustrate the three most important qualitative points that we will need for the discussion in this work.
1. The geomagnetic effects depend on the detector position. They are strongest near the equator and weakest at the magnetic poles. In equation (5) we can see that the value of the cutoff grows monotonically from a vanishing value at the magnetic pole ($`\lambda _M=\pm 90^{}`$) to a maximum value at the magnetic equator ($`\lambda _M=0^{}`$).
2. The set of allowed rigidities is largest for particles traveling toward magnetic east. In equation (5) for a fixed detector position (fixed $`\lambda _M`$) and a fixed zenith angle $`\theta `$, the cutoff is minimum (maximum) for azimuth angle $`\phi =270^{}`$ ($`\phi =90^{}`$) that corresponds to a particle traveling toward east (west).
3. For a fixed detector position and a fixed azimuth angle the set of allowed rigidities is largest (smallest) for vertical (horizontal) particles. In equation (5) for $`\lambda _M`$ and $`\phi `$ fixed, the cutoff grows monotonically with $`\theta `$.
The east–west aymmetry of the primary cosmic ray radiation is of course reflected in asymmetries for the fluxes of secondary particles generated in their showers in the atmosphere. This effect results in asymmetries of approximately the same size for all four neutrino types. This is a consequence of the fact that a primary particle of energy $`E_0`$ produces yields of $`\nu _e`$, $`\overline{\nu }_e`$, $`\nu _\mu `$ and $`\overline{\nu }_\mu `$ with approximately the same energy spectrum (and relative normalizations in the ratios $`1:1:2:2`$). The existing calculations of the neutrino fluxes only include the geomagnetic effects on the the primary flux as a source of an azimuthal asymmetry and as a consequence the predicted east–west neutrino asymmetries are all of approximately the same size.
## 3 Effects of the geomagnetic field on the shower development
### 3.1 Bending of the trajectories of charged particles
The trajectories of all charged particles in cosmic rays showers are curved because of the presence of the geomagnetic field. It will be particulary important in this work to consider the effect of the bending on the zenith angle of the particles, that is to study if the particle trajectories are bent “upward” or “downward”. We can define a (position dependent) system of coordinates with the $`z`$ axis pointing up, the $`x`$ axis pointing toward magnetic north and the $`y`$ axis completing a right–handed system pointing toward the magnetic west. In this system the magnetic field has by definition components $`\stackrel{}{B}(\left|B_{\mathrm{hor}}\right|,0,B_{\mathrm{vert}})`$ with $`B_{\mathrm{vert}}`$ positive (negative) in the southern (northern) hemisphere. The equation of motion of a charged particle in the geomagnetic field for the vertical component is:
$$\frac{dp_z}{dt}=q\beta _y\left|B_{\mathrm{hor}}\right|=+q\beta _{\mathrm{east}}\left|B_{\mathrm{hor}}\right|,$$
(6)
where $`q`$ is the electric charge and $`\stackrel{}{\beta }`$ the velocity of the particle. Positively charged particles ($`q>0`$) traveling east ($`\beta _y<0`$) have $`dp_z/dt>0`$, and are bent up. Negatively charged particles with the same initial direction are are bent down. The reverse happens for particle traveling in the opposite (west–ward) direction ($`\beta _y>0`$): negatively charged particle are bent up while positively charged ones are bent down.
An example of the bending of charged particles is shown in fig. 1. The figure is in true scale and represents a projection of the space above the Super–Kamiokande detector. The center of the SK detector (latitude 36.42 N, longitude 137.31 E, and altitude 371.8 meters above sea level ) corresponds in the figuere to the point with coordinates (0,0.3718). The thick line represents the “sea level” surface of the Earth in the vicinity of SK detector. The $`z`$ axis corresponds to the vertical axis that passes through the SK center, and the horizontal axis points to geographical west. The thin solid lines describe the projections in the plane of the figure of the trajectories of three charged particles. The particles are part of a cosmic ray shower and are “connected” by a production chain $`p\pi ^+\mu ^+`$. The trajectories are calculated integrating the equations of motion for charged particles in the International Geomagnetic Reference Field (IGRF ) for the year 2000
and correspond to a 10 GeV proton, a 7 GeV $`\pi ^+`$ and a 5 GeV $`\mu ^+`$. The pion production and decay points are marked by two diamonds. The muon emits a neutrino (the trajectory and its continuation beyond is indicated by a dotted line) that intersects the SK detector. For illustration both the pion and the muon decay after exactly three proper lifetimes. Both the pion and the muon are emitted exactly collinearly with respect to the primary (or parent) particle, so the angle between the neutrino and the primary directions is entirely due to the effects of the geomagnetic field. Since all three particle are positive and travel toward the east direction they are bent upward.
In the following we will study the effects of the bending of the shower particles in the magnetic field on the fluxes of atmospheric neutrinos. These effects have been ignored in all previous calculations. In our numerical work we have included the bending of all charged particles in the shower. In the qualitative discussion that follows we will discuss separately the effect of the binding on the primary particles, and on the muons produced in the showers. These are in fact the most important sources of the effects we have found. Charged pions have a a lifetime that is two orders of magnitude shorter than $`\tau _\mu `$ (and a mass similar to $`m_\mu `$), therefore their bending in the magnetic field is much less important than for muons. Most of the neutrinos are produced in the chain decay of mesons produced in the first interaction of a primary. Also for those produced deeper in the shower most of the interesting effects is the result of the bending of the primary particles and the muons.
Note that we can expect that the effects of the field will be especially important for the neutrinos that come from muon decay, and significantly smaller for the muon (anti)–neutrinos produced directly in meson decay.
### 3.2 The bending of the primary particles
It is incorrect to think that all effects of the magnetic field on the primary particles are taken into account by the study of allowed and forbidden trajectories discussed in section 2. The bending of the primary particles is a factor in determining the average altitude of the primary interactions, and the zenith angle distribution of the primaries at the interaction point. These two effects are not taken into account in the existing calculations. Note that the zenith angle (that is the polar angle with respect the local vertical) of a particle can be calculated as:
$$\mathrm{cos}\theta =\frac{\stackrel{}{p}\stackrel{}{r}}{\left|\stackrel{}{p}\right|\left|\stackrel{}{r}\right|}$$
(7)
and is not a constant even for particles traveling along a straight line. The value of the zenith angle along a particle trajectory will of course be determined by the bending in the magnetic field.
As an illustration let us study the trajectories of two protons both on allowed trajectories that cross an imaginary spherical surface above the Earth with the same zenith angle but different azimuth angles, with one particle traveling toward east, and second one toward west. If the two trajectories are approximated as straight lines the distributions of altitude of the interaction point for the two particles (determined by the interaction length $`\lambda _{\mathrm{int}}`$, the inclination of the trajectory and the altitude profile of the air density) will of course be identical. It is however easy to see that since the east–going particle is bent “up”, and the west–going one is bent “down”, the first particle will on average interact interact higher, and more “horizontally”:
$$h_{\mathrm{int}}_E>h_{\mathrm{int}}_W,$$
(8)
$$\mathrm{cos}\theta _0^{\mathrm{int}}_E<\mathrm{cos}\theta _0^{\mathrm{int}}_W$$
(9)
Both the inequalities (8) and (9) result in an enhancement of the east–west effect for the neutrinos. The average number of neutrinos produced in a shower is larger when the primary particle interacts higher because the muons have a longer pathlength for decay, and the shower develops in a medium of lower density where the decays of mesons is enhanced, and muons lose less energy before decay. For the same reasons more inclined showers produce more neutrinos.
In summary we can expect that the bending of the primary particle inside the atmosphere result in an enhancement of the neutrino east–west asymmetry with respect to the estimates that only consider the bending in vacuum. This enhacement will be approximately equal for all four neutrino types, it is therefore an interesting effect, but it is cannot be the explanation for the (hint of) discrepancy between data and prediction represented by equations (2), (3) and (4).
### 3.3 The angle between the neutrino and the primary particle
For the effects that we want to consider it is essential to consider the difference in direction between a neutrino and the primary particle that generated the shower where it was produced. We introduce the notation: $`\mathrm{\Omega }_\nu =\mathrm{\Omega }_0\mathrm{\Omega }_{0\nu }`$ where $`\mathrm{\Omega }_\nu `$ is the direction of the detected neutrino, $`\mathrm{\Omega }_0`$ the direction of the primary particle, and $`\mathrm{\Omega }_{0\nu }`$ the angle between the neutrino and the primary. The angle $`\mathrm{\Omega }_{0\nu }`$ is the result of the combination of several processes. For the neutrinos that are produced directly in a meson decay we can write:
$$\mathrm{\Omega }_{0\nu }=\mathrm{\Omega }_{0\pi }\mathrm{\Omega }_{\pi \nu }$$
(10)
where $`\mathrm{\Omega }_{0\pi }`$ the direction of emission of the parent pion in the primary interaction, and $`\mathrm{\Omega }_{\pi \nu }`$ the direction of emission of the neutrino in the pion decay (this qualitative discussion obviously applies also for kaons). Similarly for neutrinos emitted in the chain decay $`\pi \mu \nu `$ we have:
$$\mathrm{\Omega }_{0\nu }=\mathrm{\Omega }_{0\pi }\mathrm{\Omega }_{\pi \mu }\mathrm{\Omega }_{\mu B}\mathrm{\Omega }_{\mu \nu }$$
(11)
where $`\mathrm{\Omega }_{\pi \mu }`$ ($`\mathrm{\Omega }_{\mu \nu }`$) is the emission direction of the muon (neutrino) in the decay of the pion (muon), and $`\mathrm{\Omega }_{\mu B}`$ is the deviation of the muon in the geomagnetic field (the deviations in the magnetic field of the other charged particles are typically two orders of magnitude smaller; they are included in our numerical work).
The average deviation (considered as a space angle) $`\alpha _j`$ for the process $`j`$ can be easily estimated in first order:
$$\alpha _{0\pi }\frac{p_{}_\pi }{p_\pi }\frac{4p_{}_\pi }{p_\nu }\frac{5.2^{}}{p_\nu \left(\mathrm{GeV}\right)},$$
(12)
where we have used the approximation $`p_\pi 4p_\nu `$ (since a charged pion energy is approximately shared equally between three neutrinos and a $`e^\pm `$) and have assumed a transverse momentum $`p_{}350`$ MeV;
$$\alpha _{\pi \nu }\frac{p^{}}{p_\nu }\frac{1.7^{}}{p_\nu \left(\mathrm{GeV}\right)}$$
(13)
where $`p^{}`$ is the center of mass momentum of the final state particles in the decay $`\pi \mu \nu `$;
$$\alpha _{\pi \mu }\frac{p^{}}{p_\mu }\frac{3p^{}}{p_\nu }\frac{0.6^{}}{p_\nu \left(\mathrm{GeV}\right)}$$
(14)
this is approximately 3 times smaller that (13) because the muons have on average three times the energy of their neutrino decay products;
$$\alpha _{\mu \nu }\frac{m_\mu /3}{p_\nu }\frac{2.0^{}}{p_\nu \left(\mathrm{GeV}\right)}$$
(15)
and for the deviation of the muons in the magnetic field:
$$\alpha _{\mu B}\frac{L_\mu }{R_\mu }\left(\tau _\mu \frac{p_\mu }{m_\mu }\right)\left(\frac{eB}{p_\mu }\right)10.7^{}B\left(\mathrm{Gauss}\right)$$
(16)
where $`L_\mu `$ is the muon decay length, $`R_\mu `$ the muon gyro–radius in the geomagnetic field, and $`B`$ the value of the field. In this estimate we have neglected the muon energy loss and we have also assumed that the muons decay before hitting the ground.
There are two fundamental differences between the effect of the magnetic bending and the other contributions to $`\mathrm{\Omega }_{0\nu }`$. The first difference is that all deviations, except the magnetic bending one scale $`p_\nu ^1`$ reflecting the forward boost in the interaction or decay of relativistic particles. The deviation of the muons in the geomagnetic field is (in first order) independent from the muon momentum, because both the magnetic rigidity and the decay pathlength increase proportionally to $`p_\mu `$. Higher momentum muon bend less in the field, but live longer and the field can act for a longer time.
The second and most important difference is that all deviations (except the bending one) are azimuthally symmetric, therefore if we consider the deviation in any plane, the average value vanishes. The deviation in the magnetic field is of course not an “average value”, all muon of the same initial momentum have exactly the same deviation, and the deviation is not azimuthally symmmetric but happens in a well defined plane. If we for example consider the average zenith angle of the neutrinos produced in the showers generated by primaries of fixed direction $`\mathrm{\Omega }_0(\mathrm{cos}\theta _0,\phi )`$ we find that the average contribution of all sources vanishes, except for the effect of the bending in the field.
$$\theta _\nu =\theta _0+\theta _{\mu B}$$
(17)
This effect plays a crucial role in determining the neutrino east–west asymmetry.
## 4 Magnetic bending of $`\mu ^\pm `$ and the $`\nu `$ east–west asymmetries
The effect of the geomagnetic field on the trajectory of charged particles (equation (6)) implies (a scheme of the results is also shown in fig. 2):
* Positively charged particle ($`q>0`$) traveling toward east ($`\beta _y<0`$) are bent “up”.
* Positively charged particle ($`q>0`$) traveling toward west ($`\beta _y>0`$) are bent “down”.
* Negatively charged particle ($`q>0`$) traveling toward east ($`\beta _y<0`$) are bent “down”.
* Negatively charged particle ($`q>0`$) traveling toward west ($`\beta _y>0`$) are bent “up”.
From equation (17) we can then deduce that neutrinos produced in the decay of positive (negative) muons in showers of east–ward (west–ward) traveling primaries will be on average more “horizontal” than the primary particle; conversely the neutrinos produced in the decay of positive (negative) muons in showers of west–ward (east–ward) traveling primaries will be on average more “vertical” than the primary particle.
The key remark is now, that this systematic difference between the zenith angle of the neutrino and the primary results for the cases (a) and (d) in an enhancement and for the cases (b) and (c) in a suppression of the neutrino flux.
Two effects play a role in the enhancement or suppression of the flux. The first effect is purely geometrical and would be present even for an exactly isotropic primary flux (this is not so academic as it may sound, because the primary flux is actually isotropic for all rigidities above 60 GV). In this case the rate of primary interactions with zenith angle $`\theta _0`$ is $`\mathrm{cos}\theta _0`$, and there is an excess of vertical showers over horizontal ones. Therefore if the zenith angle of the primary particles is smaller (larger) than the neutrino one we can expect a larger (smaller) event rate. This is illustated schematically in the diagram of fig. 2. The neutrinos detected by an observer (indicated by a diamond) located on the surface of the Earth (shown as a thick line) looking in the solid angle $`d\mathrm{\Omega }_{\mathrm{west}}`$ (indicated by a dashed cone) are produced in a “patch” (outlined with a thick line) of atmosphere (that in the diagram is represented by a thin spherical shell) subtended by the solid angle $`d\mathrm{\Omega }`$. It can be seen from the figure that this patch of atmosphere “looks” larger (it has a larger projected area) for the primary particles if the production chain is $`p\pi ^+\mu ^+\nu `$ (the neutrinos can be $`\nu _e`$ and $`\overline{\nu }_\mu `$) than in the case when the neutrinos are produced directly in meson decay. This is a simple consequence of the fact that the projected area of the element of atmosphere subtended by the solid angle $`d\mathrm{\Omega }`$ is proportional to $`\mathrm{cos}\theta _0`$. For a fixed neutrino zenith angle, the average value of the primary zenith angle is not a constant but depends on the production mechanism. In the case $`p\pi ^+\mu ^+\nu `$ the relation between the zenith angle of the primary and of the neutrinos is $`\theta _\nu <\theta _0`$, in the case of neutrinos directly from meson decay: $`\theta _\nu \theta _0`$. Similarly the patch of atmosphere “looks” smaller for the primaries when the production process is $`p\pi ^{}\mu ^{}\nu `$ (the neutrinos are $`\overline{\nu }_e`$ and $`\nu _\mu `$) since in this case the relation between the zenith angle of the neutrinos and the primary is: $`\theta _\nu >\theta _0`$. The bending in different directions of positive and negative muons results in one case in an enhancement, and in the other in a suppression of the neutrino fluxes. This discussion can be repeated for a solid angle $`d\mathrm{\Omega }_{\mathrm{east}}`$ in the opposite hemisphere, with the crucial difference that the production chains that are enhanced and suppressed are interchanged.
If now we consider the fact that the primary cosmic rays that reach the Earth’s atmosphere are not isotropic we find actually that the enhancement and suppression are stronger that what is obtained with the geometrical argument outlined above. This can be uderstood noting that the set of set of forbidden trajectories becomes larger when the zenith angle increases. This can be easily seen from the Störmer formula (5) where for fixed $`\lambda _M`$ and $`\phi `$ the rigidity cutoff grows monotonically with the zenith angle $`\theta `$.
In summary: the effect of the magnetic bending of the muons results in an enhancement for the flux of the neutrinos produced in the decay of $`\mu ^+`$ (that is $`\nu _e`$ and $`\overline{\nu }_\mu `$) when the particles travel toward east, and a suppression when the particles travel toward west. The net effect is a positive contribution to the east–west asymmmetry. The effect of magnetic bending is obviously opposite for the flux of neutrinos produced in the decay of $`\mu ^{}`$ ($`\overline{\nu }_e`$ and $`\nu _\mu `$). The flux is suppressed when the when the particles travel toward east, and enhanced when the particles travel toward west. The net effect is a negative contribution to the east–west asymmmetry.
The effect of the magnetic bending of the muons has to be combined with the asymmetry generated by the geomagnetic effects on the primary particles, that to a good approximation result in an equal contribution to the asymmetry of all neutrino types. Moreover, for the neutrinos produced directly in meson decay, the effects of the geomagnetic field during the development the shower are negligible, and for them the only significant source of asymmetry are the geomagnetic effects on the primary trajectory.
### 4.1 Qualitative predictions
We can finally put all results together and give some simple, non–trivial and testable predictions.
1. With increasing statistics, the Super–Kamiokande detector will detect east–west asymmetries that are not in agreement with predictions based on one–dimensional calculations.
2. The asymmetries for the four neutrino types are in the relation:
$$A_{\overline{\nu }_e}<A_{\nu _\mu }<A_\nu ^{1\mathrm{D}}<A_{\overline{\nu }_\mu }<A_{\nu _e}$$
(18)
where $`A_\nu ^{1\mathrm{D}}`$ is the asymmetry predicted in a one–dimensional approximation, that is approximately equal for all neutrino types.
3. The neutrino type with the largest asymmetry is the $`\nu _e`$. The majority of these neutrinos are produced in the chain decay $`p\pi ^+\mu ^+\nu _e`$, and for all these neutrinos the contributions of the geomagnetic effects on the primary particle and on the muon add to each other.
4. The east–west asymmetry for $`\overline{\nu }_\mu `$ is smaller that for $`A_{\nu _e}`$ but larger than the 1–D prediction. This is a consequence of the fact that approximately one half of the $`\overline{\nu }_\mu `$’s are produced in $`\mu ^+`$ and have a large asymmetry (approximately equal to $`A_{\nu _e}`$) and for these neutrinos the geomagnetic effects on the primary particles and on the muons add to each other. The other half of the $`\overline{\nu }_\mu `$’s is produced in pion decay, and will have an asymmetry close to what is predicted in the 1–D approximation.
5. Developing a similar argument we can conclude that the asymmetry for $`\nu _\mu `$’s will be smaller that the 1–D prediction, because for approximately half of these neutrinos (those produced in the chain decay $`p\pi ^{}\mu ^{}\nu _\mu `$) the effects of the bending of the muons gives a negative contribution to the asymmetry.
6. The asymmetry for $`\overline{\nu }_e`$ will be the smallest one since most of the particles are produced in the decay of negative muons, and for nearly all $`\overline{\nu }_e`$’s the effect of the bending of the muons must be subtracted from the asymmetry generated by geomagnetic effects on the primary particles. Numerically we actually will find a deficit of east–going $`\overline{\nu }_e`$’s, that is an asymmetry that has changed sign. The effect of the muon bending overcompensates for the effect of the primary flux.
7. Even if the detector is not capable to distinguish neutrinos from anti–neutrinos, the effect of the magnetic bending of the muons is still measurable. For $`e`$–like events the net effect is an enhancement of the asymmetry with respect to the 1–D expectation:
$$A_e>A_e^{1\mathrm{D}}$$
(19)
This can be simply deduced as a consequence of three steps: (i) the effect of muon magnetic bending on the asymmetry is positive for $`\nu _e`$’s and negative for $`\overline{\nu }_e`$’s; (ii) in the neutrino flux the $`\nu _e/\overline{\nu }_e`$ ratio is approximately unity (more accurately $`1.2`$, reflecting the $`\pi ^+/\pi ^{}`$ ratio in the final states of proton interactions); (iii) the cross section for neutrinos is larger than for anti–neutrinos.
8. With a very similar argument, we can estimate that the asymmetry of $`\mu `$–like events must be smaller than the 1–D calculation:
$$A_\mu <A_\mu ^{1\mathrm{D}}$$
(20)
The inequality is opposite with respect to the $`e`$–like case, since the effect of the muon magnetic bending on the asymmetry is negative for neutrinos (with the larger cross section) and positive for anti–neutrinos.
## 5 Numerical results from a full 3–D calculation
To go beyond this qualitative discussion we have performed a full three-dimensional montecarlo calculation of the atmospheric neutrino fluxes, including the bending in the geomagnetic field of all charged particles in the cosmic ray showers. This montecarlo calculation has been already described in , where it was used to illustrate effects of a three–dimensional calculation on the predicted zenith angle distibutions of $`e`$–like and $`\mu `$–like events. We refer to for more details on the calculation method, here we will only briefly summarize the main points. We have used a straightforward, direct method (see the appendix for a comparison with other approaches). The first step of the calculation is the generation of isotropic fluxes of cosmic rays (we used the results of for the energy spectrum and mass composition) at “the top of the atmosphere” (at a radius $`R=R_{}+80`$ Km). The geomagnetic effects on the primary flux are computed studying the past trajectories of each generated particle. Primary particles on forbidden trajectories are rejected. Primary particles on allowed trajectories are propagated in the magnetic field (the International Geomagnetic Reference Field for the year 2000 ). About 97% of the particles interact in the air (the remaining fraction only grazes the Earth and continues its travel). The air is described as having a spherically symmetric distribution $`\rho \left(r\right)`$ obtained from a fit to the US standard atmosphere. The final states of the hadronic interactions was generated using the model (and montecarlo algorithms) of Hillas . All secondary charged particles are propagated along curved trajectories in the magnetic field, and a shower is generated developing a “tree” with a standard technique. For the propagation of muons we have included the (crucially important) energy loss for ionization in the air (and the trajectory correctly takes into account the variation of the muon momentum). Energy loss was neglected for all other particles. Multiple scattering was neglected for all particles. When neutrinos are produced their trajectories (simple straight lines) are studied. Each neutrino either “misses” the Earth, or crosses its surface twice, first as down–going and then as up–going. At each intersection its position and direction of the neutrinos is recorded. To compute the azimuth of the neutrinos, for each point on the surface of the Earth we have defined a “local” system of coordinates as discussed in section 3: the $`z`$ axis points up, the $`x`$ axis toward magnetic “north”, and the $`y`$ axis toward magnetic “west”. In this system the magnetic field has components $`\left(\right|B_{\mathrm{hor}}|,0,B_{\mathrm{vert}}`$). The azimuth $`\phi `$ (following the SK convention) is defined as the direction where the particle is going: a particle with $`\phi =90^{}`$ ($`270^{}`$) is traveling from east toward west (from west toward east).
The neutrinos generated with these algorithms are distributed over the entire surface of the Earth, with important non–uniformities due to the effects of the magnetic field. We have collected the neutrinos in five different regions of equal area selected according to the magnetic latitude: an equatorial region, two intermediate and two polar regions. The results in the two (north and south) polar regions and the two intermediate regions are essentially undistinguishable, and here we present only the results for the northern regions.
Our results on the azimuth distributions are collected in 6 figures (from fig. 3 to fig. 8). In the figures we plot the azimuth distribution of the event rates, obtained after the convolution of the calculated fluxes with the model of the neutrino cross section from . We have selected neutrinos in the momentum interval: $`p_\nu \left(\mathrm{GeV}\right)=[0.5,3]`$) and the zenith angle region $`\mathrm{cos}\theta _\nu =[0.5,0.5]`$. No detection efficiency or experimental smearing has been included. For each of the three regions of the Earth that we have considered, we include two figures. The first figure contains (in four separate panels) the azimuth distributions of the four neutrino types. The second figure (in two separate panels) contains the azimuth distribution of $`e`$–like and $`\mu `$–like events (these distributions are obtained simply summing together the results for $`\nu _e`$ and $`\overline{\nu }_e`$ (or $`\nu _\mu `$ and $`\overline{\nu }_\mu `$). The scale of the $`y`$ axis is absolute for all figures.
We have also performed the montecarlo calculation three times.
* A first calculation (represented by thick histograms) was performed using the “full 3–D” algorithms descrivbed above, including the bending of all secondary charged particles in the showers.
* A second calculation (represented by thin histograms) was performed to reproduce the 1–D algorithms. The geomagnetic effects are calculated for the primary cosmic ray particles traveling outside the Earth’s atmosphere (exactly as in the previous case), but all particles travel along straight line trajectories for $`r<R_{}+80`$ Km.). Also all final state particles are collinear with the projectile (or parent) particle. This is achieved modeling the interactions and particle decays exactly as in the previous case, (including therefore transverse momentum), and performing as a last step a rotation of all the 3–momenta of the final state particles so that they become parallel to the projectile (for interactions) or parent (for decays) particle.
* A third calculation (represented by dashed histograms) was performed neglecting the geomagnetic effects on the primary flux (therefore considering exactly isotropic primary fluxes) and using the 1–D algorithms outlined in the previous point. With these approximations the azimuth angle distributions must be flat and independent from the detector position. The only non trivial result of the calculation is the absolute normalization of the different rates. Note that this normalization must also be independent from the geographycal region considered.
Inspecting the results of fig. 3 to fig. 8, we can make the following remarks.
1. Calculation (iii) (no geomagnetic effects) results as expected in flat, detector position independent azimuth distributions.
2. Calculation (ii) (1D with geomagnetic effects included for the primary) exhibits as expected an east–west asymmetry, that is of approximately the same size for all four neutrino types.
3. The taking into account of the effects of the bending in the magnetic field of secondary particles (and also of the primary in the atmosphere) in calculation (iii) is a non negligible correction.
4. Comparing the results of calculation (ii) and calculation (iii) one can see that the effects of the bending of secondary particles result in an enhancement of the east–west asymmetry for $`\nu _e`$, and $`\overline{\nu }_\mu `$, and a suppression of the asymmetry for $`\overline{\nu }_e`$ and $`\nu _\mu `$.
5. The asymmetry for the combination $`\left(\nu _e+\overline{\nu }_e\right)`$ is also enhanced, while the asymmetry for the combination $`\left(\nu _\mu +\overline{\nu }_\mu \right)`$ is suppressed.
In summary we can observe that the expectations of the qualitative discussion of the previous section have been confirmed in a preliminary, but detailed, quantitative calculation. The results are summarized in table 1, that gives the results of the asymmetries obtained with the calculation (ii) (or 1–D) and the calculation (iii) (or 3–D), in the different geographyical regions. For the definition of the magnetic latitude we have used as dipole axis the one that corresponds to the leading term in the multipole expansion of the IGRF field for the year 2000. With this definition the SK detector has a magnetic latitude of $`27.08^{}`$ ($`\mathrm{sin}\lambda _M^{\mathrm{SK}}=0.455`$) and is inside the “intermediate” region in the northern hemisphere.
Inspecting table 1, we can notice that as expected the east–west asymmetry depends on the magnetic latitude of the detector, and is strongest near the magnetic equator, where all geomagnetic effects are enhanced. In the 1–D calculation the asymmetries for the four neutrino types are approximately equal, with small differences that reflect the fact that the neutrino spectra produced by a primary of fixed energy are not identical for the four flavors. The full 3–D calculation clearly exhibits the enhancements and suppressions predicted qualitatively in section 4.
## 6 Summary and conclusions
In this work we have shown that the azimuth angle distributions of atmospheric neutrinos are shaped by the effects of the geomagnetic field on both the trajectories of the primary cosmic rays, and their showers in the Earth’s atmosphere.
In the existing (one–dimensional) calculations of the atmospheric neutrino fluxes the effects of the magnetic field on the shower development are ignored. In a 1–D framework the asymmetries of the $`\nu `$ fluxes reflect only the highest suppression of the primary cosmic ray flux when the primary particles travel towards west, and the predicted east–west asymmetry is approximately equal for all four neutrino types: ($`\nu _e`$, $`\overline{\nu }_e`$, $`\nu _\mu `$ and $`\overline{\nu }_\mu `$).
The situation changes when one includes the bending in the geomagnetic field of secondary particles (and in particular of muons) in the cosmic rays showers. Positively and negatively charged particles are curved in different directions and the resulting effect on the east–west asymmetry is different for different $`\nu `$ types. It is an enhancement for the neutrinos that are the product of $`\mu ^+`$ decay ($`\nu _e`$, and $`\overline{\nu }_\mu `$) and a suppression for the products of $`\mu ^{}`$ decay ($`\overline{\nu }_e`$, and $`\nu _\mu `$).
For detectors like Super–Kamiokande that cannot separate $`\nu `$’s and $`\overline{\nu }`$’s the detectable effect is an enhanced east–west asymmetry for the $`e`$–like events, and a suppressed asymmetry for the $`\mu `$–like events, as it is already observed by Super–Kamiokande. It is remarkable that the detctor is able to measure successfully subtle effects as those described here (and without the need of a prediction). In our view this is a confirmation of the high quality of the experimental data.
The (partial) failure of the existing calculations in the prediction of the east–west asymmetry is also a warning and an indication that the development of more accurate and detailed calculations of the atmospheric neutrino fluxes is needed. Several groups are working in this direction. We note that a correct prediction of the east–west effects is possible only with a fully three–dimensional calculation. In our view the success of the calculation presented here in reproducing the SK results is also a confirmation of the non–trivial geometrical effects connected with 3–D effects and the geometry of the neutrino source volume discussed in and in more detail in .
The understanding that the present calculations of atmospheric neutrinos do not describe correctly the east–west asymmetry can have some interesting consequences for the determination of the oscillation parameters from the atmospheric neutrino data (in SK and also other experiments, in particular Soudan). We note that the effects we have been discussing are obviously especially important for horizontal neutrinos, and that the effects have different sign for $`e`$–like and $`\mu `$–like events. The shape of the zenith angle distribution of the atmospheric neutrino fluxes, in particular close to the horizontal plane, and the relation between the muon and electron event rates are (at least in the opinion of this author) the most important problems in the calculation of the atmospheric neutrino fluxes. As an illustration, in the simplest (and favored) picture of $`\nu _\mu \nu _\tau `$ oscillations for values of $`\mathrm{\Delta }m^2`$ in the range indicated by the data and the range of $`E_\nu `$ characteristic of the contained events, the value of $`\mathrm{sin}^22\theta `$ can be reliably calculated comparing the “up–going” and “down–going rates”, with a relatively small systematic error and even for a poor (or even incorrect) estimate of $`\mathrm{\Delta }m^2`$. This is possible since the ratio of the rates is to a good approximation determined only the “asymptotic forms” of the oscillation probability ($`P_{\nu _\mu \nu _\tau }=0`$ or $`\mathrm{sin}^22\theta /2`$) where the value of $`\mathrm{\Delta }m^2`$ is absent. On the other hand the determination of $`\mathrm{\Delta }m^2`$ requires a detailed fit of the zenith angle distributions, with the region close to the horizontal playing a crucial role. The (nearly exact) up–down symmetry of the $`\nu `$–fluxes is of little help in this case. The comparison of the $`e`$–like and $`\mu `$–like rates remains as a powerful tool to estimate the suppression (or enhancement) of a flavor type, but of course as discussed in this work, for precise quantitative evaluations of the parameters one need to study with great care the relations between the fluxes of different neutrino types.
We postpone to a future work a quantitative estimate of the effects of a full three–dimensional calculation in the determination of the allowed region in the neutrino oscillation parameter space.
Other important sources of uncertainties in the determination of the allowed region for the oscillation parameters are the input primary spectra, where very valuable new data became recently available , and the modeling of hadronic interactions. New data on this problem would also be of great value . A detailed description of the neutrino cross section is also needed to compute the event rates, and new data would be very valuable and benefit also the long–baseline programs.
Acknowledgments Special thanks to Takaaki Kajita for very useful discussions and encouragement. I’m also grateful to Ed Kearns for discussions about the SK data, and gladly acknowledge discussions with Giuseppe Battistoni, Alfredo Ferrari, Tom Gaisser, and Yoichiro Suzuki. This work was also possible thanks to the precious help of Massimo Carboni, Kenji Kaneyuki and Atsushi Okada.
## Appendix: 3–D calculations of atmospheric neutrinos
The first calculations of the atmospheric neutrino fluxes have been performed in a one–dimensional approximation, that is assuming that the neutrinos are emitted collinearly with the primary particle. This allows an enormous reduction of the size of a montecarlo calculation, because only the (formally vanishingly small) fraction of the cosmic rays showers that has trajectories that intersect the detector has to be simulated.
The difficulty of a a full three–dimensional calculation of the neutrino fluxes has been discussed in and . In a 3–D calculation any shower can produce a neutrino that intersect the detector we are considering, and therefore all possible showers have to be studied. Since the atmospheric neutrinos are generated quasi–uniformally over the entire surface of the Earth with a total area $`5.1\times 10^8`$ Km<sup>2</sup>, only a very small fraction of these neutrinos will intersect a dector with an area of order $`10^3`$ Km<sup>2</sup>. Of course it is possible to consider a detector area greatly enlarged to improve the statistical precision of the calculation, however this cannot be done arbitrarily because in fact the $`\nu `$ flux is not exactly uniform, but it does depend on the detector location. This in a nutshell is the computational problem of a 3–D calculation.
There are several methods that have been used or are currently under study to overcome this difficulty.
### 6.1 Spherically symmetric problem
A useful first step is to consider a a simpler, sperically symmetric problem, that is obtained neglecting all effects of the geomagnetic field both for the determination of allowed and forbiddden trajectories, and in the development of the showers. This problem is explicitely spherically symmetric: all points on the surface of the Earth are equivalent, and the entire surface can be used as the “detector”. This calculation requires therefore the same computer power of a 1–D calculation, because essentially all produced neutrinos (all those that intersect the Earth’s surface) can be collected and analysed to “measure” the $`\nu `$–fluxes in a montecarlo calculation.
A calculation along these lines has been performed by G. Battistoni et al . This calculation for the first time demonstrated the existence of non–trivial geomatrical effects (see for more discussion) due to the spherical geometry of the neutrino source volume.
### 6.2 Shower translation algorithm
Unfortunately a calculation performed using the approximation of a spherically symmetric Earth is not adequate for the analysis of experimental data. The largest effect that is missing is the effect of the geomagnetic effect that “forbids” low rigidity cosmic rays from reaching the vicity of the Earth. If one neglects the effects of the geomagnetic field on the shower development it is possible to “translate” a shower from a position on the Earth to an arbitrary one, keeping the zenith angle of the primary particle (and therefore the zenith angle of all particles) as constant (this of course corresponds to a rotation of a shower seen as “rigid body” around the center of the Earth, plus an additional rotation around the new vertical axis). A single shower generated by montecarlo can then be “used” many times with enormous saving of computer time. The effects of the geomagnetic effects in the determination of the allowed and forbidden trajectories can be easily included, checking that if a given shower position corresponds to an allower of forbidden primary trajectory and rejecting forbidden trajectories. A preliminary calculation using this method for the three sites of Kamioka, Soudan and Gran Sasso has been recently made available on the web by Battistoni et al. .
### 6.3 Weighted technique
If the effects of the magnetic fields are also included in the development of the shower, the problem loses all symmetry and one has to confront the computational problem of a calculation that is intrinsically very inefficient. A solution of this problem that is under study , is to use a weighted technique, generating the cosmic ray showers with a strong bias in position and direction, so that the neutrinos produced are more likely to arrive in the vicinity of the detector in consideration, and using a weight system to estimate correctly the flux. In this method approximately half of the cosmic rays interact in a small region above the detector to simulate the down–going flux, while approximately one half are generated over the rest of the Earth’s surface. For these events the direction of the shower is chosen preferentially in a cone with an axis that “points” toward the detector. In this way it is possible to collect a sufficiently large statistics of neutrinos in a relatively small area around the detector site.
### 6.4 Direct approach
A possible solution is also to neglect any attempt at simplification, or optimization, and simply generate all cosmic showers on the Earth with a realistic distributions in position and direction, that only takes into account the geomagnetic effects studying the development of the showers in a realistic magnetic field. Such a calculation will produce a population of neutrinos that is distributed quasi–uniformally over the entire surface of the Earth, with the non–uniformities representing the expected variations in intensity due to geomagnetic effects, that produce a larger (smaller) flux in the magnetic polar (equatorial) region. In this direct (or brute force) approach the calculation is not performed for a given detector position, because all positions on the Earth are treated “democratically”. This is the technique used in and also in this work. Its only drawback is that the results that can be obtained in a reasonable time integrating then neutrino results over a large area of the Earth.
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# THE NEWTONIAN CORRESPONDENCE
(Dec. 2001)
It is shown that the field equations of general theory of relativity in the Einstein tensor form and the unimodular theory of gravity do not fulfill the correspondence principle commitment completely. The consistent formalisms are briefly discussed.
Key words: correspondence principle, general relativity, unimodular gravity.
1. INTRODUCTION
General relativity (GR) is established upon some physical principles. The principle of covariance (PC) and the correspondence principle (CP) are two of them which are considered in this letter. According to PC the field equations should have tensorial form, and GR must agree with the Newtonian gravitational theory in the limit of weak gravitational fields and low velocities by CP. Different aspects of the Newtonian limit may be classified as follows:
* \- The equation of geodesic deviation.
* \- The geodesic equation.
* \- The weak field limit of GR should give the same equations of motions as Newtonian gravity.
In GR we are dealing with second rank tensorial field equations, generally a set of ten relations, while in the Newtonian gravity we have only one Poisson equation and it seems there is no correspondence for nine of the rest. Does this mean that the PC breaks in taking the Newtonian limit? The answer is negative. In the Newtonian limit the Lorentz transformations reduce to Galileo transformations, so that $`t`$ appears as a scalar. By Newtonian correspondence we must consider weak fields and low velocities. It turns out that in the spatial components of the field equations the first non-zero term has an order of approximation higher than the corresponding one in the $`tt`$-component. Since in finding the Newtonian limit we merely keep the first order terms in the $`tt`$-component, this leads to $`0=0`$ for other components. For more clarification we may work in a system of units that $`c1`$. This explicitly shows , when the velocity of light tends to infinity, how some components of the field equation disappear. From this point of view we may say that PC is not violated but the other components have no physical information. So we may restate the item (c) as follows:
* \- The (00)-component of the field equation must reduce to the Poisson equation for a weak stationary field produced by nonrelativistic matter.
We are going to show that, from (ć) point of view, the GR field equations in the form of Einstein tensor and the field equations of the unimodular gravity do not satisfy CP. The consistent form and its consequences are discussed.
2. EINSTEIN TENSOR FORM
We restrict our discussion to the Schwarzschild space which is the solution of the field equations for spherically symmetric vacuum space around a point mass M. In the literature we have the Einstein field equations in the form Einstein tensor proportional to energy-momentum tensor i.e. :
$$R_{\mu \nu }\frac{1}{2}g_{\mu \nu }R=\frac{8\pi G}{c^4}T_{\mu \nu }$$
(1)
satisfying CP so that the (00)-component of the field equation reduces to the Poisson equation in the weak field limit. It will be shown that for the Schwarzschild metric this is not so.
The Weinberg’s argument to reach this result is based on the fact that in a nonrelativistic system $`T_{ij}T_{00}`$ , then $`|G_{ij}||G_{00}|`$ and $`R_{ij}\frac{1}{2}g_{ij}R`$. Furthermore $`g_{\alpha \beta }\eta _{\alpha \beta }`$ and the curvature scalar is given by
$$RR_{kk}R_{00}2R_{00}$$
(2)
So concludes that $`G_{00}R_{00}`$ and $`R_{00}^2g_{00}`$ . The weak point in this argument is that by making use of $`G_{ii}=0`$ in calculating $`G_{tt}`$ actually different components of the field equations are combined. In other words the Poisson equation is constructed by a proper mixing of all the available equations. This is in contrast with the original claim that the $`(tt)`$-component of the field equation in the weak field limit gives the Poisson equation.
The Einstein tensor in the weak field limit may yield to $`\mathrm{}\psi _{\mu \nu }`$ , where $`\psi _{\mu \nu }h_{\mu \nu }\frac{1}{2}\eta _{\mu \nu }h,h_{\mu \nu }=g_{\mu \nu }\eta _{\mu \nu }`$ in Minkowski coordinate , provided that the Einstein condition is satisfied as follows
$$h_{\nu ,\mu }^\mu \frac{1}{2}h_{,\nu }=0,h=\eta ^{\mu \nu }h_{\mu \nu }$$
(3)
If this condition holds , in the stationary case $`\mathrm{}`$ reduces to $`^2`$ and the Poisson equation is obtained automatically. Let us see what happens in the weak field limit of Schwarzschild metric. We have
$`h_{tt}={\displaystyle \frac{2\varphi }{c^2}},h_{x_ix_i}={\displaystyle \frac{2\varphi x_{i}^{}{}_{}{}^{2}}{r^2c^2}},h_{x_ix_j}={\displaystyle \frac{2\varphi x_ix_j}{r^2c^2}},i,j=1,2,3`$
$`\varphi ={\displaystyle \frac{GM}{r}}c^2.`$ (4)
Using (4) we get $`h=0`$ and these do not satisfy (3), i.e. Einstein condition does not hold in this case. It means that for Schwarzschild space we do not end to the Poisson equation in the weak field limit.
Since curvature tensor and its contractions are invariant quantities under a gauge transformation of $`h_{\mu \nu }`$ as follows
$`x^\mu x_{}^{^{}}{}_{}{}^{\mu }=x^\mu +ϵ\xi ^\mu `$
$`h_{\mu \nu }h_{}^{^{}}{}_{\mu \nu }{}^{}=h_{\mu \nu }2\xi _{(\mu ,\nu )}`$ (5)
it is possible to find a gauge in which Einstein condition holds. This gauge may be obtained from
$$\mathrm{}\xi _\nu =\psi _{}^{\mu }{}_{\nu ,\mu }{}^{},\psi _{}^{\mu }{}_{\nu ,\mu }{}^{}=h_{}^{\mu }{}_{\nu ,\mu }{}^{}\frac{1}{2}h_{,\nu }$$
(6)
We may conclude that the weak field limit of GR and Newtonian field equation are not in the same gauge.
In what follows we will see that this approximation although may lead to a correct prediction of reciprocal of distance for Newtonian potential but indeed does not reduce to the Poisson equation as is required.
The line element of a spherically symmetric vacuum space is
$$ds^2=B(r)c^2dt^2A(r)dr^2r^2(d\theta ^2+\mathrm{sin}^2\theta d\phi ^2)$$
(7)
where $`B(r)=A^1(r)=1+\frac{2\varphi }{c^2}`$. Using (7) the nonvanishing components of Einstein tensor are:
$$G_{rr}=\frac{B^{^{}}}{rB}+\frac{A1}{r^2}$$
(8)
$$G_{\theta \theta }=\frac{r^2B^{^{\prime \prime }}}{2AB}+\frac{r^2B^{^{}}}{AB}(\frac{A^{^{}}}{A}+\frac{B^{^{}}}{B})\frac{r}{2A}(\frac{A^{^{}}}{A}+\frac{B^{^{}}}{B})$$
(9)
$$G_{\phi \phi }=\mathrm{sin}^2\theta G_{\theta \theta }$$
(10)
$$G_{tt}=c^2[\frac{BA^{^{}}}{rA^2}+\frac{B}{r^2}(1+\frac{1}{A})]$$
(11)
prime stands for differentiation with respect to $`r`$.
In contrast to what is expected, $`G_{tt}`$ for the Schwarzschild metric merely contains the first order differentiation with respect to $`r`$ and in no way can yield to the Poisson equation in weak field limit. Therefore there is an obvious discripancy between the obtained result and the Newtonian equation. Although (11) in the limit of weak fields gives
$$G_{tt}2(\frac{\varphi ^{^{}}}{r}+\frac{\varphi }{r^2})$$
(12)
which has the same solution of reciprocal of $`r`$ as the Poisson equation possess for a particle with mass M. This can be considered as a gauge violation of CP which may be forbidden too.
3. UNIMODULAR GRAVITY
In a more plausible consideration of cosmological constant as an integration constant the unimodular gravity is actually very well motivated. If the determinant of $`g`$ is not dynamical then the action only has to be stationary with respect to variations in the metric for which $`g^{\mu \nu }\delta g_{\mu \nu }=0`$ , yielding the field equations
$$R^{\mu \nu }\frac{1}{4}g^{\mu \nu }R=\frac{8\pi G}{c^4}(T^{\mu \nu }\frac{1}{4}g^{\mu \nu }T_\alpha ^\alpha )$$
(13)
with $`T^{\mu \nu }`$ as conserved stress tensor of matter. The combination of this with Bianchi identities for the covariant derivative of the Einstein tensor gives a nontrivial consisting condition
$$\frac{1}{4}_\mu R=\frac{8\pi G}{c^4}(\frac{1}{4}_\mu T_\lambda ^\lambda )$$
(14)
Denoting the constant of integration by $`4\mathrm{\Lambda }`$ the Einstein field equations is recovered.
We also see that this form of field equations i.e. (13), regretfully does not satisfy the CP from (ć) point of view. For spherically symmetric vacuum space (7) the components of (13) are :
$$R_{rr}\frac{1}{4}g_{rr}R=\frac{B^{^{\prime \prime }}}{4B}\frac{B^{^{}}}{8B}(\frac{A^{^{}}}{A}+\frac{B^{^{}}}{B})+\frac{A1}{2r^2}\frac{A^{^{}}}{rA}$$
(15)
$$R_{\theta \theta }\frac{1}{4}g_{\theta \theta }R=\frac{r^2B^{^{\prime \prime }}}{4AB}+\frac{r^2B^{^{}}}{8AB}(\frac{A^{^{}}}{A}+\frac{B^{^{}}}{B})+\frac{1}{2}(\frac{1}{A}1)$$
(16)
$$R_{\phi \phi }\frac{1}{4}g_{\phi \phi }R=\mathrm{sin}^2\theta (R_{\theta \theta }\frac{1}{4}g_{\theta \theta }R)$$
(17)
$`R_{tt}{\displaystyle \frac{1}{4}}g_{tt}R=c^2[{\displaystyle \frac{B^{^{\prime \prime }}}{4A}}+{\displaystyle \frac{B^{^{}}}{8A}}({\displaystyle \frac{A^{^{}}}{A}}+{\displaystyle \frac{B^{^{}}}{B}}){\displaystyle \frac{B}{2rA}}({\displaystyle \frac{A^{^{}}}{A}}+{\displaystyle \frac{B^{^{}}}{B}})`$
$`{\displaystyle \frac{B}{2r^2}}(1{\displaystyle \frac{1}{A}})]`$ (18)
In the weak field limit for the (18) we get
$$R_{tt}\frac{1}{4}g_{tt}R=\frac{\varphi ^{^{\prime \prime }}}{2}+\frac{\varphi }{r^2}$$
(19)
Again for (19) we have reciprocal of $`r`$ as its solution but it is not the Poisson equation as is expected from CP.
4. CONSISTENT FORM
The ordinary field equations in the following form
$$R_{\mu \nu }=\frac{8\pi G}{c^4}(T_{\mu \nu }\frac{1}{2}g_{\mu \nu }T)$$
(20)
fulfill the CP requirement , that is the $`(00)`$\- component of (20) for weak field limit of Schwarzschild metric reduces to
$$\varphi ^{^{\prime \prime }}\frac{2\varphi ^{^{}}}{r}=4\pi GM\delta (\stackrel{}{r})$$
(21)
which in a compact form is exactly the Poisson equation
$$^2\varphi =4\pi GM\delta (\stackrel{}{r})$$
(22)
For a perfect fluid the (00)-component of the RHS of (20) in the weak field limit reduces to
$$4\pi G(\rho +3p/c^2)$$
(23)
which is equal to $`8\pi G\rho _t`$ where $`\rho _t`$ is the timelike convergence density . In the limit of slow motion , $`\rho p/c^2`$, and $`p/c^2`$ can be ignored so that $`\rho _t=\rho /2`$, and Eq.(22) gives
$$^2\varphi =4\pi G\rho $$
(24)
The reason why this discripancy has not been recognized is that in finding the Schwarzschild metric we usually solve $`R_{\mu \nu }=0`$ as field equation. We may conclude that the form of the Einstein field equations with cosmological constant consistent with CP is
$$R_{\mu \nu }=\frac{8\pi G}{c^4}(T_{\mu \nu }\frac{1}{2}g_{\mu \nu }T)\mathrm{\Lambda }g_{\mu \nu }$$
(25)
This field equation may be derived from standard actions by considering the density metric of weight $`+1`$ instead of the metric as dynamical variables which is defined as :
$$\stackrel{~}{g}_{\mu \nu }=\sqrt{g}g_{\mu \nu }$$
(26)
and we get
$$\delta I=d^4x\{\frac{c^4}{16\pi G}(R^{\mu \nu }+\mathrm{\Lambda }g^{\mu \nu })+\frac{1}{2}(T^{\mu \nu }\frac{1}{2}g^{\mu \nu }T)\}\delta \stackrel{~}{g}_{\mu \nu }$$
(27)
From (26) we have
$$\delta \stackrel{~}{g}_{\mu \nu }=\sqrt{g}\delta g_{\mu \nu }\frac{1}{2}\sqrt{g}g_{\mu \nu }g^{\alpha \beta }\delta g_{\alpha \beta }$$
(28)
Inserting (28) in (27) gives the ordinary variation of standard action with respect to the variation of the metric.
$$\delta I=d^4x\{\frac{c^4}{16\pi G}(R^{\mu \nu }\frac{1}{2}g^{\mu \nu }R+\mathrm{\Lambda }g^{\mu \nu })+\frac{1}{2}T^{\mu \nu }\}\sqrt{g}\delta g_{\mu \nu }$$
(29)
This procedure may be carried out in an elegant way by applying the Palatini approach based on the idea of treating the metric (the density metric) and the connection separately as dynamical variables which the variation with respect to the connection reveals that the connection is necessarily the metric connection.
It is evident from (29) that the common field equations (1) are obtained under the variations of $`\delta g_{\mu \nu }`$ with the condition that$`g0`$. While the consistent form (20) are resulted from (27) under the variations of $`\delta \stackrel{~}{g}_{\mu \nu }`$ without any condition.
5. REMARKS
Let us summarize the significant results.
* \- It is shown that how (ć) statement may be explicitly obtained from (c) statement in the mentioned CP classification without violating PC.
* \- Einstein field equations in the common form (1) of Einstein tensor proportional to the energy-momentum tensor do not fulfill the CP from (ć) point of view.
* \- The unimodular gravity field equations (13) do not satisfy the (ć) statement.
* \- The alternative field equations (20) which are mathematically equivalent to the Einstein common field equations (1) satisfy the CP commitments completely. This means that indeed these two forms are not physically equivalent. In Cartesian spatial coordinates the Poisson equation may be obtained from all the components of this form of field equations.
* \- The failure of unimodular model in this study ceases the interpretation of the cosmological constant as an integration constant, i.e. it is a universal constant of nature.
* \- Derivation of Eq.(1) from Lagrangian formalism (29) requires the constraint $`g0`$. Thus the resulted field equations are restricted and are not necessarily defined for the whole space.
* \- By taking the density metric tensors (26) as dynamical variables the obtained field equations from Lagrangian formalism (27) are free from any constraint and holds everywhere.
Accordingly, we should accept to carry out recasting of the GR field equations.
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# ■("SPhT-00/013"@"HD–THEP–00–09") Chiral symmetry restoration and axial vector renormalization for Wilson fermions
## 1 Introduction
Any realization of fermions on the lattice has to respect the constraints imposed by the Nielsen-Ninomiya theorem . Whereas Wilson fermions break chiral symmetry explicitly, Ginsparg-Wilson fermions have an exact chiral symmetry on the lattice that is generated by composite local lattice operators . In both cases, the continuum chiral flavor mixing symmetry and the anomaly are to be properly reproduced as the cutoff is removed.
In a recent paper it was shown that under very general conditions on the lattice Dirac operator, which, in particular, are satisfied both for Wilson and for Ginsparg-Wilson fermions, the axial anomaly is correctly generated in the continuum limit. The main conditions are gauge invariance, absense of doublers, and locality on the lattice in a more general sense. The origin of the anomaly is traced back to an irrelevant, local lattice operator in the axial vector Ward identity.
For Ginsparg-Wilson fermions, the composite operator in the chiral transformation which ensures an exact flavor mixing symmetry on the lattice, stays irrelevant under renormalization . As a consequence, the axial vector current does not require renormalization. On the other hand, for Wilson fermions it is not obvious, although widely believed, that the chiral symmetry becomes restored in the continuum limit. Below we give a short but strict proof of this assertion to all orders of perturbation theory, based on lattice power counting for massless theories . Although we explicitely refer to Wilson fermions, the result is as general as that for the anomaly generation given in . As we shall show, the only role played by the irrelevant, symmetry breaking operator in the flavor mixing axial vector Ward identity is to give rise to a unique multiplicative renormalization $`Z_j`$ of the axial vector current, ensuring that chiral symmetry is restored in the continuum limit. We compute the one-loop contribution to $`Z_j`$ as a function of the Wilson parameter. The result is largely independent of the particular lattice regularization of the current.
## 2 General framework
Although our general proof will be given for QED, it generalizes in an obvious way to non-abelian gauge theories with massless fermions.
### 2.1 Renormalized lattice QED
The action for renormalized QED is given by
$$S(A,\psi ,\overline{\psi })=S_W(U)+S_f(U,\psi ,\overline{\psi })+S_{gf}(A).$$
(1)
$`S_W(U)`$ is e.g. the Wilson plaquette action
$$S_W(U)=Z_A\frac{1}{2g^2}\underset{xa^4}{}\underset{\mu \nu =0}{\overset{3}{}}\left(1U(x;\mu )U(x+a\widehat{\mu };\nu )U(x+a\widehat{\nu };\mu )^1U(x;\nu )^1\right),$$
(2)
with $`g`$ the renormalized gauge coupling constant and $`U(x,\mu )=\mathrm{exp}(iagA_\mu (x))`$ $``$ $`U(1)`$. The fermion action is given by
$$S_f=a^4\underset{xa^4}{}Z_\psi \overline{\psi }(x)\left(D[U]+m_0\right)\psi (x),$$
(3)
with $`\psi `$ a 2-flavor Dirac spinor field and with $`D[U]`$ the Wilson Dirac operator,
$`D[U]\psi (x)`$ $`=`$ $`{\displaystyle \frac{1}{2a}}{\displaystyle \underset{\mu =0}{\overset{3}{}}}[(\gamma _\mu r)U(x;\mu )\psi (x+a\widehat{\mu })`$ (4)
$`(\gamma _\mu +r)U(xa\widehat{\mu };\mu )^1\psi (xa\widehat{\mu })+2r\psi (x)].`$
$`m_0`$ is the bare fermion mass, which for massless fields must be tuned to its critical value of $`O(g^2)`$. $`S_{gf}`$ denotes the gauge fixing action. For concreteness we choose the Lorentz gauge,
$$S_{gf}(A)=a^4\underset{xa^4}{}\frac{\lambda }{2}\left(\underset{\mu =0}{\overset{3}{}}\frac{1}{a}\widehat{}_\mu ^{}A_\mu (x)\right)^2,$$
(5)
with $`\lambda >0`$ the gauge fixing parameter. Here and in the following, $`\widehat{}_\mu `$ and $`\widehat{}_\mu ^{}`$ denote the forward and backward lattice difference operators, respectively,
$$\widehat{}_\mu f(x)=f(x+a\widehat{\mu })f(x),\widehat{}_\mu ^{}f(x)=f(x)f(xa\widehat{\mu }),$$
(6)
where $`\widehat{\mu }`$ is the unit vector in $`\mu `$ direction.
The generating functional $`W`$ of the connected correlation functions is given by
$`\mathrm{exp}W(J,\eta ,\overline{\eta })={\displaystyle \underset{x}{}\left(d\psi (x)d\overline{\psi }(x)\underset{\mu }{}dA_\mu (x)\right)}`$
$`\mathrm{exp}\left\{S(A,\psi ,\overline{\psi })+S_c(A,\psi ,\overline{\psi };J,\eta ,\overline{\eta })\right\}`$ (7)
with source term $`S_c`$
$`S_c(A,\psi ,\overline{\psi };J,\eta ,\overline{\eta })=a^4{\displaystyle \underset{x}{}}\left\{{\displaystyle \underset{\mu }{}}J_\mu (x)A_\mu (x)+\overline{\eta }(x)\psi (x)+\overline{\psi }(x)\eta (x)\right\}.`$ (8)
The vertex functional $`\mathrm{\Gamma }`$ is obtained by a Legendre transformation
$$W(J,\eta ,\overline{\eta })=\mathrm{\Gamma }(𝒜,\psi ,\overline{\psi })+a^4\underset{x}{}\left(\underset{\mu }{}J_\mu (x)𝒜_\mu (x)+\overline{\eta }(x)\psi (x)+\overline{\psi }(x)\eta (x)\right),$$
(9)
where
$$a^4𝒜_\mu (x)=\frac{W}{J_\mu (x)},a^4\psi (x)=\frac{W}{\overline{\eta }(x)},a^4\overline{\psi }(x)=\frac{W}{\eta (x)}.$$
(10)
By $`\stackrel{~}{\mathrm{\Gamma }}^{(n,m)}(k,l)`$ we denote the momentum space vertex function of $`n`$ fermion pairs and $`m`$ gauge fields, with their collected momenta denoted by $`k`$ and $`l`$, respectively. For $`Q`$ any composite local lattice operator, we write $`\stackrel{~}{\mathrm{\Gamma }}_Q^{(n,m)}(q;k,l)`$ for the vertex function with one insertion of $`Q`$, with $`q`$ its momentum. Momentum conservation is implied. Massless fermions require that
$$\mathrm{tr}\stackrel{~}{\mathrm{\Gamma }}^{(1,0)}(k=0)=\mathrm{\hspace{0.33em}0}$$
(11)
to be achieved by tuning $`m_0`$, where the trace is taken in spinor space. $`Z_A`$ and $`Z_\psi `$ are uniquely determined by appropriate normalization conditions at non-exceptional momenta, e.g. by
$$.\frac{i}{4}\frac{}{\stackrel{~}{k}_0}\mathrm{tr}\gamma _0\stackrel{~}{\mathrm{\Gamma }}^{(1,0)}(k)|_{\overline{k}}=\mathrm{\hspace{0.33em}1},.\frac{1}{2}\frac{}{\widehat{k}_0}\stackrel{~}{\mathrm{\Gamma }}_{11}^{(0,2)}(k)|_{\overline{k}}=\widehat{\overline{\mu }},$$
(12)
where $`\overline{k}=(\overline{\mu }0,0,0,0)`$, $`\widehat{k}=(2/a)\mathrm{sin}(ka/2)`$, $`\stackrel{~}{k}=(1/a)\mathrm{sin}(ka)`$.
### 2.2 Symmetries
Below we make explicit reference to the charge conjugation symmetry
$$\mathrm{\Gamma }(𝒜^C,\psi ^C,\overline{\psi }^C)=\mathrm{\Gamma }(𝒜,\psi ,\overline{\psi }),$$
(13)
where
$$𝒜_\mu ^C(x)=𝒜_\mu (x),\psi ^C(x)=C\overline{\psi }(x)^T,\overline{\psi }^C(x)=\psi (x)^TC^1.$$
(14)
The superscript $`T`$ denotes transposition and $`C`$ the charge conjugation matrix satisfying
$$C^1\gamma _\mu C=\gamma _\mu ^T,\mu =0,\mathrm{},3.$$
(15)
Furthermore, applying a gauge transformation leads to the local Ward identity
$`i{\displaystyle \underset{\mu =0}{\overset{3}{}}}{\displaystyle \frac{1}{a}}\widehat{}_\mu ^{}{\displaystyle \frac{\mathrm{\Gamma }}{𝒜_\mu (x)}}+\left[g\overline{\psi }(x){\displaystyle \frac{\mathrm{\Gamma }}{\overline{\psi }(x)}}g\psi (x){\displaystyle \frac{\mathrm{\Gamma }}{\psi (x)}}\right]`$
$`i\lambda a{\displaystyle \underset{\mu ,\nu =0}{\overset{3}{}}}\widehat{}_\nu ^{}\widehat{}_\nu \widehat{}_\mu ^{}𝒜_\mu (x)=\mathrm{\hspace{0.33em}0}.`$ (16)
It implies that the renormalized action is of the form as stated above.
## 3 Chiral symmetry breaking and symmetry restoration
Chiral symmetry is broken by the Wilson Dirac operator. Under a local, flavor mixing chiral transformation
$$\delta \psi (x)=iϵ(x)\sigma _\alpha \gamma _5\psi (x),\delta \overline{\psi }(x)=iϵ(x)\overline{\psi }(x)\gamma _5\sigma _\alpha ,$$
(17)
where $`\sigma _\alpha `$, $`\alpha =1,2,3`$, denote the Pauli matrices acting in flavour space, the action transforms according to
$$\delta S=a^4\underset{x}{}iϵ(x)\left(\underset{\mu }{}\frac{1}{a}\widehat{}_\mu ^{}j_{\mu \alpha }(x)+\mathrm{\Delta }_\alpha (x)+\mathrm{\hspace{0.33em}2}m_0𝒫_\alpha (x)\right),$$
(18)
with
$$𝒫_\alpha (x)=Z_\psi \overline{\psi }(x)\gamma _5\sigma _\alpha \psi (x).$$
(19)
The gauge invariant local operators $`j_\mu `$ and $`\mathrm{\Delta }`$ are not uniquely determined by (18). In general, $`\mathrm{\Delta }`$ is a local lattice operator which is classically irrelevant, that is,
$$\underset{a0}{lim}\mathrm{\Delta }_\alpha (x)=\mathrm{\hspace{0.33em}0}.$$
(20)
It has UV degree 4 and IR degree 5. A convenient representation of $`j_{\mu \alpha }`$ and of $`\mathrm{\Delta }_\alpha `$ is given by
$`j_{\mu \alpha }(x)`$ $`=`$ $`Z_\psi {\displaystyle \frac{1}{2}}(\overline{\psi }(x)(\gamma _\mu +s)\gamma _5\sigma _\alpha U(x;\mu )\psi (x+a\widehat{\mu })`$
$`+\overline{\psi }(x+a\widehat{\mu })(\gamma _\mu s)\gamma _5\sigma _\alpha U(x;\mu )^1\psi (x)),`$
$`\mathrm{\Delta }_\alpha (x)`$ $`=`$ $`Z_\psi {\displaystyle \frac{1}{2a}}{\displaystyle \underset{\mu =0}{\overset{3}{}}}\{(rs)\overline{\psi }(x)\gamma _5\sigma _\alpha `$
$`\left[U(x;\mu )\psi (x+a\widehat{\mu })+U(xa\widehat{\mu };\mu )^1\psi (xa\widehat{\mu })2\psi (x)\right]`$
$`+(r+s)\left[\overline{\psi }(x+a\widehat{\mu })U(x;\mu )^1+\overline{\psi }(xa\widehat{\mu })U(xa\widehat{\mu };\mu )2\overline{\psi }(x)\right]`$
$`\gamma _5\sigma _\alpha \psi (x)\},`$
where $`r`$ is the Wilson parameter, and $`s`$ some arbitrary real but otherwise fixed constant. The following discussion on renormalization does not depend on a particular choice of s.
We add to the source part of the action $`S_c`$ a term
$$a^4\underset{x}{}\underset{\alpha =1}{\overset{3}{}}\left(\underset{\mu =0}{\overset{3}{}}G_{\mu \alpha }(x)j_{\mu \alpha }(x)+F_\alpha (x)\left[\mathrm{\Delta }_\alpha (x)+2m_0𝒫_\alpha (x)\right]\right),$$
(22)
and denote the corresponding vertex functional by $`\mathrm{\Gamma }^{}(𝒜,\psi ,\overline{\psi };G,F)`$. Obviously, $`\mathrm{\Gamma }^{}(𝒜,\psi ,\overline{\psi };G=0,F=0)`$ $`=`$ $`\mathrm{\Gamma }(A,\psi ,\overline{\psi })`$. Then (18) implies that $`\mathrm{\Gamma }^{}`$ satisfies the axial vector current Ward identity
$`{\displaystyle \underset{\mu }{}}{\displaystyle \frac{1}{a}}\widehat{}_\mu ^{}{\displaystyle \frac{\mathrm{\Gamma }^{}}{a^4G_{\mu \alpha }(x)}}+\left\{{\displaystyle \frac{\mathrm{\Gamma }^{}}{a^4\psi (x)}}\sigma _\alpha \gamma _5\psi (x)\overline{\psi }(x)\sigma _\alpha \gamma _5{\displaystyle \frac{\mathrm{\Gamma }^{}}{a^4\overline{\psi }(x)}}\right\}`$ (23)
$`={\displaystyle \frac{\mathrm{\Gamma }^{}}{a^4F_\alpha (x)}}+O(F,G).`$
The functional identity (23) is equivalent to the infinite set of momentum space Ward identities
$$i\underset{\mu =0}{\overset{3}{}}\widehat{q}_\mu \stackrel{~}{\mathrm{\Gamma }}_{j_{\mu \alpha }}^{(n,m)}(q;k,l)\stackrel{~}{\mathrm{\Gamma }}_{QED}^{(n,m)}(k,l)=\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{\Delta }_\alpha }^{(n,m)}(q;k,l)+2m_0\stackrel{~}{\mathrm{\Gamma }}_{𝒫_\alpha }^{(n,m)}(q;k,l).$$
(24)
Here we have written $`\stackrel{~}{\mathrm{\Gamma }}_{QED}^{(n,m)}(k,l)`$ for the pure QED part, which is a linear combination of $`\stackrel{~}{\mathrm{\Gamma }}^{(n,m)}(k,l)`$ with $`\gamma _5\sigma _\alpha `$ attached to the various external fermion lines, but with no composite operator inserted. According to the renormalization prescription of QED, it is UV finite and universal in the continuum limit.
QED is already renormalized, but because of $`\mathrm{\Delta }0`$, the axial vector current $`j_\mu `$ requires additional renormalization. This renormalization is multiplicative. That is, there exists a renormalization constant $`Z_j`$ such that
$$\stackrel{~}{\mathrm{\Gamma }}_{j_{\mu \alpha }R}^{(n,m)}(q;k,l)=Z_j\stackrel{~}{\mathrm{\Gamma }}_{j_{\mu \alpha }}^{(n,m)}(q;k,l)$$
(25)
is finite in the continuum limit, for all $`n`$ and $`m`$. The renormalized current satisfies the Ward identities
$$i\underset{\mu =0}{\overset{3}{}}\widehat{q}_\mu \stackrel{~}{\mathrm{\Gamma }}_{j_{\mu \alpha }R}^{(n,m)}(q;k,l)\stackrel{~}{\mathrm{\Gamma }}_{QED}^{(n,m)}(k,l)=\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{\Delta }_\alpha R}^{(n,m)}(q;k,l),$$
(26)
where
$`\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{\Delta }_\alpha R}^{(n,m)}(q;k,l)=\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{\Delta }_\alpha }^{(n,m)}(q;k,l)`$
$`+\left[\stackrel{~}{\mathrm{\Gamma }}_{𝒫_\alpha }^{(n,m)}(q;k,l)+i\left(Z_j1\right){\displaystyle \underset{\mu =0}{\overset{3}{}}}\widehat{q}_\mu \stackrel{~}{\mathrm{\Gamma }}_{j_{\mu \alpha }}^{(n,m)}(q;k,l)\right]`$ (27)
are the renormalized vertex functions with one $`\mathrm{\Delta }_\alpha `$-insertion. Because of
$$m_0=O(g^2)\text{and}Z_j1=O(g^2),$$
(28)
the part in brackets on the right hand side of (3) is equivalently obtained by adding local counter terms to the lattice source action $`S_c`$, that is, in (22), the square bracket is replaced by
$$\mathrm{\Delta }_\alpha (x)+2m_0𝒫_\alpha (x)+(Z_j1)\frac{1}{a}\underset{\mu =0}{\overset{3}{}}\widehat{_\mu ^{}}j_{\mu \alpha }(x).$$
(29)
We now show that for the particular choice of $`Z_j`$, these counter terms provide precisely overall Taylor subtractions at zero momentum, for all correlation functions with one $`\mathrm{\Delta }`$-insertion, according to their overall ultaviolet lattice divergence degrees. Together with (12), because $`\mathrm{\Delta }`$ is an irrelevant local lattice operator, this then implies that
$$\underset{a0}{lim}\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{\Delta }_\alpha R}^{(n,m)}(q;k,l)=\mathrm{\hspace{0.33em}0},$$
(30)
to all orders of perturbation theory, and for all $`n`$ and $`m`$ . The renormalized axial vector current becomes conserved in the continuum limit.
For the proof of this assertion, we recall that $`\mathrm{\Delta }`$ is a local operator of IR degree 5. This implies that $`\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{\Delta }R}^{(1,1)}`$ is continuous at zero momentum and $`\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{\Delta }R}^{(1,0)}`$ is once continuously differentiable at zero momentum. (These are the only vertex functions with one $`\mathrm{\Delta }`$-insertion that require overall UV subtractions, with overall (lattice) divergence degrees $`0`$ and $`1`$, respectively.)
First, charge conjugation symmetry implies that
$$\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{\Delta }_\alpha R}^{(1,1)}(0)=\mathrm{\hspace{0.33em}0}.$$
(31)
Furthermore, for the massless theory, satisfying (11), we obtain from (26) with $`n=1`$, $`m=0`$
$$\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{\Delta }_\alpha R}^{(1,0)}(0)=\mathrm{\hspace{0.33em}0},$$
(32)
because $`\stackrel{~}{\mathrm{\Gamma }}_{j_{\mu \alpha }R}^{(1,0)}(q;k,l)`$ is at most logarithmically infrared divergent if all momenta are sent to zero. Again, using charge conjugation symmetry, we know that for small momenta
$$\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{\Delta }_\alpha R}^{(1,0)}(q;k)=i\underset{\mu =0}{\overset{3}{}}\gamma _\mu q_\mu \gamma _5\sigma _\alpha +o(q,k)\text{as}q,k0,$$
(33)
with (infrared) finite constant $`z_\mathrm{\Delta }`$. Hence, order by order, $`Z_j`$ is uniquely determined by the requirement that
$$\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{\Delta }_\alpha R}^{(1,0)}(q;k)=o(q,k)\text{as}q,k0.$$
(34)
This completes the proof.
## 4 Axial Vector Renormalization Constant in One-Loop order
To one-loop order the renormalization constant $`Z_j`$ is determined from the condition (34), where $`\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{\Delta }R}^{(1,0)}`$ is given according to (3) by
$$\left[\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{\Delta }_\alpha R}^{(1,0)}(q;k)\right]_{1loop}=\left[\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{\Delta }_\alpha }^{(1,0)}(q;k)\right]_{1loop}+\mathrm{\hspace{0.33em}2}m_0\gamma _5\sigma _\alpha +i(Z_j1)\underset{\mu =0}{\overset{3}{}}\widehat{q}_\mu \gamma _\mu \gamma _5\sigma _\alpha .$$
(35)
The lattice Feynman diagrams that contribute to the first term on the right hand side are listed below. The vertices that correspond to the $`\mathrm{\Delta }`$-insertion are given in Appendix B.
(36)
To implement condition (34) we expand the right hand side of (35) around vanishing momentum up to first order. After a tedious but straight forward calculation one finds that
$$T_1\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{\Delta }_\alpha R}^{(1,0)}(q;k)=i\left[\zeta _\mathrm{\Delta }+\left(Z_j1\right)\right]\underset{\mu =0}{\overset{3}{}}q_\mu \gamma _\mu \gamma _5\sigma _\alpha ,$$
(37)
where $`T_1`$ denotes the Taylor expansion to order $`1`$ at zero momentum. $`\zeta _\mathrm{\Delta }`$ is given by the following (r-dependent) expression.
$$\zeta _\mathrm{\Delta }=Cg^2r^2_\pi ^\pi \frac{d^4\mathrm{}}{(2\pi )^4}\frac{h_\sigma (\mathrm{})}{(\stackrel{~}{\mathrm{}}^2+\widehat{M}_r^2)^2},$$
(38)
with $`C=1`$ for $`U(1)`$ and $`C=(N^21)/(2N)`$ for $`SU(N)`$, and
$`h_\sigma (\mathrm{})=\mathrm{cos}\mathrm{}_\sigma \left(\stackrel{~}{\mathrm{}}^2\stackrel{~}{\mathrm{}}_\sigma ^2+{\displaystyle \frac{\widehat{\mathrm{}}^{\mathrm{\hspace{0.17em}2}}}{2}}\eta _\sigma (\mathrm{})\right)\stackrel{~}{\mathrm{}}_\sigma ^2\eta _\sigma (\mathrm{})`$
$`+{\displaystyle \frac{\stackrel{~}{\mathrm{}}^2}{2}}+{\displaystyle \frac{1}{4}}r^2\widehat{\mathrm{}}^{\mathrm{\hspace{0.17em}2}}\left(\widehat{\mathrm{}}^{\mathrm{\hspace{0.17em}2}}\mathrm{cos}^2{\displaystyle \frac{\mathrm{}_\sigma }{2}}\stackrel{~}{\mathrm{}}_\sigma ^2\right),`$ (39)
where
$`\widehat{\mathrm{}}_\mu =\mathrm{\hspace{0.33em}2}\mathrm{sin}{\displaystyle \frac{\mathrm{}_\mu }{2}},\widehat{\mathrm{}}^{\mathrm{\hspace{0.17em}2}}={\displaystyle \underset{\mu =0}{\overset{3}{}}}\widehat{\mathrm{}}_\mu ^{\mathrm{\hspace{0.17em}2}},\stackrel{~}{\mathrm{}}_\mu =\mathrm{sin}\mathrm{}_\mu ,\stackrel{~}{\mathrm{}}^2={\displaystyle \underset{\mu =0}{\overset{3}{}}}\stackrel{~}{\mathrm{}}_\mu ^2,`$
$`\widehat{M}_r={\displaystyle \frac{r}{2}}\widehat{\mathrm{}}^{\mathrm{\hspace{0.17em}2}},\eta _\sigma (\mathrm{})=\mathrm{\hspace{0.33em}2}\mathrm{cos}^2{\displaystyle \frac{\mathrm{}_\sigma }{2}}{\displaystyle \underset{\mu =0}{\overset{3}{}}}\mathrm{cos}^2{\displaystyle \frac{\mathrm{}_\mu }{2}}.`$ (40)
$`\sigma `$ is any one of the indices $`0,\mathrm{},3`$. Condition (34) now determines $`Z_j`$ to be
$$Z_j=\mathrm{\hspace{0.33em}1}+\zeta _\mathrm{\Delta }.$$
(41)
$`\zeta _\mathrm{\Delta }`$ does not depend on a particular realization of the axial vector current $`j_{\mu \alpha }`$, Eqn. (3), that is, it is independent of the parameter $`s`$.
Note that the (finite) renormalization constant of the axial vector current is solely determined from diagrams involving the insertion of a classically irrelevant operator. The integral (38) is evaluated numerically. The dependence of $`\zeta _\mathrm{\Delta }`$ on the Wilson parameter is shown in Fig. 1. From this figure we see that $`\zeta _\mathrm{\Delta }`$ is a monotonically decreasing function of $`r>0`$. There exists no non-zero value of $`r`$ for which the axial current is not renormalized. In particular, for the commonly used value $`r=1`$ we obtain $`\zeta _\mathrm{\Delta }=0.0549Cg^2`$.
## 5 Conclusion
In this paper we have investigated the renormalization of the flavour mixing axial vector current for massless gauge theories with Wilson fermions. The corresponding axial vector Ward identity involves a symmetry breaking lattice operator $`\mathrm{\Delta }`$ which is local and classically irrelvant. Using the lattice power counting theorem for massless field theories, we have shown that $`\mathrm{\Delta }`$ uniquely determines the renormalization constant of the axial vector current in such a way that the chiral symmetry becomes restored in the continuum limit, to all orders of perturbation theory.
We have computed the renormalization constant to one-loop order. It is largely independent of a particular lattice realization of the current and non-vanishing whenever the Wilson parameter $`r0`$.
Although we have considered Wilson fermions, the result of symmetry restoration in the continuum limit is quite general. It holds for any lattice Dirac operator that satisfies a general set of conditions. These conditions are gauge invariance and charge conjugation symmetry, absense of doublers, and locality in the more general sense as stated in .
## Appendix A Small momentum behavior
The infrared properties of $`\mathrm{\Delta }`$ ensure that the vertex functions $`\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{\Delta }R}^{(1,1)}(q;k,l)`$ and $`\stackrel{~}{\mathrm{\Gamma }}_{\mathrm{\Delta }R}^{(1,0)}(q;k)`$ are continuous and once continuously differentiable at zero momentum, respectively. Their regular parts are obtained from the small momentum behavior of that part of the vertex functional $`\mathrm{\Gamma }^{}`$ that is linear in the source $`F`$. In momentum space it reads
$`{\displaystyle }{\displaystyle \frac{d^4q}{(2\pi )^4}}{\displaystyle \frac{d^4k_1}{(2\pi )^4}}{\displaystyle \frac{d^4k_2}{(2\pi )^4}}{\displaystyle \underset{\mu \alpha }{}}\stackrel{~}{F}_{\mu \alpha }(q)\{(2\pi )^4\delta (q+k_1+k_2)`$
$`\stackrel{~}{\overline{\psi }}(k_1)\left[i(\rho _rk_2+\rho _lk_1)_\mu \gamma _\mu +\eta \right]\gamma _5\sigma _\alpha \stackrel{~}{\psi }(k_2)`$ (42)
$`+\xi \stackrel{~}{\overline{\psi }}(k_1)\gamma _\mu A_\mu (qk_1k_2)\gamma _5\sigma _\alpha \stackrel{~}{\psi }(k_2)\}`$ (43)
with c-numbers $`\rho _r,\rho _l,\eta `$ and $`\xi `$. Applying the transformation (14) yields the same expression with
$$\rho _r\rho _l,\eta \eta ,\xi \xi .$$
(44)
The symmetry (13) thus implies that $`\rho _r=\rho _l`$ and $`\xi =0`$. The vanishing of $`\eta `$ is implied by the chiral Ward identity. This implies the statements (31)-(33).
## Appendix B Feynman rules
We state the Feynman rules for the insertion of one $`\mathrm{\Delta }_\alpha `$-operator, Eqn. (3), that are required for the computation of the axial vector current renormalization constant $`Z_j`$ to one-loop order. For simplicity the rules are given for gauge group U(1).
(45)
(46)
(47)
where
$$c_{k,\mu }=\mathrm{cos}\frac{k_\mu a}{2},\widehat{k}_\mu =\frac{2}{a}\mathrm{sin}\frac{k_\mu a}{2}.$$
(48)
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# References
KANAZAWA-00-02
March, 2000
Radiative symmetry breaking and Higgs mass bound in the NMSSM
Y. Daikoku <sup>1</sup><sup>1</sup>1e-mail: daikoku@hep.s.kanazawa-u.ac.jp and D. Suematsu <sup>2</sup><sup>2</sup>2e-mail: suematsu@hep.s.kanazawa-u.ac.jp
Institute for Theoretical Physics, Faculty of Science, Kanazawa University,
Kanazawa 920-1192, Japan
We study the upper mass bound of the lightest neutral Higgs scalar in the NMSSM using the RGE analysis. We require the successful occurence of the electroweak radiative symmetry breaking to restrict the parameter space. As a result the upper mass bound $`m_{h^0}`$ is largely restricted compared with the one estimated without imposing this condition. We point out some features of $`m_{h^0}`$ related to the initial value of $`h_t`$ and discuss why the models with more extra matters $`\mathrm{𝟓}+\overline{\mathrm{𝟓}}`$ of SU(5) could bring the larger maximum value of $`m_{h^0}`$.
The existence of a rather light CP-even neutral Higgs scalar is a common feature of the supersymmetric extension of the standard model . It is an important evidence of this kind of models and it is very useful to know its possible upper bound for the judgement of the consistency of the models.
The next to the minimal supersymmetric standard model (NMSSM) is the simplest extension of the minimal supersymmetric standard model (MSSM) . In this model a singlet chiral supermultiplet $`S`$ is introduced and a $`\mu `$ term in the MSSM is replaced by a Yukawa coupling $`\lambda SH_1H_2`$ with the usual Higgs doublet chiral superfields $`H_1`$ and $`H_2`$. The superpotential of the Higgs sector in this model is expressed as
$$W_{\mathrm{NMSSM}}=\lambda SH_1H_2+\frac{1}{3}\kappa S^3+\mathrm{}.$$
(1)
If the scalar component $`\stackrel{~}{S}`$ of $`S`$ gets a vacuum expectation value, the $`\mu `$ scale appears as $`\lambda \stackrel{~}{S}`$. In this model the $`\mu `$ problem in the MSSM is potentially solvable when the tree level $`\mu `$ term is forbidden due to a suitable symmetry .
An interesting feature of this model is the fact that the mass bound of the lightest neutral Higgs scalar can be estimated with no dependence on the soft supersymmetry breaking parameters at the tree level ,
$$m_{h^0}^{(0)2}m_Z^2\left[\mathrm{cos}^22\beta +\frac{2\lambda ^2}{g_1^2+g_2^2}\mathrm{sin}^22\beta \right],$$
(2)
where $`\mathrm{tan}\beta =H_2/H_1`$. This bound is mainly controled by the value of $`\mathrm{tan}\beta `$ and the bound of $`\lambda `$. Its dependence on the soft supersymmetry breaking parameters appears through the loop correction to the effective potential as a result of the large top Yukawa coupling $`h_t`$ . All of these effects are related through the renormalization group equations (RGEs). Many works have been done on this aspect . In Ref. it has been suggested that the additional extra matters such as $`\mathrm{𝟓}+\overline{\mathrm{𝟓}}`$ of SU(5) can heavily affect the mass bound by changing the running of gauge couplings, $`h_t`$ and $`\lambda `$ but not affecting the scale of the gauge coupling unification, which is a great success of the MSSM.
In the NMSSM the radiative generation of nonzero $`\stackrel{~}{S}`$ or the $`\mu `$ scale is a very important aspect. From this point of view it seems to be rather crucial to study how the requirement of the successful occurence of the radiative symmetry breaking affects on the upper mass bound of the lightest neutral Higgs scalar. In this note we investigate this problem by finding the radiatively induced minimum of the effective potential parameterized by $`(\mathrm{tan}\beta ,\stackrel{~}{S})`$ in the models with extra matters $`n(\mathrm{𝟓}+\overline{\mathrm{𝟓}})`$ of SU(5).
The one-loop effective potential due to the large top Yukawa coupling is,
$$V_1=\frac{1}{64\pi ^2}\left[12m_t^4\left(\mathrm{ln}\frac{m_t^2}{Q^2}\frac{3}{2}\right)+\underset{i=1}{\overset{2}{}}6\stackrel{~}{m}_{t_i}^4\left(\mathrm{ln}\frac{\stackrel{~}{m}_{t_i}^2}{Q^2}\frac{3}{2}\right)\right],$$
(3)
where $`Q`$ is a renormalization point and $`\stackrel{~}{m}_{t_i}^2`$ is the eigenvalue of the stop mass matrix
$$\left(\begin{array}{cc}\stackrel{~}{m}_Q^2+m_t^2& m_t(A_t+\lambda \stackrel{~}{S}\mathrm{cot}\beta )\\ m_t(A_t+\lambda \stackrel{~}{S}\mathrm{cot}\beta )& \stackrel{~}{m}_{\overline{T}}^2+m_t^2\end{array}\right).$$
(4)
The correction to Eq. (1) due to this one-loop effective potential is expressed by using these mass eigenvalues as
$$\mathrm{\Delta }m_{h^0}^2=\frac{1}{2}\left(\frac{^2V_1}{v_1^2}\frac{1}{v_1}\frac{V_1}{v_1}\right)\mathrm{cos}^2\beta +\frac{1}{2}\frac{^2V_1}{v_1v_2}\mathrm{sin}2\beta +\frac{1}{2}\left(\frac{^2V_1}{v_2^2}\frac{1}{v_2}\frac{V_1}{v_2}\right)\mathrm{sin}^2\beta .$$
(5)
Here we used the potential minimum condition to eliminate the soft scalar masses of Higgs fields and took a field basis which is used to derive Eq. (2). The mass matrix in Eq. (4) depends on $`(\mathrm{tan}\beta ,\stackrel{~}{S})`$ other than the Yukawa couplings and the soft SUSY breaking parameters.
In order to fix this matrix we need to determine these values as the ones at the potential minimum. It is a nontrivial problem whether such values of $`(\mathrm{tan}\beta ,\stackrel{~}{S})`$ can be radiatively realized starting from certain sets of the Yukawa couplings and the soft SUSY breaking parameters. Our task is to estimate $`m_{h^0}^2(m_{h^0}^{(0)2}+\mathrm{\Delta }m_{h^0}^2)`$ numerically for the parameter sets which can radiatively realize the phenomenologically acceptable potential minimum. We determine such parameter sets so as to satisfy the following conditions:
(i) starting from the suitable initial values of parameters, the radiative symmetry breaking occurs successfully and the following phenomenologically required condition is satisfied at the potential minimum,
$$H_1^2+H_2^2=(174\mathrm{GeV})^2,m_t=174\mathrm{GeV},$$
(6)
(ii) $`m_{h^0}^2`$ which corresponds to the one of diagonal elements of the $`3\times 3`$ neutral Higgs mass matrix should be smaller than other two diagonal components ,
(iii) the experimental mass bounds on the charged Higgs bosons $`m_{H^\pm }`$, charginos $`m_{\chi ^\pm }`$, stops $`\stackrel{~}{m}_{t_i}`$ and gluinos $`M_3`$ are satisfied. These masses, except for $`M_3`$, are dependent on $`\lambda `$ and $`\stackrel{~}{S}`$. We require the following values for them:
$$m_{H^\pm }>65\mathrm{GeV},m_{\chi ^\pm }>72\mathrm{GeV}\stackrel{~}{m}_{t_2}>67\mathrm{GeV},M_3>173\mathrm{GeV},$$
(7)
(iv) the vacuum should be a color conserving one .
In this study we solve a set of RGEs which are composed of two-loop ones for dimensionless couplings and one-loop ones for dimensional SUSY breaking parameters, for simplicity. As the initial conditions for the SUSY breaking parameters we take
$$\stackrel{~}{m}_{\varphi _i}^2=(\gamma _i\stackrel{~}{m})^2,M_a=M,A_t=A_\kappa =A_\lambda =A,$$
(8)
where $`\stackrel{~}{m}`$ is the universal soft scalar mass. We introduce the nonuniversality represented by $`\gamma _i`$ only among soft scalar masses of $`H_1,H_2`$ and $`S`$ to make it easy to find the radiative symmetry breaking solutions. This will be taken as $`0.8\gamma _i1.2`$. These initial conditions are assumed to be applied at the scale $`M_X`$ where the coupling unification of SU(2)<sub>L</sub> and U(1)<sub>Y</sub> occurs. We donot require the regolous coupling unification of SU(3)<sub>C</sub> but only impose the realization of the low energy experimental value . The initial values of the parameters are surveyed through the following region,
$`0h_t1.2(0.1),2.0\kappa 0(0.2),0\lambda 3.0(0.2),`$
$`0M/M_S0.8(0.3),0\stackrel{~}{m}/M_S,|A|/M_S3.0(0.5),`$ (9)
where in the parentheses we give the interval which we use in the survey of these parameter regions.<sup>1</sup><sup>1</sup>1 Since the sign of $`\kappa `$ and $`A`$ affects the scalar potential, we need to investigate both sign of them. However, a negative $`\kappa `$ seems to cover almost solutions for the positive $`\kappa `$. Here we give the only result in the case of the negative $`\kappa `$. We also assume that the RGEs of the model are changed from the supersymmetric ones to the nonsupersymmetric ones at a supersymmetry breaking scale $`M_S`$, for which we take $`M_S=1`$ TeV as a typical numerical value .
To estimate the one-loop effect Eq. (5), it is necessary to know the values of $`(\mathrm{tan}\beta ,\stackrel{~}{S})`$ at the potential minimum. If we impose the radiative symmetry breakling condition, such a potential minimum has to be realized as a result of RGEs solution. In order to see the effect of the radiative symmetry breaking condition on the upper mass bound $`m_{h^0}`$ we calculate it under two situations. We impose the full condition of (i) in a case (I). On the other hand, in a case (II) we require only Eq. (6) but donot require that it is realized at the potential minimum<sup>2</sup><sup>2</sup>2 In other words, this case corresponds to the situation that Eq. (6) may be satisfied at the minimum which is obtained from the unnatural initial soft scalar masses.. We take the number of extra matters as $`n=3`$. The perturbative unification of the gauge coupling requires $`n4`$. Later we will also discuss other cases.
At first, we show how the radiative symmetry breaking condition restricts the parameters strongly relevant to $`m_{h^0}`$. In Fig. 1 we give the plots of the solutions in $`(\mathrm{tan}\beta ,\stackrel{~}{S})`$ and $`(\mathrm{tan}\beta ,\lambda )`$ planes. Since we assume that only the top Yukawa coupling is large, the consistent $`\mathrm{tan}\beta `$ cannot be so large<sup>3</sup><sup>3</sup>3To estimate $`\mathrm{tan}\beta `$ we take account of the translation of the pole mass to the running mass .. We take it as $`\mathrm{tan}\beta 15.0`$. In the case (II), $`\stackrel{~}{S}`$ can take the value in such a wide range as $`70\mathrm{GeV}{}_{}{}^{<}\stackrel{~}{S}{}_{}{}^{<}45\mathrm{TeV}`$, which is not plotted in Fig. 1(a). These show that the value of $`\lambda (m_t)`$ and $`\stackrel{~}{S}`$ are heavily restricted in the case (I) compared with the case (II). We should remind that $`\lambda (m_t)`$ and $`\stackrel{~}{S}`$ affect $`m_{h^0}`$ through Eqs. (2) and (4). As seen from Fig. 1, the $`\mu (\lambda \stackrel{~}{S})`$ scale can take a value in a rather wide range to cause the successful radiative symmetry breaking.
In Fig. 2 we show the allowed region of $`m_{h^0}`$ in both cases for the corresponding parameter sets to the ones of Fig. 1. This shows that the imposition of the consistent occurence of the radiative symmetry breaking can strongly affect the estimation of $`m_{h^0}`$. The boundary value of $`m_{h^0}`$ can be changed by a few percent to ten percent. We can also find some parameter dependences of $`m_{h^0}`$ in these figures. In Fig. 2(b) we restrict the initial soft scalar mass as $`\stackrel{~}{m}=1`$ TeV. The larger values of soft scalar mass $`\stackrel{~}{m}`$ and $`\stackrel{~}{S}`$ tend to realize the larger value of $`m_{h^0}`$. The value of $`\stackrel{~}{S}`$ determines the off-diagonal component of the stop mass matrix (4). This shows that the larger stop mixing tends to make $`m_{h^0}`$ larger.
It is useful to comment on the existence of two branches of the solutions which show a very different behavior in Figs. 1 and 2. This is a common feature in both case (I) and (II) and also in the models with a different $`n`$. For simplicity, we take the case (II) as an example to discuss this feature here. We refer the one of the smaller $`\mathrm{tan}\beta `$ as a branch B1 and the one of the larger $`\mathrm{tan}\beta `$ as a branch B2. They are divided by an initial value of $`h_t`$ as shown in Table 1. The interesting feature of two branches can be clearly seen in Fig. 1(b) and Table 1. The branch B2 corresponds to the smoothly extending solutions of $`\lambda (m_t)`$ to the large $`\mathrm{tan}\beta `$ region and the branch B1 comes from the ones which are confied in the small $`\mathrm{tan}\beta `$. The branch B2 has a similar lower bound of $`\mathrm{tan}\beta `$ and the similar maximum value of $`m_{h^0}143`$GeV at $`\mathrm{tan}\beta {}_{}{}^{>}9`$ for the different $`n`$ values. If we note that at the larger $`\mathrm{tan}\beta `$ region the second term of Eq. (2) can be neglected and Eq. (2) reduces to the one of the MSSM, this behavior can be understood. On the other hand, the branch B1 strongly depends on $`n`$. There are solutions with a little bit larger maximum value of $`\lambda `$ at the smaller $`\mathrm{tan}\beta `$ according to the increase of $`n`$. This results in the larger maximum value of $`m_{h^0}`$ for the larger $`n`$. The reason can be mainly found in the $`\beta `$ dependence of Eq. (2). Decreasing $`n`$, the maximum value of $`\mathrm{tan}\beta `$ of the branch B1 increases, where the maximum value of $`\lambda (m_t)`$ is realized. The branch B1 disappears finally at $`n=0`$ since it is difficult to realize $`0.95{}_{}{}^{<}h_{t}^{}(m_t){}_{}{}^{<}1.35`$ corresponding to $`1{}_{}{}^{<}\mathrm{tan}\beta {}_{}{}^{<}15`$ starting from the initial value of $`h_t`$ used here. We should remind that the value of $`h_t(m_t)`$ is strictly restricted by Eq. (6). In the case (I) we also find the similar qualitative feature discussed here in the case (II).
Finally we want to present one comment. There is a rather big difference of the density of the radiative symmetry breaking solutions in the models with the different $`n`$. For example, in the present parameter setting the number of solutions in $`n=4`$ is very smaller than the ones in $`n=3`$. The finer tuning of marameters seems to be necessary to find the solution in the $`n=4`$ case compared with the $`n=3`$ case.
In summary we studied the upper bound $`m_{h^0}`$ of the lightest neutral Higgs scalar mass in the NMSSM using the RGE analysis. We required the successful occurence of the electroweak radiative symmetry breaking starting from the suitable initial values of parameters. This condition substantially constrains the allowable parameter space and as a result the mass bound $`m_{h^0}`$ is heavily restricted compared with the one obtained without imposing this condition. We discussed the typical feature related to the initial value of $`h_t`$ and also why the larger $`n`$ models could bring the larger maximum value of $`m_{h^0}`$.
This work has been partly supported by the a Grant-in-Aid for Scientific Research from the Ministry of Education, Science and Culture(#11640267 and #11127206).
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# FROM CONSTITUENT QUARK TO HADRON STRUCTURE IN THE NEXT-to-LEADING ORDER: NUCLEON AND PION
## I INTRODUCTION
Our knowledge of hadronic structure is based on the hadronic spectroscopy and the Deep Inelastic Scattering (DIS) data. In the former picture quarks are massive particles and their bound states describe the static properties of hadrons; while the interpretation of DIS data relies upon the quarks of QCD Lagrangian with very small mass. The hadronic structure in this picture is intimately connected with the presence of a large number of partons(quarks and gluons). The two types of quarks not only differ in mass, but also on other important properties. for example, the color charge of quark field in QCD Lagrangian is ill-defined and is not gauge invariant, reflecting the color of gluons in an interacting theory. On the other hand, color associated with a Constituent Quark (CQ) is a well defined entity. It has been recently shown that one can perturbatively dress a QCD Lagrangian field to all orders and construct a CQ in conformity with the color confinement. From this point of view a CQ is defined as a quasi-particle emerging from the dressing of valence quark with gluons and $`q\overline{q}`$ pairs in QCD.
Of course, the concept of CQ as an intermediate step between the quarks of QCD Lagrangian and hadrons is not new. In fact, in the context of $`SU(6)XO(3)`$ long before Altarelli and Cabibo have used them. R.C. Hwa in his elaborated work termed them as valon, extended it and showed its application to many physical processes . In Ref. it is suggested that the concept of dressed quark and gluon might be useful in the area of jet physics and heavy quark effective theory. Despite the ever presence of CQ no one has calculated its content and partonic structure without resorting to hadronic data and the process of deconvolution. The purpose of this paper is threefold: (a) we will evaluate the structure of a CQ in the Next-to-Leading Order (NLO) in QCD. (b) we will verify its conformity with the structure function data of nucleon and pion for which there are ample data available. (c) In the process, however, we will notice that additional refinements are needed to account for the violation of Gottfried Sum Rule (GSR) and the effect of binding of CQ’s to form a physical hadron.
## II FORMALISM
By definition, a CQ is a universal building block for every hadron, that is, its structure is common to all hadrons and generated perturbatively. Once its structure is evaluated, in principle it would permit to calculate the structure of any hadron. In doing so we will follow the philosophy that in a DIS experiment at high enough $`Q^2`$ it is the structure of a CQ which is being probed and at sufficiently low value of $`Q^2`$ this structure cannot be resolved thus, a CQ behaves as valence quark and hadron is viewed as the bound state of its CQ’s. Under these criteria partons of DIS experiments are components of CQ and at high $`Q^2`$ one can write for a U-type CQ its structure as follows:
$$F_2^U(z,Q^2)=\frac{4}{9}z(G_{\frac{u}{U}}+G_{\frac{\overline{u}}{U}})+\frac{1}{9}z(G_{\frac{d}{U}}+G_{\frac{\overline{d}}{U}}+G_{\frac{s}{U}}+G_{\frac{\overline{s}}{U}})+\mathrm{}$$
(1)
where all the functions on the right-hand side are the probability functions for quarks having momentum fraction $`z`$ of a U-type constituent quark at $`Q^2`$. Similar expression can be written for a D-type CQ. Defining the singlet (S) and nonsinglet (NS) CQ distribution functions as:
$$G^S=\underset{i=1}{\overset{f}{}}(G_{\frac{q_i}{CQ}}+G_{\frac{\overline{q}_i}{CQ}})=G_f+(2f1)G_{uf}$$
(2)
$$G^{NS}=\underset{i=1}{\overset{f}{}}(G_{\frac{q_i}{CQ}}G_{\frac{\overline{q}_i}{CQ}})=G_fG_{uf}$$
(3)
where $`G_f`$ is the favored distribution describing the structure function of a quark within a CQ of the same flavor while unfavored distribution $`G_{uf}`$ describes the structure function of any quark of different flavor within the CQ. $`f`$ is the number of active flavors. In terms of singlet and nonsinglet distributions they are as follows:
$$G_f=\frac{1}{2f}(G^S+(2f1)G^{NS})$$
(4)
$$G_{uf}=\frac{1}{2f}(G^SG^{NS})$$
(5)
Having expressed all the structure functions of CQ’s in terms of $`G_f`$ and $`G_{uf}`$ we now go to the moment space and define the moments of these distributions as:
$$M(n,Q^2)=_0^1x^{n2}F(x,Q^2)𝑑x$$
(6)
$$M_i(n,Q^2)=_0^1x^{n1}G_i(n,Q^2)𝑑x$$
(7)
the subscription $`i`$ stands for S, NS. Using charge symmetry, in the following we will refer to CQ distribution only in proton.
In the NLO approximation the dependence of the running coupling constant, $`\alpha `$ on $`Q^2`$ is given by:
$$\alpha (Q^2)=\frac{4\pi }{\beta _0ln(\frac{Q^2}{\mathrm{\Lambda }^2})}\left(1\frac{\beta _1lnln(\frac{Q^2}{\mathrm{\Lambda }^2})}{\beta _0^2ln(\frac{Q^2}{\mathrm{\Lambda }^2})}\right)$$
(8)
with $`\beta _0=\frac{1}{3}(332f)`$ and $`\beta _1=102\frac{38f}{3}`$. The moments of NS and S in the NLO are:
$$M^{NS}(n,Q^2)=[1+\frac{\alpha (Q^2)\alpha (Q_0^2)}{4\pi }\left(\frac{\gamma _{NS}^{(1)N}}{2\beta _0}\frac{\beta _1\gamma _{qq}^{(0)N}}{2\beta _0^2}\right)](\frac{\alpha _s(Q^2)}{\alpha _s(Q_0^2)})^{\gamma _{qq}^{(0)N/2\beta _0}}$$
(9)
$`M^S(n,Q^2)=\{({\displaystyle \frac{\alpha _s(Q^2)}{\alpha _s(Q_0^2)}})^{\frac{\lambda _{}^N}{2\beta _0}}[p_{}^N{\displaystyle \frac{1}{2\beta _0}}{\displaystyle \frac{\alpha _s(Q_0^2)\alpha _s(Q^2)}{4\pi }}p_{}^N\gamma ^Np_{}^N({\displaystyle \frac{\alpha _s(Q_0^2)}{4\pi }}{\displaystyle \frac{\alpha _s(Q^2)}{4\pi }})^{\frac{\lambda _+^N\lambda _{}^n}{2\beta _0}})`$ (10)
$`{\displaystyle \frac{p_{}^N\gamma ^Np_{}^N}{2\beta _0^2+\lambda _+^N\lambda _{}^N}}]+(+)\}`$ (11)
where $`\gamma ^N=\gamma ^{(1)N}\frac{\beta _1}{\beta _0}\gamma ^{(0)N}`$ and $`\gamma ^{(0)N}`$ and $`\gamma ^{(1)N}`$ are anomalous dimension metrices and $`\lambda _\pm ^N`$ denote the eigenvalues of one-loop anomalous dimension matrix $`\gamma ^{(0)N}`$.
$$\lambda _\pm ^N=\frac{1}{2}[\gamma _{qq}^{(0)N}+\gamma _{gg}^{(0)N}\pm \sqrt{(\gamma _{gg}^{(0)N}\gamma _{qq}^{(0)N})^2+4\gamma _{qg}^{(0)N}\gamma _{gq}^{(0)N}}]$$
(12)
and $`p_\pm ^N`$ is given by:
$$p_\pm ^N=\pm (\gamma ^{(0)N}\gamma _\pm ^N)/(\lambda _+^N\lambda _{}^N)$$
(13)
The quantities $`d_\pm ^{(0)}`$, the leading order anomalous dimentions, given in Ref., are related to $`\lambda _\pm `$ in our notations. $`t`$ is the evolution parameter defined as:
$$t=\mathrm{𝑙𝑛}\frac{\mathrm{𝑙𝑛}\frac{Q^\mathit{2}}{\Lambda ^\mathit{2}}}{\mathrm{𝑙𝑛}\frac{Q_\mathit{0}^\mathit{2}}{\Lambda ^\mathit{2}}}$$
(14)
The coefficients $`\gamma _{kl}^{(1,0)N}`$ and can be found in .
We have taken our initial scale $`Q_0^2=0.283`$ $`GeV^2`$ and $`\mathrm{\Lambda }=0.22`$ GeV. It seems that evolution of parton distributions from such a low value of $`Q_0^2`$ is not justified theoretically. The above value of $`Q_0`$ corresponds to a distance of 0.36 fm which is roughly equal to or slightly less than the radius of a CQ. It may be objected that such a distances are probably too large for a meaningful pure perturbative treatment. We note that $`F_2^{CQ}(z,Q^2)`$ has the property that it becomes $`\delta (z1)`$ as $`Q^2`$ is extrapolated to $`Q_0^2`$, which is beyond the region of validity. This mathematical boundary condition signifies that the internal structure of a CQ cannot be resolved at $`Q_0`$ in the NLO approximation. Consequently, when this property is applied to Eq.(19) bellow, the structure function of the nucleon becomes directly related to $`xG_{\frac{CQ}{P}}(x)`$ at those values of $`Q_0`$, that is, $`Q_0`$ is the leading order effective value at which the hadron can be considered as consisting only of three (two) CQ’s, for baryons (mesons). In fact our results are only meaningful for $`Q_0^20.4`$ $`GeV^2`$. As it is stated above, the moments of the CQ structure function, $`F_2^{CQ}(z,Q^2)`$ are expressed completely in terms of the evolution parameter, $`t`$, of Eq. (13). From the theoretical standpoint, both $`\mathrm{\Lambda }`$ and $`Q_0`$ depend on the order of the moments $`n`$. In this work we have assumed that they are independent of $`n`$, hence introducing some degrees of approximation to the $`Q^2`$ evolution of the valence and sea quarks. However, on one hand there are other contributions like target-mass effects, which add uncertainties to the theoretical predictions of perturbative QCD, while on the other hand since we are dealing with the CQ, there is no experimental data to invalidate an $`n`$ independent $`\mathrm{\Lambda }`$ assumption. The moments of valence and sea quarks in a CQ are:
$$M_{\frac{valence}{CQ}}=M^{NS}(n,Q^2)$$
(15)
$$M_{\frac{sea}{CQ}}=\frac{1}{2f}(M^SM^{NS})$$
(16)
where $`M^{S,NS}`$ are given above. Evaluating $`M_{\frac{valence}{CQ}}`$ and $`M_{\frac{sea}{CQ}}`$ at any $`Q^2`$ or $`t`$ is now straight forward. Using Inverse Mellin Transform techniques, following forms for the valence and sea quark distributions inside a CQ is obtained in the NLO:
$$zq_{\frac{val.}{CQ}}(z,Q^2)=az^b(1z)^c$$
(17)
$$zq_{\frac{sea}{CQ}}(z,Q^2)=\alpha z^\beta (1z)^\gamma [1+\eta z+\xi z^{0.5}]$$
(18)
The parameters $`a`$, $`b`$, $`c`$ , $`\alpha `$, etc. are functions of $`Q^2`$ through the evolution parameter $`t`$. The same form as in Eq.(17) is obtained for Gluon distribution in a CQ but only with different parameters. Functional form of them is a polynomial of order three in $`t`$ and are given in the appendix. We notice that the following sum rule reflecting the fact that each CQ contains only one valence quark is satisfied for all values of $`Q^2`$:
$$_0^1q_{\frac{val.}{CQ}}(z,Q^2)𝑑z=1.$$
(19)
Substituting these results in Eq.(1) completes the evaluation of a constituent quark structure function in NLO. In Figure (1) various parton distributions inside a CQ is plotted.
## III HADRONIC STRUCTURE
In previous section we calculated the NLO structure of a CQ. In this section we will use the convolution theorem to calculate the structure function of proton, $`F_2^p(x,Q^2)`$, and that of a pion. Let us denote the structure function of a CQ by $`F_2^{CQ}(z,Q^2)`$ and the probability of finding a CQ carrying momentum fraction $`y`$ of the hadron by $`G_{\frac{CQ}{h}}(y)`$. The corresponding structure function of the hadron, using the convolution theorem is as follows:
$$F_2^h(x,Q^2)=\underset{CQ}{}_x^1\frac{dy}{y}G_{\frac{CQ}{h}}(y)F_2^{CQ}(\frac{x}{y},Q^2)$$
(20)
where summation runs over the number of CQ’s in a particular hadron. Also notice that $`G_{\frac{CQ}{h}}(y)`$ is independent of the nature of the probe and its $`Q^2`$ value. $`G_{\frac{CQ}{h}}(y)`$ in effect, describes the wave function of hadron in CQ representation containing all the complications due to confinement. From the theoretical point of view this function cannot be evaluated accurately. To facilitate phenomenological analysis, following Ref. we assume a simple form for the exclusive CQ distribution in proton and pion as follows:
$$G_{UUD/p}(y_1,y_2,y_3)=l(y_1y_2)^my_3^n\delta (y_1+y_2+y_31)$$
(21)
$$G_{\overline{U}D/\pi ^{}}(y_1,y_2)=qy_1^\mu y_2^\nu \delta (y_1+y_21)$$
(22)
Integrating over unwanted momenta, we can arrive at inclusive distribution of individual CQ:
$$G_{U/p}(y)=\frac{1}{B(a+1,b+a+2)}y^a(1y)^b$$
(23)
$$G_{D/p}(y)=\frac{1}{B(b+1,2a+2)}y^b(1y)^{2a+1}$$
(24)
$$G_{\overline{U}/\pi ^{}}(y)=\frac{1}{B(\mu +1,\nu +1)}y^\mu (1y)^\nu $$
(25)
similarly expression for $`G_{D/\pi ^{}}`$ with the interchange of $`\mu \nu `$. In the above equations $`B(i,j)`$ is Euler Beta function and its arguments are fixed using the sum rule:
$$_0^1G_{\frac{CQ}{h}}(y)𝑑y=1$$
(26)
where $`CQ=U,D,\overline{U}`$ and $`h=p,\pi ^{}`$. Numerical values are: $`\mu =0.01`$, $`\nu =0.06`$, $`a=0.65`$ and $`b=0.35`$. In Figure (2) the CQ distributions in proton and $`\pi ^{}`$ are shown. We stress that CQ distributions in hadrons are independent of $`Q^2`$ and the nature of probe being used. Now it is possible to determine various parton distributions in a hadron. For proton we can write:
$$q_{val./p}(x,Q^2)=2_x^1\frac{dy}{y}G_{U/p}(y)q_{val./U}(x,Q^2)+_x^1\frac{dy}{y}G_{D/p}(y)q_{val./D}(x,Q^2)=u_{val./p}(x,Q^2)+d_{val./p}(x,Q^2)$$
(27)
$$q_{sea/p}(x,Q^2)=2_x^1\frac{dy}{y}G_{U/p}(y)q_{sea/U}(x,Q^2)+_x^1\frac{dy}{y}G_{D/p}(y)q_{sea/D}(x,Q^2)$$
(28)
The above equation represents the contribution of constituent quarks to the nucleon sea. Comparing with data on proton structure functions shows that the results fall short of representing the experimental data by a 3-5 percent. This is due in our opinion, to the fact that CQ is not free in a hadron but they interact with each other in forming the bound states. That means, there are some soft gluons in the nucleon besides the CQ’s. In the process of formation of bound state, a CQ emits gluon which in turn decays into $`\overline{q}q`$ pairs which gives a residual component to the partons in a hadron. In our picture there is no room in the CQ structure for breaking the $`SU(2)`$ symmetry of the sea but after creation of $`\overline{q}q`$ pair from the emitted gluon , these quarks can recombine with CQ to fluctuate into meson-nucleon state which breaks the symmetry of the nucleon see. In what follows we will compute this component, its contribution to the nucleon structure function and violation of Gottfried sum rule following the preswcription used in . In order to distinguish these partons from those confined inside the CQ, we will term them as inherent partons. Although this component is intimately related to the bound state problem, and hence it has a non-perturbative origin, for not so small values of $`Q^2`$ we will calculate it perturbatively for the process of $`CQCQ+\mathrm{𝑔𝑙𝑢𝑜𝑛}\overline{q}q`$ at an initial value of $`Q^2=0.65GeV^2`$ where $`\alpha _s`$ is still small enough. The corresponding splitting functions are as follows:
$$P_{gq}(z)=\frac{4}{3}\frac{1+(1z)^2}{z}$$
(29)
$$P_{qg}(z)=\frac{1}{2}(z^2+(1+z)^2)$$
(30)
For the joint probability distribution of the process at hand, we get:
$$q_{inh.}(x,Q^2)=\overline{q}_{inh}(x,Q^2)=N\frac{\alpha _s^2}{(2\pi )^2}_x^1\frac{dy}{y}P_{qg}(\frac{x}{y})_y^1\frac{dz}{z}P_{gq}(\frac{y}{z})G_{CQ}(z)$$
(31)
The splitting functions and the $`q_{inh}(x,Q^2)=\overline{q}_{inh}(x,Q^2)`$, above are that of the leading order rather than NLO. We do not expect it should make much of a difference since, its contribution to the whole structure function is only a few percent as can be seen in Figure (3). In the above equation $`N`$ is a factor depending on $`Q^2`$ and $`G_{CQ}`$ is the constituent quark distribution in the proton given previously. The same process, however, also can be a source of $`SU(2)`$ symmetry breaking of nucleon sea and resulting in $`u_{sea}d_{sea}`$, and hence the violation of Gottfried sum rule (GSR). There are several explanations forthis observation such as flavor asymmetry of the nucleon sea , isospin symmetry breaking between proton and neutron, Pauli blocking, etc. One of these explanations fits well within our model. It was proposed by Eichten, Inchliffe and Quigg that valence quark fluctuates into quark and a pion. In other words a nucleon can fluctuate into a meson-nucleon state. This idea is appealing in our model and can be calculated rather easily. In our model after a pair of inherent $`q\overline{q}`$ created a $`\overline{u}`$ can couple to a D-type CQ to form an intermediate $`\pi ^{}=D\overline{u}`$ while the $`u`$ quark combines with the other two U-type CQ’s to form a $`\mathrm{\Delta }^{++}`$. This is the lowest $`u\overline{u}`$ fluctuation. Similarly a $`d\overline{d}`$ can fluctuate into the $`\pi ^+n`$ state. Since $`\mathrm{\Delta }^{++}`$ state is more massive than $`n`$ state, then the probability of $`d\overline{d}`$ fluctuation will dominate over $`u\overline{u}`$ fluctuation which naturally leads to an excess of $`d\overline{d}`$ pairs over $`u\overline{u}`$ in the proton sea. This process is depicted in Figure (4). Probability of formation of a meson-barion state can be written as in Ref. :
$$P_{MB}(x)=_0^1\frac{dy}{y}_0^1\frac{dz}{z}F(y,z)R(y,z;x)$$
(32)
where $`F(y,z)`$ is the joint probability of finding a CQ with momentum fraction $`y`$ and an inherent quark or anti-quark of momentum fraction $`z`$ in the proton. $`R(y,z;x)`$ is the probability of recombining a CQ of momentum $`y`$ with an inherent quark of momentum $`z`$ to form a meson of momentum fraction $`x`$ in the proton. The evaluation of both of these probability functions are discussed in for a more general case and an earlier, but pioneering, version also proposed in . In the present case these functions are much simpler. Guided by works done in Ref. we can write :
$$F(y,z)=\mathrm{\Omega }yG_{\frac{CQ}{p}}(y)z\overline{q}_{inh.}(z)(1yz)^\delta $$
(33)
$$R(y,z;x)=\rho y^az^b\delta (y+z1)$$
(34)
Here we take $`a=b=1`$ reflecting that two CQ’s in meson almost equally share its momentum. The exponent $`\delta `$ is fixed for the $`n`$ and $`\mathrm{\Delta }^{++}`$ states using the data from E866 experiment and the mass ratio of $`\mathrm{\Delta }`$ to $`n`$. They turn out to be approximately 18 and 13 respectively. $`\mathrm{\Omega }`$ and $`\rho `$ are the normalization constants also fixed by data. It is now possible to evaluate $`\overline{u}_M`$ and $`\overline{d}_M`$ quarks associated with the formation of meson states:
$$\overline{d}_M(x,Q^2)=_x^1\frac{dy}{y}[P_{\pi n}(y)+\frac{1}{6}P_{\pi \mathrm{\Delta }^{++}}(y)]d_\pi (\frac{x}{y},Q^2)$$
(35)
$$\overline{u}_M(x,Q^2)=\frac{1}{2}_x^1\frac{dy}{y}P_{\pi \mathrm{\Delta }^{++}}(y)u_\pi (\frac{x}{y},Q^2)$$
(36)
where $`u_\pi `$ and $`d_\pi `$ are the valence quark probability densities in the pion at scale $`Q_0^2`$. The coefficients $`\frac{1}{2}`$ and $`\frac{1}{6}`$ are due to Isospin consideration. Using Eqs. (16, 17, 24) we can calculate various parton distributions in a the pion. Those pertinent to Eqs. (34, 35) are:
$$\overline{u}_{val.}^\pi ^{}(x,Q^2)=_x^1G_{\overline{U}/\pi ^{}}(y)\overline{u}_{val./\overline{U}}(\frac{x}{y},Q^2)\frac{dy}{y}$$
(37)
$$d_{val.}^\pi ^{}(x,Q^2)=_x^1G_{D/\pi ^{}}(y)d_{val./D}(\frac{x}{y},Q^2)\frac{dy}{y}$$
(38)
$$\overline{u}_{val./\overline{U}}=u_{val./U}$$
(39)
There are some data on the valence structure function of $`\pi ^{}`$ . Defining valence structure function of $`\pi ^{}`$ as:
$$F_{val.}^\pi ^{}=x\overline{u}_{val.}^\pi ^{}=xd_{val.}^\pi ^{}$$
(40)
we present the results of our calculation for $`F_{val.}^\pi ^{}`$ in Figure (5) along with the experimental data at $`Q^2`$ around 6 $`Gev^2`$. Returning to the $`\overline{d}`$ excess over $`\overline{u}`$ in proton we can write all the contributions as:
$$(\frac{\overline{d}}{\overline{u}})_{proton}=\frac{\overline{d}_M+\overline{d}_{inh.+CQ}}{\overline{u}_M+\overline{u}_{inh.+CQ}}$$
(41)
NuSea collaboration at FermiLab E866 experiment has published their results for integral of $`\overline{d}\overline{u}`$ and $`\frac{\overline{d}}{\overline{u}}`$ at $`Q=7.35`$ $`GeV`$. With the procedure described, we have calculated these values at the same $`Q`$ and for the range of $`x`$ as experimentally measured: $`x=[0.02,0.35]`$. The results of the model is:
$$_{0.02}^{0.345}𝑑x(\overline{d}\overline{u})=0.085$$
(42)
to be compared with the experimental value of $`0.068\pm 0.0106`$. We get for the entire range in $`x`$: $`_0^1𝑑x(\overline{d}\overline{u})=0.103`$ while the experimentally extrapolated value is $`0.1\pm 0.018`$ which are in excellent agreement with our calculations. This gives a value for the Gottfried sum rule of:
$$S_G=_0^1[F_2^p(x)F_2^n(x)]\frac{dx}{x}=\frac{1}{3}\frac{2}{3}_0^1𝑑x[\overline{d}(x)\overline{u}(x)]=0.264.$$
(43)
at $`Q=7.35`$ $`GeV`$. The NMC result is $`S_G=0.235\pm 0.026`$ which is at much lower value of $`Q^2=4`$ $`GeV^2`$. In Figure (6), $`\overline{d}(x)\overline{u}(x)`$ and $`\frac{\overline{d}}{\overline{u}}`$ in proton are shown as a function of $`x`$ at $`Q=7.35`$ $`GeV`$ along with the measured results.
We are now in a position to present the results for the proton structure function, $`F_2^p`$. In Eqs.(16,17) we presented the form of parton distributions in a CQ. Using those relations with the numerical values given in appendix then from Eq.(1) the structure function of CQ is obtained. In Eqs. (22, 23) the shape of CQ distributions in proton is given. With the help of Eq. (19) now all the ingredients are in place to calculate $`F_2^p`$. In Figure (7) the results are shown at many values of $`Q^2`$. As it is evident they fall a few percent short of representing the data. However, as mentioned earlier, there is an additional contribution from the inherent partons to $`F_2^p`$ which is calculated in Eq. (30). Adding this component represents the data rather well and can be seen from Figure (7). The data points are from . For the purpose of comparing our results with other calculations, we have also included in Figure (7), the GRV’s NLO results as well as the prediction of CTEQ4M . Notice that we have taken the number of active flavors to be three for $`Q^25GeV^2`$ and four flavors elsewhere. In Figure (8) the gluon distribution predicted by the model is presented along with those from Ref..
## IV summary and conclusion
In this paper we have used the notion of constituent quark as a well defined entity being common to all hadrons. Its structure can be calculated in QCD perturbatively to all orders. A CQ receives its own structure by dressing a valence quark with gluon and $`q\overline{q}`$ pairs in QCD. We have calculated its structure function in the Next-to-Leading order for the first time. Considering a hadron as the bound states of these CQs we have used the convolution theorem to extract the hadronic structure functions for proton and pion. Besides the CQ structure contribution to the hadrons, there is also a nonperturbative component while contributing only a few percent to the overall structure of hadrons becomes crucial in explaining the violation of Gottfried sum rule and the excess of $`\overline{d}`$ over $`\overline{u}`$ in the nucleon sea. A mechanism is devised for this purpose and necessary calculations are outlined. We have presented the results and compared them with all available relevant data and with the work of others. We found that our results are in good agreement with the data.
## V APPENDIX
In this appendix we will give the functional form of parameters of Eqs. (16, 17) in terms of the evolution parameter, $`t`$. This will completely determines partonic structure of CQ and their evolution. The results are valid for three and four flavors, although the flavor number is not explicitly present but they have entered in through the calculation of moments. As we explained in the text, we have taken the number of flavors to be three for $`Q^25GeV^2`$ and four for higher $`Q^2`$ values.
I) Valence quark in CQ (Eq. 16):
$`a=0.1512+1.785t1.145t^2+0.2168t^3`$
$`b=1.4601.137t+0.471t^20.089t^3`$
$`c=1.031+1.037t0.023t^2+0.0075t^3`$
II) Sea quark in CQ (Eq. 17):
$`\alpha =0.0700.213t+0.247t^20.080t^3`$
$`\beta =0.3361.703t+1.495t^20.455t^3`$
$`\gamma =20.526+57.495t46.892t^2+12.057t^3`$
$`\eta =3.1879.141t+10.000t^23.306t^3`$
$`\xi =7.914+19.177t18.023t^2+5.279t^3`$
$`N=1.023+0.124t2.306t^2+1.965t^3`$
III) Gluon in CQ (Eq. 17)
$`\alpha =0.8261.643t+1.856t^20.564t^3`$
$`\beta =0.3281.363t+0.950t^20.242t^3`$
$`\gamma =0.482+1.528t0.223t^20.023t^3`$
$`\eta =0.4803.386t+4.616t^21.441t^3`$
$`\xi =2.375+6.873t7.458t^2+2.161t^3`$
$`N=2.2476.903t+6.879t^21.876t^3`$
## VI Figure Caption
Figure-1. Moments of partons in a CQ at $`Q^2=20`$ $`GeV^2`$ as a function of $`z`$.
Figure-2. Parton distributions in proton and $`\pi ^{}`$ at $`Q^2=20`$ $`GeV^2`$ as a function of $`x`$.
Figure-3. Contribution of inherent component to a sea parton distribution. The dashed-dotted line is that of CQ and the solid line represents the sum of the two components.
Figure-4. Processes responsible to $`SU(2)`$ symmetry breaking in the nucleon sea and violation of Gottfried sum rule.
Figure-5. Pion valence structure function as a functuon of $`x`$ at $`Q^2=5.5`$ $`(Gev/c)^2`$. Solid curve is the result of the model calculations and the data points are from Ref..
Figure-6. The ratio $`\frac{\overline{d}}{\overline{u}}`$ and the difference $`\overline{d}\overline{u}`$ as a function of $`X`$. The solid line in the model calculation and the dotted line is the prediction of CTEQ4M. Data are from Ref. .
Figure-7. Proton structure function $`F_2^p`$ as a function of $`x`$ calculated using the model and compared with the data from Ref. for different $`Q^2`$ values. The thin line is the prediction of GRV Ref. and the dashed line is that of CTEQ4M Ref..
Figure-8. The gluon distribution in proton as a function of $`x`$ at $`Q^2=20`$ $`GeV^2`$. We have also shown the prediction of GRV (dashed-dotted line) and CTEQ4M (dashed line). The data points are from H1 collaboration.
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# Vacuum structure in QCD and Correlation functions*footnote **footnote * Talk prepared for DAE Nuclear Physics Symposium, December 27th–31st, 1999,Chandigarh, India
## I Introduction
Quantum chromodynamics (QCD) is now accepted to be the theory of strong interaction in terms of quarks and gluons, and, at a secondary level, of hadrons. It is a nonabelian gauge theory with $`SU(3)`$ as the gauge group. This nonabelianess of the interaction leads to two important consequences. Firstly, the interaction become weak at high momentum (above several GeV) transfer processes where perturbation theory is applicable. At low energies and momenta, relevant to most of nuclear physics, the quark gluon coupling strength becomes large and an expansion in powers of this coupling is not useful. The basic difficulty appears to be an understanding of the ground state properties of QCD or its vacuum structure which plays an important role for the related physics .
Conceptually, low energy QCD has many common feature with condensed matter physics. The vacuum here appear to be consisting of having fluctuating quarks and gluon fields with average properties being described by condensates of quarks and gluons. The quark condensate $`\overline{q}q`$, i.e. expectation value of quark scalar densities, plays an important role in the context of strong interactions. It displays the breaking of chiral symmetry. This symmetry is crucial to our understanding of nucleons,nuclei and dense matter. The other quantity that charcterise the QCD vacuum is the gluon condensate. The energy density of vacuum is lowered by presence of electric and magnetic gluon fields in the ground state. These condensates were introduced in the context of QCD sumrules and the values are estimated from the charmonia spectroscopy as $`\alpha _s/\pi GG=2\alpha _s/\pi ^2(𝐄^2𝐁^2)=0.012GeV^4`$ and $`\overline{q}q=(250MeV)^3`$ .
Thus QCD vacuum poses a rich complicated many body problem. As in other many body system ground state correlation functions give an insight to the ground state structure. One might be interested in asking questions like what happens to a meson pair when placed in such a condensate medium or what happens to interquark interaction as a function of their spatial separation. We might remind ourselves that nucleon scattering phase shifts gives information regading inter-nucleon forces complementary to that infered from their bound state namely from the properties of deuteron. N-N scattering allows one to study different components of nuclear forces (spin-spin, tensor,$`\mathrm{}`$) at different spatial separation in much more detail than deuteron observables which is an composite effects of all channels. In QCD we however do not have free quarks or gluons due to confinement. None the less, one can infer about inter-quark interaction as a function of their separation in different channels by studying propagators and correlation functions of hadron currents in a more detailed manner than the composite effects reflected in hadron bound state. The correlation functions that we shall consider here are spacelike separated correlation functions. To be precise we shall take the correlation functions to be defined at equal time so that their separation is purely spatial. Such correlation functions have several appealing features. They describe different physics at different spatial separations; they can be calculated in some channels phenomenologically from $`e^+e^{}`$hadrons data or from $`\tau `$ decay experimental data. Moreover they can be evaluated in lattice simulations or in some model for vacuum.
Our appraoch here shall be assuming a structure for the vacuum and then examining its consequences regarding correlator phenomenology. More precisely, we shall use phenomenological results of correlation functions of hadronic currents to guide us towards a “true” structure of QCD vacuum. We organise this note as follows. In section II, we shall discuss a construct for the vacuum in terms of quark and gluon condensates , and the resulting correlation functions . It appears that condensates alone do not give rise to correct phenomenology of correlators unless one includes fluctuations of such condensates . Inclusion of the fluctuation fields yields correlation functions consistent with the phenomenology of correlation functions. Such a structure of vacuum of zero temeperature is then generalised to finite temperature in section III. Here we shall also compute the finite temperature correlators to determine temperature dependent hadron properties . In section IV, we shall discuss the vacuum structure at high density and the QCD vacuum with diquark condensate giving rise to color superconductivity. Finally, we summarize our results with some remarks and discussions in section V.
## II An ansatz for the QCD vacuum and correlation functions
A variational ansatz was considered in Ref. for the QCD vacuum with an explicit construct involving quark antiquark pairs and gluon pairs. The trial ansatz was given as
$$|vac=\mathrm{exp}(B_F^{}B_F)(B_G^{}B_G)|0$$
(1)
where, the Bogoliubov pair creation operators for the quarks and the gluons are given respectively as
$$B_F^{}=q_I^0(𝐤)^{}(𝝈𝐤)h(𝐤)\stackrel{~}{q}_I^0(𝐤)𝑑𝐤,$$
(3)
and,
$$B_G^{}=a_i^a(𝐤)^{}g(𝐤)a_i^a(𝐤)𝑑𝐤.$$
(4)
In the above $`q_I^0,\stackrel{~}{q}_I^0`$ are two component quark and antiquark annihilation operators respectively. The subscript $`0`$ indicates that they annihilate the peruturbative vacuum $`|0`$ i.e. $`q_I^0|0=0=\stackrel{~}{q}^{}|0`$. $`a_i^a`$ is the gluon annihilation operator. The operators satisfy the quantum algebra given in Coulomb gauge as
$$[q_{Ir}^{i0}(𝐤),q_{Is}^{0j}(𝐤^{})^{}]_+=\delta ^{ij}\delta _{rs}\delta (𝐤𝐤^{}),$$
(6)
and,
$$[a_i^a(𝐤),a_j^b(𝐤^{})]=\delta ^{ab}\left(\delta _{ij}𝐤_i𝐤_j/𝐤^2\right)\delta (𝐤𝐤^{}).$$
(7)
Finally, $`h(𝐤)`$ and $`g(𝐤)`$ are two trial functions associated with the quark antiquark condensates and gluon pairs respectively. Clearly, a construct as in Eq.(1)has an obvious parallel to BCS theory of superconductivity. Such a structure for vacuum eventually reduces to Bogoliubov transformation for the operators. Then one can calculate the, energy functional–the expectation value of the Hamiltonian which is a functional of the condensate functions. Since the functions cannot be determined through functional minimisations one can parametrise the condensate functions with some trial functions simple enough to manipulate numerically as well as reasonable enough to simulate correct physical behaviour. In Ref., the following choices were made (with $`k=|𝐤|`$)
$$\mathrm{tan}2h(𝐤)=\frac{A^{}}{(\mathrm{exp}(R^2k^2)1)^{1/2}}$$
(9)
$$\mathrm{sinh}g(𝐤)=A\mathrm{exp}(bk^2)$$
(10)
In Ref. the energy density was minimised with respect to the condensate parameters subjected to the constraint that the pion decay constant $`f_\pi `$ and the gluon condensate value $`(\alpha _s/\pi )GG`$ turns out to be the experimental values of 93 MeV and $`0.012GeV^4`$ respectively. The result of such a minimisation then leads to instabilty of perturbative vacuum to condensate formation beyond a critical coupling of $`\alpha _s^c=0.6`$. For $`\alpha _s=1.28`$, the charge radius of pion comes out correctly. The values of $`A^{}`$ and $`R`$ turns out to be $`A^{}1`$ and $`R0.96fm.`$ Further some of the baryonic properties like charge radius of proton, magnetic moments of proton and neutron turns out to be close to their corresponding experimental values. Further, the bag constant– the energy difference between the perturbative and the nonperturbative vacuum turns out to be $`ϵ_0=(140MeV)^4`$which appears to be in general agreement with the phenomenological value of this parameter.
With such a description of the vacuum in terms of condensates let us look at the correlation functions and propagators in a condensed medium. The equal time propagator is given as
$`S_{\alpha \beta }(𝐱)`$ $`=`$ $`{\displaystyle \frac{1}{2}}[q_\alpha (𝐱),\overline{q}_\beta (\mathrm{𝟎})]`$ (11)
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{(2\pi )^3}}{\displaystyle e^{i𝐤𝐱}\left[\mathrm{sin}2h(𝐤)𝜸𝐤\mathrm{cos}2h(𝐤)\right]}`$ (12)
Clearly, free massless propagator corresponds to $`h(𝐤)0`$ limit of above equation and is given as $`S_0=(i/2\pi ^2)(𝜸𝐱/x^4)`$. Different components of the propagator given in Eq.(LABEL:prop) can be analysed and it turns out to be qualitatively similar to those obtained from other nonperturbative calculations like an instanton liquid model for QCD vacuum. Further, a small x expansion of Eq.(LABEL:prop) yields the propagtor as that one would get in the operator product expansion in the vacuum saturation approximation.
We shall next consider the correlation functions of mesonic currents of generic form $`J(𝐱)=\overline{q}_\alpha (𝐱)\mathrm{\Gamma }_{\alpha \beta }q_\beta (𝐱)`$, where $`\alpha ,\beta `$ are spinor indices and $`\mathrm{\Gamma }`$ is a 4$`\times `$4 matrix in Dirac space $`(1,\gamma _5,\gamma _\mu ,\gamma _\mu \gamma _5)`$. The equal time correlation function is defined as
$$R(𝐱)=\frac{1}{2}vac|J(𝐱)\overline{J}(\mathrm{𝟎})|vac+vac|\overline{J}(\mathrm{𝟎})J(𝐱)|vac$$
(14)
With the definition of condensate vacuum as in Eq.(1) and the propagator given in Eq.(LABEL:prop), Eq.(14) reduces to
$$R(x)=Tr[S(x)\mathrm{\Gamma }^{}S(x)\mathrm{\Gamma }]$$
(15)
where, $`\mathrm{\Gamma }^{}=\gamma _0\mathrm{\Gamma }^{}\gamma _0`$. We shall be normalising the correlation functions with that of free massless correlation function which is given as parallel to Eq(15) as
$$R_0(x)=Tr[S_0(x)\mathrm{\Gamma }^{}S_0(x)\mathrm{\Gamma }]$$
(16)
The ratio $`R(x)/R_0(x)`$ can then be evaluated for different channels. It turns out that correlation function so obtained has similar qualitative behaviour as one can obtain from phenomenology in all chnnels except for the pseudoscalar channel . Phenomenologically, in this channel there is a strong attraction with the ratio becoming about 100 at a separation of about a fermi. The calculated value however turns out to be as low as 1.2 around that value. All these results depend very weakly on the functional form of $`h(𝐤)`$.
In view of this outcome, it is obvious that some crucial physics is missing from the model of vacuum considered in Eq.(1) and has to be supplemented by additional effects. In the present framework this means that quark propagators alone do not describe the correlation functions and there has to be contributions from irreducible four point structure of the vacuum. This can be thought of as a manifestation of fluctuation of the condensates. Thus we have
$`T\overline{q}_\alpha (𝐱)q_\beta (𝐱)q_\gamma (\mathrm{𝟎})q_\delta (\mathrm{𝟎})`$ $`=`$ $`S_{\beta \gamma }(𝐱)S_{\delta \alpha }(𝐱)+:\overline{q}_\alpha (𝐱)q_\beta (𝐱)q_\gamma (\mathrm{𝟎})q_\delta (\mathrm{𝟎}):`$ (17)
$`=`$ $`S_{\beta \gamma }(𝐱)S_{\delta \alpha }(𝐱)+\mathrm{\Sigma }_{\beta \gamma }(𝐱)\mathrm{\Sigma }_{\delta \alpha }(zbfx)`$ (18)
where, we have introduced the composite fields $`\mathrm{\Sigma }`$ to include the effects of fluctuation with $`vac|\mathrm{\Sigma }\mathrm{\Sigma }|vac=0`$ but $`\mathrm{\Omega }|\mathrm{\Sigma }\mathrm{\Sigma }|\mathrm{\Omega }0`$, where, $`|\mathrm{\Omega }`$is the “new improved” QCD vacuum including the condensate fluctuations. The correlation function then takes the form
$$R(x)=[Tr[S(x)\mathrm{\Gamma }^{}S(x)\mathrm{\Gamma }]+Tr[\mathrm{\Sigma }(𝐱))\mathrm{\Gamma }^{}\mathrm{\Sigma }(x)\mathrm{\Gamma }]]$$
(20)
The structure of $`\mathrm{\Sigma }`$ field should be such that it contributes mostly to the pseudoscalar channel and should not affect the other channels very much. Such a condition restricts the composite field to be of the form
$$\mathrm{\Sigma }_{\alpha \beta }(𝐱)=\mu _1^2(\gamma ^i\gamma ^jϵ_{ijk}\varphi ^k(𝐱)+\mu _2^2\delta _{\alpha \beta }\varphi (𝐱)$$
(21)
where, we have introduced the scalar and vector fileds $`\varphi `$ and $`\varphi ^k`$ such that
$$\mathrm{\Omega }|\varphi ^i(𝐱)\varphi ^j(\mathrm{𝟎})|\mathrm{\Omega }=\delta ^{ij}g_V(𝐱)$$
(23)
$$\mathrm{\Omega }|\varphi (𝐱)\varphi (\mathrm{𝟎})|\mathrm{\Omega }=g_S(𝐱)$$
(24)
From general considerations, we may write down the functions $`g_V(𝐱),g_S(𝐱)`$ as
$$g_V(𝐱)=\frac{1}{2\pi ^2x}\left[\mu _1K_1(\mu _1x)\mu _3K_1(\mu _3x)\right]$$
(26)
$$g_S(𝐱)=\frac{1}{2\pi ^2x}\left[\mu _2K_1(\mu _4x)\mu _5K_1(\mu _6x)\right]$$
(27)
Using Eq.s(20,21,26,27) one can have the expression for the correlation functions in different channels. Explicitly different currents in different channels and the contributions from the propagator and the fluctuation fields are shown separately in Table I. We have also included here the currents for the baryons. The resulting correlation functions normalised to correlation functions obtained by treating the quarks as massless and noninteracting are plotted in Fig.1.
To get the hadron parameters from the correlation functions one relates the correlation function through a dispersion relation to a spectral density function. For example in the vector channel the dispersion relation for the correlation function reduces to
$$\mathrm{\Pi }_{\mu \mu }(x)\mathrm{\Omega }|T(J_\mu (x)\overline{J}_\mu (0))|\mathrm{\Omega }=\frac{1}{4\pi ^2}_0^{\mathrm{}}𝑑sR_i(s)D(s^{1/2},x)$$
(28)
with, the function $`D(m,x)=(m/4\pi ^2x)K_1(mx)`$ is a propagator of a particle of mass $`m`$. The function $`R_i(x)`$ is the normalised spectral density function related to the cross section of $`e^+e^{}`$ annihilation to hadrons. In presence of interaction, in general we do not know how to calculate the spectral density function . However we do know their qualitative behaviour based on experimental information. To extract hadron parameters one parametrises the spectral function and then determine the parameters. A useful parametrisation valid at zero temperature is in terms of a Dirac delta function at the pole mass of the hadron accompanied by a step function continuum at higher energy. Specifically, for $`\rho `$ meson the spectral density function is parametrised by
$$R_\rho (s)=3\lambda ^2\delta (sM^2)+\frac{3s}{4\pi ^2}\theta (ss_0).$$
(29)
where, M is the bound state mass, $`\lambda `$ is the coupling of the current to the bound state and $`s_0`$ is the threshold for continuum contribution. The fitted parametrs to the correlation functions obatined by us (solid curve in Fig.1 are given in Table II
As may be evident from this the results are comparable with that of lattice calculation and instanton liquid model calculations . Thus to be consistent with the data QCD vacuum must not only have condensates but also have their fluctuations.
## III QCD vacuum and correlation functions at finite temperature
As is well known the QCD vacuum state changes with temperature. Lattice Monte Carlo simulations suggest that chiral symmetry is restored around 150 MeV. Thus it is interesting to look at correlation functions at finite temperature in the context of behaviour of hadrons around the chiral phase transition . It may be noted that there is little phenomenological information in this regime but there are several theoretical studies using operator product expansion (OPE) and sum rule methods as well as using instanton liquid model for QCD ground state . In particular, we shall generalise the method considerd in the previous section to include temperature effects. The zero temperature vacuum defined in Eq.(1) can be generalised to finite temperature using the methodology of thermofield dynamics. Here the thermal average of an operator is obtained as an expectation value of the operator over the thermal vacuum. The thermal vacuum is obtained from the zero temperature vacuum by a thermal Bogoliubov transformation in an extended Hilbert space involving extra field operators (thermal doubling of operators). Explicitly, the thermal vacuum is given as
$$|vac,\beta =\mathrm{exp}(d𝐤\theta (𝐤,\beta )(q_{I}^{}{}_{}{}^{}(𝐤)\underset{¯}{q_I}^{}(𝐤)+\stackrel{~}{q}_I(𝐤)\underset{¯}{\overset{~}{q}}_I(𝐤))h.c.)|vac$$
(30)
where, the underlined oprators are the operators corresponding to extra of Hilbert space. Further, the $`q_I^{}s`$ are the opeartors refers to the fact that they are the quasi paricle operators corresponding to basis defined by the vacuum of Eq.(1) i.e. $`q_I|vac=0=\stackrel{~}{q}_I^{}|vac`$ and finally $`\theta (k,\beta )`$ is the function for the thermal Bogoliubov transformation and is related to the number density function given as
$$\mathrm{sin}^2\theta (𝐤,\beta )=\frac{1}{\mathrm{exp}(\beta ϵ)+1}$$
(31)
where, $`ϵ(𝐤)`$ is the single particle energy given as $`ϵ(𝐤)=\sqrt{𝐤^2+m(k)^2}`$. In the presence of condensate the dynamical mass is given as $`m(k)=k\mathrm{tan}2h(k)`$ . As before the equal time propagator can be calculated including the temperature effects as
$`S_{\alpha \beta }(𝐱)`$ $`=`$ $`vac,\beta |{\displaystyle \frac{1}{2}}[q_\alpha (𝐱),\overline{q}_\beta (\mathrm{𝟎})]|vac,\beta `$ (32)
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{(2\pi )^3}}{\displaystyle e^{i𝐤𝐱}\mathrm{cos}2\theta (𝐤,\beta )\left[\mathrm{sin}2h(𝐤)\gamma 𝐤\mathrm{cos}2h(𝐤)\right]}`$ (33)
As in Section III, we shall take a gaussian ansatz for the condensate function $`\mathrm{sin}2h(𝐤)=\mathrm{exp}(R(T)^2𝐤^2/2`$, with the condensate scale parameter $`R`$, now being temperature dependent. In order to determine $`R(T)`$ or equivalently the ratio $`S(T)=R(T=0)/R(T)`$, we first evaluate our expression of the order parameter (the condensate value) at finite temperature. In terms of the dimensionless variable $`\eta =Rk`$, this is given as
$$\frac{\overline{q}q_T}{\overline{q}q_{T=0}}=S(T)^3\left[12\sqrt{\frac{2}{\pi }}e^{\eta ^2/2}\mathrm{sin}^2(z,\eta )\eta ^2𝑑\eta \right],$$
(35)
where $`\mathrm{sin}^2\theta (z,\eta )=1/(\mathrm{exp}(zϵ(\eta ))+1))`$, with $`z=\beta /R(T)`$ and $`ϵ(\eta )=\eta /\mathrm{cos}2h(\eta )`$.
We can obtain $`S(T)=R(T=0)/R(T)`$ if we know the temperature dependence of the order parameter on the left hand side of Eq.(35). As there are no phenomenological inputs for this, we shall proceed in the following manner to take the temperature dependence of the quark condensate. For low temperatures we shall take the results from chiral perturbation theory (CHPT) which is expected to be valid at least for small temperatures. For higher temperatures near the critical temperature, lattice simulations seem to yield the universal behaviour with a large correlation length associated with a second order phase transition for two flavor massless QCD. We shall use such a critical behaviour to consider the temperature dependence of the order parameter near the critical temperature.
For intermediate regime we shall take a smooth interpolation between the two. The resulting behaviour of the temperature dependance of the quark condensate and the ratio $`R(0)/R(T)`$ determined using Eq.(35) are shown in Fig.s 2.
With the temperature dependence of $`R(T)`$ known as above, the propagator function of eq.(LABEL:propt) now completely defined. The finite temperature correlation function is now given as, parallel to Eq.(20)
$$R(𝐱,T)=[Tr[S(𝐱,T)\mathrm{\Gamma }^{}S(𝐱,T)\mathrm{\Gamma }]+Tr[\mathrm{\Sigma }(𝐱))\mathrm{\Gamma }^{}\mathrm{\Sigma }(𝐱)\mathrm{\Gamma }]]_T.$$
(36)
The second term in the above corresponds to contributions from the fluctuations at finite temperature that can be written interms of the functions $`g_S(𝐱,T),g_V(𝐱,T)`$ as earlier but now being temperature dependent. We do not know how to calculate it except for a general property that the effect of the four point structure should decrease with temperature. We take here a simple ansatz for the temperature dependence of $`g_V`$ and $`g_S`$,
$$g_{S,V}(𝐱,T)=\left(\frac{\overline{q}q_T}{\overline{q}q_{T=0}}\right)^2g_{S,V}(𝐱,T=0)$$
(37)
The parameters $`\mu _i`$ are choosen to have the same values as of zero temperature while fitting the mesonic and baryonic correlation function. The normalised correlation functions $`R(𝐱,T)/R_0(𝐱,T=0)`$ are plotted in Fig.4a and Fig 4b respectively.
As expected (on physical grounds) the amplitude of the correlator decreases with increasing temperature. The peak of the vector correlator shifts towards the right after $`T=0.9T_c`$. We might remind ourselves that the position of the peak of the correlator is inversely proportional to the mass of the particle in the relevant channel.
To extract hadronic properties at finite temperature, the correlators are parametrised in terms of a spectral density function. This is a generalisation of eq.(29) to finite temperature and is given as ,
$$\rho ^V(s)=3\lambda _\rho ^2\delta (sM_\rho ^2)+\frac{3s}{4\pi ^2}\mathrm{tanh}\left[\frac{\sqrt{s}}{4T}\right]\theta (ss_o)+T^2S_\rho \delta (s)$$
(38)
$$\rho ^P(s)=\lambda _\pi ^2\delta (sM_\pi ^2)+\frac{3s}{8\pi ^2}\mathrm{tanh}\left[\frac{\sqrt{s}}{4T}\right]\theta (ss_o)$$
(39)
Eq.(39) corresponds to spectral density function for pseudoscalar channel. The last term in Eq.(38) The last term in Eq. (38) is the scattering term for soft thermal dissociations (mainly through pions), which exists only at finite temperature and can be taken as $`S_\rho \frac{T^2}{9}`$.
The mass, threshold and coupling are then extracted and the results are plotted in Fig. 4 for the vector and in the pseudoscalar channel. As can be seen from Fig.s3, with increase in temperature, the correlation functions have a lower peak indicating lack of correlations with temperature. In the vector channel the mass of the $`\rho `$ meson appears to decrease for temperatures beyond 120 MeV. The threshold for the continuum also decreases around the same temperature. The behaviour with temperature of these quantities is qualitatively similar to that found by Hatsuda et al. We have also plotted the temperature dependence of the coupling of the bound state to the current which decreases with temperature but rather slowly as compared to mass or the threshold for the continuum. The temperature dependence of these parameters can be used to calculate the lepton pair production rate from $`\rho `$ in the context of ultra relativistic heavy ion collision experiments to estimate vector meson mass shift in the medium.
In the pseudoscalar channel the mass remains almost constant till the critical temperature whereas the thershold and the coupling decrase with the temparature. We have found that in the pseudoscalar channel, the contribution to the correlation function mostly comes from the fluctuating fields. Further, the temperature behaviour as taken in Eq.(37) essentially does not shift the position of the peak whereas the magnitude of the correlator decreases. That is reflected in the above behaviour of the parameters in the pseudoscalar channel. We may note here that similar behaviour of pion mass becoming almost insensitive to temperature below the critical temperature was also observed in Ref. where correlation functions were calculated in a QCD motivated effective theory namely the Nambu- Jona Lasinio model.
## IV QCD vacuum at finite densities
Let us now go over to explore the ground state structure in high density QCD. Unlike high temperatures, rather little is known about QCD at finite baryon densities from first principles like lattice QCD simulation due to technical problems. However, different low energy models of QCD seem to indicate a rich phase structure in this domain. In particular the color superconductivity phase at high density has attracted much attention. Although it was known for quite sometime , recent studies indicated that the superconducting gap could be as large as 100 MeV and that has lead to an extensive literature on this subject in recent past regarding its consequences in neutron stars as well as for heavy ion collisions.
As per BCS theorem, an arbitrarily weak attractive interaction makes the fermi sea of quarks unsatble at high densities. In deed, the color current current interaction is attractive in the scalar $`(\overline{\mathrm{𝟑}})`$ and axial vector channel. This can lead to formation of Cooper pairs and associated superconductivity in the color space. To discuss it in a nonperturbative manner, we take the trial ansatz over the chiral condensed vacuum of Eq.(1) as
$$|\omega =\mathrm{exp}[\frac{1}{2}[q_r^{ia}(𝐤)^{}f(𝐤)q_r^{jb}(𝐤)^{}ϵ_{ij}ϵ_{3ab}+\stackrel{~}{q}_r^{ia}(𝐤)f_1(𝐤)\stackrel{~}{q}_r^{jb}(𝐤)ϵ_{ij}ϵ_{3ab}]d𝐤.]|vac$$
(40)
In the above, $`i,j`$ are flavor indices, $`a,b`$ are the color indices and $`r(=\pm 1/2)`$ is the spin index. We shall consider here two flavor and SU(3) color. Here we have also introduced two trial functions $`f(𝐤)`$ and $`f_1(𝐤)`$ respectively for the diquark and diantiquark channel. As may be noted the state constructed in Eq.(40) is spin singlet and is antisymmetric in color and flavor.
Next, to include the effect of temperature and density, we obtain the state at finite temperature and density $`|\omega (\beta ,\mu )`$ by a thermal Bogoliubov transformation over the state $`|\omega `$ using thermofield dynamics as in Eq.(30),
$$|\omega (\beta ,\mu )=\mathrm{exp}(q_I(𝐤)^{}\theta _{}(𝐤,\beta ,\mu )\underset{¯}{q}_I(𝐤)^{}+\stackrel{~}{q}_I(𝐤)\theta _+(𝐤,\beta ,\mu )\underset{¯}{\overset{~}{q}}_I(𝐤)d𝐤h.c.)|\omega $$
(41)
In the above, the ansatz functions $`\theta _\pm (𝐤,\beta ,\mu )`$, as before, will be related to quark and antiquark distributions. We might note here that the trial ansatz given in Eq.(41) actually invloves five functions - $`h(𝐤)`$, for the quark anti quark condensates, $`f(𝐤)`$ and $`f_1(𝐤)`$ describing respectively the diquark and diantiquark condensates and $`\theta _\pm (𝐤,\beta ,\mu )`$ to include the temperature and density effects. All these functions are to be obtained by minimising the thermodynamic potential. This will involve an assumption about the effective hamiltonian . For the purpose of illustration we shall consider a hamiltonian of Nambu-Jonalasinio type given as
$$=\psi ^{}(i𝜶\mathbf{})\psi +\frac{g^2}{2}J_\mu ^aJ^{\mu a}.$$
(42)
with $`J_\mu ^a=\overline{q}\gamma ^\mu T^aq`$. One can then calculate the energy functional $`ϵ=\omega (\beta ,\mu )||\omega (\beta ,\mu )`$ and the thermodynamic potential $`=ϵ(1/\beta )S\mu N`$, using the Bogoliubov technique. The details are given in Ref. . The free energy is a functional of all the five functions $`h,f,f_1,\theta _\pm `$. However with the point interaction of the the Hamiltonian of eq.(42) one can determine them. This results in two coupled gap equations to be solved in a self consistent manner and are given as
$$\frac{4g^2}{3}\frac{1}{(2\pi )^3}𝑑𝐤\frac{1}{\sqrt{𝐤^2+M^2}}\left(\frac{\xi _{}}{\omega _{}}tanh(\frac{\beta \omega _{}}{2})+\frac{\xi _+}{\omega _+}tanh(\frac{\beta \omega _+}{2})\right)=1$$
(44)
$$\frac{4g^2}{3}\frac{1}{(2\pi )^3}𝑑𝐤\left(\frac{tanh(\frac{\beta \omega _{}}{2})}{\omega _{}}+\frac{tanh(\frac{\beta \omega _+}{2})}{\omega _+}\right)=1$$
(45)
where, $`\omega _\pm =\sqrt{\mathrm{\Delta }^2+\xi _\pm ^2}`$ and $`\xi _\pm =(E\pm \nu )`$. Here, $`E`$ is the energy of the quasi paricles given as $`E=\sqrt{(k^2+M^2)}`$, and,$`\nu `$ is the chemical potential in presence of interaction given as $`\nu =\mu (4g^2/3)(N/12)`$, $`N`$ being the quark number density. Eq.(44) is the mass gap equation in presence of diquark condensates and Eq.(45) is the superconducting gap equation which is a relativistic generalisation of BCS gap equation . In the limit of no diquark condensates Eq.(44) reduces to that obtained in Ref. except for the numerical factors before the integrand. This is due to the fact that the approximation the the later case has been a mean field approximation unlike the case here where approximation lies only with the ansatz for the ground state. Eq.(44) without the diquark condensate contributions is also have the same structure as in Ref. in the limit of the formfactors intoduced in the later case reduces to a sharp cutoff in the momentum. Similarly in the limit of chiral condensate going to zero, the super conducting gap equation (45) is similar to that obtained in Ref. or Ref..
The solution of these equation at zero temperature is shown in Fig.(5) for the mass gap and in Fig.(6) for the superconducting gap. For the sake of comparision we have also plotted the mass gap without the diquark condensates. The coupling here taken as $`g^256GeV^2`$ and the cutoff$`\mathrm{\Lambda }0.67GeV`$. these values are taken so as to give the same transition temperature as in ref. With these couplings, the mass gap at zero temperature and density is about 490 MeV and the mass gap vanishes at fermi momentum $`k_f=\sqrt{\nu ^2M^2}=400MeV`$. The corresponding critical quark number density is about $`1.7/fm^3`$.
Presence of diquark condensate does not change these values very much. The diquark condensate increases with number density and becomes maximum of about 90 MeV beyond which the effect of the cutoff is felt and it vanishes for fermi momentum around 600 MeV. We also plot the equation of state (EOS) in Fig.7.
While plotting the EOS, we have added the bag constant $`ϵ_0`$, which is the difference in energy density of the perturbative vacuum and the nonperturbative vacuum with condensate at zero temperature and density, to the pressure. With the present parameters the bag constant turns out to be $`(200MeV)^4`$. The pressure has a cusp like structure and becomes negative at finite densities. The portion of the curve that goes down with $`k_f`$ corresponds to the nontrivial solution to the mass gap and the portion that increases with density correspond to zero mass solution. The negative pressure indicates mechanical instability and can have the interpretation that uniform nonzero density matter will break up in to droplets of finite density in which chiral symmetry is restored surrounded by empty space with zero pressure and density . It is tempting to identify the droplets of quark matter with nucleons within which density is nozero and $`\overline{q}q=0`$ – a fact reminscent of bag models . Nothing in the model however says that the droplets have quark number three. The equation of state does not change much in presence of diquark condensates. This should be expected because the effect of diquark condensate is small in the region where chiral condensate is non vanishing. Thus gross structural properties of the neutron stars are not likely to be affected by diquark condensates. However, cooling of neutron stars shall be expected to be very much affected by such a gap of about 100 MeV.
## V summary and discussions
We have looked into the structure of QCD vacuum in a nonperturbative manner with a variational ansatz. This has been done for zero temperature as well as finite temperature and densities. The input has been equal time algebra for the field operators and the ansatz for the ground state. At zero temperature, the Bogoliubov type pairing ansatz involving both quark and gluon condensates becomes energetically favourable beyond a critical coupling. It also gave some of the low energy hadronic properties like pion and proton charge raddi, proton,neutron magnetic moments and the bag constant to be around their phenomenological values for the coupling value of $`1.28`$. Thus $`\alpha _s=1.28`$ effectively corresponds to the QCD coupling constant for the vacuum configuration. With optimized renormalisation group equations, it has been seen that $`\alpha _s(Q)`$ does not go to infinity as $`Q`$ decreases below 300 MeV, but freezes to a constant value around unity. Our analysis seems to remind us of a similar situation.
We next evaluated the correlation functions of hadronic currents in such a condensed medium. However it appears that to have quantitative agreement with correlator phenomenology, particularly in the pseudoscalar channel, these cannot be described by the propagators alone but must of necessity have the fluctuations of the condensates. This may be looked upon as combination of two effects –(i) an effective way of incorporating gluon condensate effects and (ii)the existance of explicit four point structure in QCD vacuum. In some ways these fluctuations may be related to the “hidden contributions” discussed by Suryak .
It is worthwhile pointing out that OPE and our approach are based on intrinsically different assumptions. The former is an expansion which separates short distance (Wilson coefficients) and long distance (condensates) physics. In our method we assume an explicit vacuum structure in terms of quark condensate (two point function) plus an irreducible four point function. Having made an ansatz for the vacuum, we do not make any further approximation in the evaluation of the correlators. The approach is phenomenological in the sense that the values of the parameters in the four point function are chosen to reproduce the behaviour of the correlators. As emphasized by Shuryak and Sch$`\ddot{a}`$fer and Shuryak OPE is able to quantitatively describe the zero temperature pion correlation functions for small x (upto 0.25 fm) but underpredicts it for large x. In our work the agreement is quantitative with experimentally deduced mesonic correlation function and lattice results for the whole range. This covers the small x values, resonance region and large x domain where the correlator vanishes. We have however more parameters. To the extent that the large x behaviour of the pionic correlator depends on gluon condensates in OPE our parametrization would imply an effective way of including gluon condensate effects.
For studying the correlators at finite temperature, we assume that the two point as well as the four point function vanish at $`T=T_C`$. The parameters in the four point function are assumed constant at their $`T=0`$ values – no additional T dependence is given to them. We see that our results are broadly in agreement with those of Hatsuda etal . All the same, it is not possible to carry out a term by term comparison of our results with OPE results. This is because the assumptions are different in the two approaches. We recall that Hatsuda etal attribute the decrease in the $`\rho `$ mass to contributions coming from four point function. They also find variation in the gluon condensates to be less than 5% over the temperature region that is considered. In our work we do not include the gluon condensates and the four point structure vanishes as $`TT_C`$. This is not the case for gluon condensate ($`GG`$) in lattice calculations or in the dilute pion gas model of Ref. . Consequently, we may be tempted to infer that the decrease in rho mass in our model is due to “genuine” four point function effects (as in Hatsuda etal ) and not from the “effective” gluon condensate contribution. One, however, has to be very cautious. This is because in the present work the parameters in the four point function were kept constant at their $`T=0`$ values. It should be recalled that these values were determined so as to correctly describe the behaviour of the correlators (in particular pion) at T=0. Hence the parameters do reflect some effective gluon condensate effects. The results at finite temperature would certainly be modified if these parameters are given significant T dependence. Our parametrization is such that if the parameters decrease with T then the contribution to the correlator will decrease. The crucial question however still remains as to the behaviour of $`GG`$ for $`T<T_C`$. If it varies very little in this range then our assumption of the parameters remaining constant would be reasonable.
We would like to add here that the present analysis will be valid for temperatures below the critical temperature. Above the critical temperature there have been calculations essentially using finite temperature perturbative QCD in random phase approximations (RPA) . However, in the region above $`T_C`$, nonperurbative features have been known to exist from studies in lattice QCD simulations. In view of this, one may have to carry out a hard thermal loop calculation where a partial resummation is done. Alternatively, one may use other nonperturbative approaches such as QCD sum rules at finite temperature or RPA approach in an instanton liquid model for the QCD vacuum .
The vacuum structure at finite densities was looked into taking into account the possibility of diquark condensates in a Nambu-Jonalasinio model. As earlier, the approximation lies here only in the ansatz. The resulting gap equation for the superconducting gap is a relativistic generalisation of the BCS gap equation. The mass gap equation in the presence of superconducting gap is a new feature of the present calculation. Because of the point interaction as in $`_{int}`$ of Eq.(42)we could solve for the gap functions explicitly. In the presence of realistic potentials as in Ref.or Ref. one will have to solve an integral equation for the gap functions. Such a calculation is in progress . It will be interesting to see how the results for the superconducting gap would be affected in presence of quark antiquark condensates as compared to the resummed perturbative QCD calculations at finite densities . The equation of state here did not change very much. Thus the global properties of neutron stars shall not be affected in an appreciable manner. However, a gap of about 100 MeV can have its implications on neutron star cooloing , magnetic fields of pulsars and their thermal evolution . Future theoretical studies of QCD at finite baryon densities may reveal to which extent these newly established features of quark matter considered in effective models shall have their correspondence in a more complete treatment of QCD at finite baryonic density and shed light on whether one may expect any distinguishing feature in the global properties of neutron/quark stars.
## ACKNOWLEDGMENTS
I would like to thank S.P. Misra, Amruta Mishra, J.C. Parikh and Varun Sheel for an enjoyble and fruitful collaboration over the years.
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# Cohomology of congruence subgroups of SL₄(ℤ)
## 1. Introduction
### 1.1.
Let $`n1`$, and let $`\mathrm{\Gamma }`$ be a congruence subgroup of $`\mathrm{SL}_n()`$ of level $`N`$. Let $`S_N`$ be the subsemigroup of the integral matrices in $`\mathrm{GL}_n()`$ such that $`(\mathrm{\Gamma },S_N)`$ is a Hecke pair.
We denote by $`(N)`$ the $``$-algebra of double cosets $`\mathrm{\Gamma }S_N\mathrm{\Gamma }`$. This algebra acts on the cohomology and homology of $`\mathrm{\Gamma }`$ with any coefficient $`S_N`$-module. When a double coset is acting, we call the map defined by its action a Hecke operator. In this paper we will work only with the trivial coefficient module $``$.
Let $`l`$ be a prime not dividing $`N`$, and let $`D(l,k)`$ be the diagonal matrix with diagonal $`(1,\mathrm{},1,l,\mathrm{},l)`$, where the number of $`l`$’s is $`k`$. Then $`(N)`$ contains all double cosets of the form $`\mathrm{\Gamma }D(l,k)\mathrm{\Gamma }`$, and we denote the corresponding Hecke operator by $`T(l,k)`$. Fix a prime $`p`$ not dividing $`N`$, and an embedding of $`_p`$ into $``$. Let $`G_{}=\mathrm{Gal}(\overline{}/)`$.
###### Definition 1.2.
Let $`𝒱`$ be an $`(N)`$-module and suppose $`\beta 𝒱`$ is an eigenclass for the action of $`(N)`$. For $`l`$ prime to $`N`$, write $`T(l,k)(\beta )=a(l,k)\beta `$ , where the $`a(l,k),k=0,\mathrm{},n`$ are algebraic integers. Let $`\rho `$ be a continuous semisimple representation $`\rho :G_{}\mathrm{GL}_n(_p)`$, unramified outside $`pN`$, such that
(1)
$$\underset{k}{}(1)^kl^{k(k1)/2}a(l,k)X^k=det(I\rho (\mathrm{Frob}_l)X)$$
for all $`l`$ not dividing $`pN`$. Then we shall say that $`\rho `$ is *attached* to $`\beta `$.
For example, theorems of Eichler, Shimura and Deligne imply that if $`n=2`$ and $`𝒱`$ is the Hecke-module of classical holomorphic modular cuspforms for $`\mathrm{SL}_2()`$, then there always exists a $`\rho `$ attached to any Hecke eigenform $`\beta 𝒱`$.
Standard conjectures (for example in ) state that if $`𝒱`$ is the cuspidal cohomology of $`\mathrm{\Gamma }`$ with trivial coefficients, then any Hecke eigenclass should have an attached $`p`$-adic Galois representation. In fact, one has the stronger conjecture that there should be an attached motive (cf. ). We can also extend the conjecture of to include all the cohomology—in principle, the theory of Eisenstein series should allow one to reduce the extended conjecture to the one for cuspidal cohomology.
In a series of papers , the first author with a number of coworkers has tested this conjecture and a mod $`p`$ variant of it when $`n=3`$. Other tests for $`n=3`$ can be found in the work of van Geemen and Top with van der Kallen and Verberkmoes . The purpose of this paper is to make the first computational tests of this conjecture for $`n=4`$.
### 1.3.
We work here with $`\mathrm{\Gamma }=\mathrm{\Gamma }_0(N)`$, defined as the subgroup of $`\mathrm{SL}_n()`$ consisting of the matrices with last row congruent to $`(0,\mathrm{},0,)`$ modulo $`N`$.
Conceptually, there is no big difference in computing the cohomology of $`\mathrm{\Gamma }`$ between the cases of $`n=3`$—or even $`n=2`$—and the case of general $`n`$. One needs to write down a simplicial complex $`C`$ homotopic to the chains of a $`B\mathrm{\Gamma }(1)`$-space, and compute the cohomology.
When $`n=2`$ or $`3`$, the most interesting part of cohomology, namely the cuspidal part, occurs in the top dimension of $`C`$. Therefore to compute it all we have to compute is the cokernel of a coboundary map. However, when $`n=4`$ the cuspidal cohomology occurs in dimensions 4 and 5, whereas the virtual cohomological dimension of $`\mathrm{\Gamma }`$, and hence smallest possible dimension of $`C`$, is 6. So now the cohomology is a subquotient, which adds considerably to the complexity of the computer programs. The cuspidal cohomology in degree 4 is dual to that in degree 5, so we concentrate on computing the latter in this paper.
But the big difference between $`n=2,3`$ and $`n=4`$ occurs when we try to compute the Hecke action on the cohomology. In the top dimension we can use the Ash-Rudolph algorithm or its variants , as was done for $`n=3`$ in the works cited above. However, for $`n=4`$, where we look just below the top dimension, a brand new idea was necessary. This is due to the second author, and is the subject of . Thus this paper is also a test of the algorithms proposed in , and they pass with flying colors. It is an open question whether the algorithms of terminate in a finite number of steps; in practice, though, they always terminate quickly, and we used them here without problems.
Our computations of $`H^5(\mathrm{\Gamma }_0(N);)`$, detailed in §5, were made for $`N53`$.<sup>1</sup><sup>1</sup>1Actually, to avoid numerical instability in floating-point computations, we replace $``$ with the finite field $`𝔽_{31991}`$ (cf. §5.3). No cusp forms were discovered, but some interesting phenomenology of the boundary cohomology was observed, as discussed in §6. This leads to some open questions about the cohomology of the boundary, which are discussed there. Since we have not completed the Hecke computations for some high levels near $`53`$, it is possible that we do have a cusp form that we haven’t yet identified as such. There is a bound due to Fermigier that states that if $`N<31`$, there cannot be any cuspidal cohomology.
For high levels our computations produced large sparse matrices, as large as $`110464\times 30836`$ for level $`48`$. To perform linear algebra with these matrices, we used a version of the Lanczos algorithm mod $`p`$, in the spirit of LaMacchia-Odlyzko (see also ).
### 1.4.
We thank Eric Conrad for some assistance with programming, and Peter Woit for excellent computing support. We thank David Ginzburg for conversations at the beginning of this project.
## 2. Background
### 2.1.
Let $`V`$ be the $``$-vector space of all symmetric $`n\times n`$ matrices, and let $`CV`$ be the cone of positive-definite matrices. Then the group $`G=\mathrm{SL}_n()`$ acts on $`C`$ on the left by $`(g,c)gcg^t`$, and the stabilizer of any given point is isomorphic to $`\mathrm{SO}_n`$.
Let $`X`$ be $`C`$ mod homotheties. The $`G`$-action on $`C`$ commutes with the homotheties and induces a transitive $`G`$-action on $`X`$. The stabilizer of any given point of $`X`$ is again isomorphic to $`\mathrm{SO}_n`$. After choosing a basepoint, we may identify $`X`$ with the global riemannian symmetric space $`\mathrm{SL}_n()/\mathrm{SO}_n`$, a contractible, noncompact, smooth manifold of real dimension $`d=n(n+1)/21`$.
The group $`\mathrm{SL}_n()`$ acts on $`X`$ via the $`G`$-action, and does so properly discontinuously. Hence if $`\mathrm{\Gamma }\mathrm{SL}_n()`$ is any finite-index subgroup, the quotient $`\mathrm{\Gamma }\backslash X`$ is a real noncompact manifold except for at most finitely many quotient singularities. We may then identify the complex group cohomology $`H^{}(\mathrm{\Gamma };)`$ with $`H^{}(\mathrm{\Gamma }\backslash X;)`$. Although the dimension of $`\mathrm{\Gamma }\backslash X`$ is $`d`$, it can be shown that $`H^i(\mathrm{\Gamma }\backslash X;)=0`$ if $`i>dn+1`$ \[12, Theorem 11.4.4\]. The number $`\nu =dn+1`$ is called the *virtual cohomological dimension* of $`\mathrm{\Gamma }`$.
In this paper we will always take $`\mathrm{\Gamma }`$ to be the congruence subgroup $`\mathrm{\Gamma }_0(N)`$ of matrices whose last row is congruent to $`(0,\mathrm{},0,)modN`$.
### 2.2.
Recall that a point in $`^n`$ is said to be *primitive* if the greatest common divisor of its coordinates is $`1`$. In particular, a primitive point is nonzero. Let $`𝒫^n`$ be the set of primitive points. Any $`v𝒫`$, written as a column vector, determines a rank-one symmetric matrix $`q(v)\overline{C}`$ by $`q(v)=vv^t`$. The *Voronoǐ polyhedron* $`\mathrm{\Pi }`$ is the closed convex hull of the points $`q(v)`$, as $`v`$ ranges over $`𝒫`$.
Note that, by construction, $`\mathrm{SL}_n()`$ acts on $`\mathrm{\Pi }`$. The cones over the faces of $`\mathrm{\Pi }`$ form a fan $`𝒱`$ that induces a $`\mathrm{\Gamma }`$-admissible decomposition of $`C`$ \[2, p. 117\]. Essentially, this means that $`\mathrm{\Gamma }`$ acts on $`𝒱`$; that each cone is spanned by a *finite* collection of points $`q(v)`$ where $`v𝒫`$; and that there are only finitely many $`\mathrm{\Gamma }`$-orbits in $`𝒱`$. The fan $`𝒱`$ provides a reduction theory for $`C`$ in the following sense: any point $`xC`$ is contained in a unique $`\sigma (x)𝒱`$, and the set $`\{\gamma \mathrm{SL}_n()\gamma \sigma (x)=\sigma (x)\}`$ is finite.
### 2.3.
We summarize facts about the well-rounded retract of . There is a deformation retraction $`CC`$ that is equivariant under the actions of both $`\overline{\mathrm{\Gamma }}=\mathrm{SL}_n()`$ and the homotheties. Its image modulo homotheties is the *well-rounded retract* $`W`$ in $`X`$. The well-rounded retract is contractible, since it is a deformation retract of the contractible space $`X`$. Hence the cohomology of $`\mathrm{\Gamma }\backslash X`$ with coefficients in $``$ is canonically isomorphic to the equivariant cohomology $`H_\mathrm{\Gamma }^i(W;)`$ where $`\mathrm{\Gamma }`$ acts trivially on the coefficient module $``$. This is in turn canonically isomorphic to the complex cohomology $`H^i(\mathrm{\Gamma }\backslash W;)`$, since $``$ has characteristic zero, which moreover is isomorphic to $`H^i(\mathrm{\Gamma };)`$. We will focus on computation of the equivariant cohomology. The dimension of $`W`$ equals the virtual cohomological dimension $`\nu `$, and the quotient $`\mathrm{\Gamma }\backslash W`$ is compact.
The well-rounded retract $`W`$ is naturally a locally finite cell complex, the cells being convex polytopes in $`V`$. The group $`\mathrm{SL}_n()`$ preserves the cell structure, and the stabilizer of each cell in $`\mathrm{SL}_n()`$ is finite. The theory of cores and co-cores in \[7, Chapter 2\] shows that the cells in $`W`$ are in a one-to-one, inclusion-reversing correspondence with the cones in the Voronoǐ fan $`𝒱`$. By abuse of notation, the cell in $`W`$ corresponding to $`\sigma `$ will still be denoted $`\sigma `$.
### 2.4.
One can give a more precise description of the combinatorics of the cones in $`𝒱`$ and the cells in $`W`$. To each $`\sigma 𝒱`$, we define the set
$$M(\sigma )=\{v𝒫q(v)\text{ is a vertex of the face of }\mathrm{\Pi }\text{ generating }\sigma \text{.}\}$$
We associate the same set $`M(\sigma )`$ to the corresponding cell in $`W`$, and call $`M(\sigma )`$ the set of *minimal vectors of* $`\sigma `$ (because of how $`W`$ is constructed in ). Since $`\mathrm{\Pi }`$ is the convex hull of the $`q(v)`$’s, it is clear that $`\mu :\sigma M(\sigma )`$ is an inclusion-preserving (respectively, inclusion-reversing) bijection between the cones in $`𝒱`$ (resp., cells in $`W`$) and a collection of finite subsets of $`𝒫`$. In principle, this reduces the study of the combinatorics of $`𝒱`$ and $`W`$ to the study of the image of $`\mu `$. For instance, face relations $`\tau \sigma `$ in $`𝒱`$ are read off from subset relations $`M(\tau )M(\sigma )`$. In practice, determining the image of $`\mu `$ calls for explicit computations with real quadratic forms, computations whose difficulty grows exponentially as a function of $`n`$. The computations have been carried out completely for $`n5`$ by various authors.
For the rest of this subsection, we set $`n=4`$ and give more details. We state results for the well-rounded retract; these imply their analogues for $`𝒱`$. The image of $`\mu `$ was computed independently by and (in essence) . The cells of $`W`$ fall into eighteen equivalence classes modulo $`\mathrm{SL}_n()`$. Let $`T`$ (for “type”) be a variable running through these eighteen classes. This partitions the set of cells of $`W`$ into eighteen pieces called the $`W_T`$. Any $`\sigma W_T`$ is said to be *of type $`T`$.* In each $`W_T`$, we fix one representative cell $`\sigma _T`$, the *standard cell of type $`T`$*. The $`M(\sigma _T)`$’s are written down explicitly in .<sup>2</sup><sup>2</sup>2Since $`vM(\sigma )vM(\sigma )`$, it is customary to write down only one member of the pair $`\pm v`$. This determines the image of $`\mu `$, since the image is the union of the $`\mathrm{SL}_n()`$-translates of the eighteen $`M(\sigma _T)`$’s.
### 2.5.
To compute the action of the Hecke operators on cohomology, the well-rounded retract is insufficient, since the operators do not act cellularly. To ameliorate this, we use the sharbly complex. The material in this subsection closely follows .
###### Definition 2.6.
The *sharbly complex* is the chain complex $`\{S_{},\}`$ given by the following data:
1. For $`k0`$, $`S_k`$ is the module of formal $``$-linear combinations of basis elements $`𝐮=[v_1,\mathrm{},v_{n+k}]`$, where each $`v_i𝒫`$, mod the relations:
1. If $`\tau `$ is a permutation on $`(n+k)`$ letters, then
$$[v_1,\mathrm{},v_{n+k}]=\mathrm{sgn}(\tau )[\tau (v_1),\mathrm{},\tau (v_{n+k})],$$
where $`\mathrm{sgn}(\tau )`$ is the sign of $`\tau `$.
2. If $`q=\pm 1`$, then
$$[qv_1,v_2\mathrm{},v_{n+k}]=[v_1,\mathrm{},v_{n+k}].$$
3. If the rank of the matrix $`(v_1,\mathrm{},v_{n+k})`$ is less than $`n`$, then $`𝐮=0`$.
2. The boundary map $`:S_kS_{k1}`$ is
$$[v_1,\mathrm{},v_{n+k}]\underset{i=1}{\overset{n+k}{}}(1)^i[v_1,\mathrm{},\widehat{v_i},\mathrm{},v_{n+k}].$$
The basis elements $`𝐮=[v_1,\mathrm{},v_{n+k}]`$ are called *$`k`$-sharblies*. By abuse of notation, we will often use the same symbol $`𝐮`$ to denote a $`k`$-sharbly and the $`k`$-sharbly chain $`1𝐮`$. The obvious left action of $`\mathrm{\Gamma }`$ on $`S_{}`$ commutes with $``$.
For any $`k0`$, let $`(S_k)_\mathrm{\Gamma }`$ be the module of $`\mathrm{\Gamma }`$-coinvariants. This is the quotient of $`S_k`$ by the relations of the form $`\gamma 𝐮𝐮`$, where $`\gamma \mathrm{\Gamma }`$, $`𝐮S_k`$. This is also a complex with the induced boundary, which we denote by $`_\mathrm{\Gamma }`$. It is known (cf. ) that $`H^{\nu k}(\mathrm{\Gamma };)`$ is naturally isomorphic to $`H_k((S_{})_\mathrm{\Gamma })`$.
Let $`𝐮=[v_1,\mathrm{},v_{n+k}]`$ be a $`k`$-sharbly. Let $`𝐮`$ be
$$Max|det(v_{i_1},\mathrm{},v_{i_n})|,$$
where the maximum is taken over all $`n`$-fold subsets $`\{i_1,\mathrm{},i_n\}\{1,\mathrm{},n+k\}`$. Note that this quantity is well-defined mod the relations in Definition 2.6. We extend this notion to sharbly chains $`\xi =n(𝐮)𝐮`$ by setting $`\xi `$ to be the maximum of $`𝐮`$, as $`𝐮`$ ranges over all sharblies in the support of $`\xi `$. We say that $`\xi `$ is *reduced* if $`\xi =1`$. It is known (cf. ) that for $`\mathrm{\Gamma }\mathrm{SL}_4()`$, the group $`H^5(\mathrm{\Gamma };)`$ is spanned by reduced $`1`$-sharbly cycles.
### 2.7.
Since the generators of the sharbly complex are indexed by sets of primitive vectors, it is clear that there is a close relationship between $`S_{}`$ and the chain complex associated to $`W`$, although of course $`S_{}`$ is much bigger. Both complexes compute $`H^{}(\mathrm{\Gamma };)`$. We refer to for a discussion of this, phrased in terms of the fan $`𝒱`$. The main advantage of $`(S_{})_\mathrm{\Gamma }`$ is that it admits a Hecke action. Specifically, let $`\xi =n(𝐮)𝐮`$ be a sharbly cycle mod $`\mathrm{\Gamma }`$, and consider the Hecke operator $`T(l,k)`$ associated to the double coset $`\mathrm{\Gamma }D(l,k)\mathrm{\Gamma }`$ (cf. §1.1). Write
$$\mathrm{\Gamma }D(l,k)\mathrm{\Gamma }=\underset{g\mathrm{\Omega }}{}\mathrm{\Gamma }g,$$
a finite (disjoint) union. Then
(2)
$$T(l,k)(\xi )=\underset{g\mathrm{\Omega },𝐮}{}n(𝐮)g𝐮.$$
Since $`\mathrm{\Omega }\mathrm{SL}_n()`$ in general, the Hecke-image of a reduced sharbly isn’t usually reduced.
## 3. Implementation details
### 3.1.
We state our results for general $`n`$ as much as possible, though our main case of interest is $`n=4`$. We have working programs for $`n4`$. Though we focus on $`\mathrm{SL}_n()`$, analogous results hold for $`\mathrm{GL}_n()`$, and we have working programs for both $`\mathrm{SL}`$ and $`\mathrm{GL}`$.
Section 3 is very technical. The reader may wish to skip to §4 or §5.
### 3.2.
Let $`\overline{\mathrm{\Gamma }}=\mathrm{SL}_n()`$. Recall that $`W_T`$ is the $`\overline{\mathrm{\Gamma }}`$-orbit of cells of type $`T`$ in $`W`$. Let $`\sigma _T`$ be a fixed representative cell in $`W_T`$. The stabilizer in $`\overline{\mathrm{\Gamma }}`$ of $`\sigma _T`$ is denoted $`\overline{\mathrm{\Gamma }}_{\sigma _T}`$, or $`\overline{\mathrm{\Gamma }}_T`$ for short. This is a finite group that is straightforward to compute, since the minimal vectors $`M(\sigma _T)`$ are known. Our program maintains a database of the $`\overline{\mathrm{\Gamma }}_T`$.
Standard facts about stabilizers give the following:
###### Proposition 3.3.
There is a one-to-one correspondence between cells $`\sigma W_T`$ and cosets $`\overline{\mathrm{\Gamma }}/\overline{\mathrm{\Gamma }}_T`$, given by $`\gamma \sigma _T\gamma \overline{\mathrm{\Gamma }}_T`$ for any $`\gamma \overline{\mathrm{\Gamma }}`$ such that $`\sigma =\gamma \sigma _T`$.
Under the smaller group $`\mathrm{\Gamma }`$, the $`\overline{\mathrm{\Gamma }}`$-orbit $`W_T`$ breaks up into suborbits. If $`\mathrm{\Gamma }`$ were a torsion-free group, $`\mathrm{\Gamma }\backslash W`$ would be a finite cell complex, its cells would be given exactly by the $`\mathrm{\Gamma }`$-suborbits, and we could compute $`H^i(\mathrm{\Gamma }\backslash W)`$ by the standard methods for cell complexes. In our case, $`\mathrm{\Gamma }`$ is not torsion-free, but $`\mathrm{\Gamma }\backslash W`$ can be thought of as a finite “orbifold cell” complex, whose elements are the quotients of cells by finite groups; the $`\mathrm{\Gamma }`$-suborbits are in one-to-one correspondence with the orbifold cells.
The goal of this subsection is to understand the $`\mathrm{\Gamma }`$-suborbits in terms of the actions of the $`\overline{\mathrm{\Gamma }}_T`$ on finite projective spaces. By $`^{n1}=^{n1}(/N)`$, we mean the set of vectors $`(x_1,\mathrm{},x_n)(/N)^n`$ that are primitive in the sense that the ideal $`(x_1,\mathrm{},x_n)`$ in $`/N`$ is $`(1)`$, modulo the equivalence relation given by scalar multiplication by the units $`(/N)^\times `$ of $`/N`$. When $`N`$ is a prime, this is the usual projective space over the field of $`N`$ elements. As usual, the equivalence class of the vector $`(x_1,\mathrm{},x_n)`$ is denoted $`𝐚=[x_1:\mathrm{}:x_n]`$. We view these $`n`$-tuples as rows rather than columns; $`\overline{\mathrm{\Gamma }}`$ acts on the right on $`^{n1}(/N)`$ in the obvious way.
We define the bottom row map $`𝔟:\overline{\mathrm{\Gamma }}^{n1}(/N)`$ as follows. For a matrix $`\gamma \overline{\mathrm{\Gamma }}`$, the bottom row of $`\gamma `$ is a primitive vector in $`^n`$. Let $`𝔟(\gamma )`$ be the equivalence class of this image in $`^{n1}(/N)`$.
###### Lemma 3.4.
The bottom row map $`𝔟:\overline{\mathrm{\Gamma }}^{n1}(/N)`$ induces a bijection between $`\mathrm{\Gamma }\backslash \overline{\mathrm{\Gamma }}`$ and $`^{n1}(/N)`$, given by
$$\mathrm{\Gamma }\gamma 𝔟(\gamma ).$$
The map is equivariant for the right action of $`\overline{\mathrm{\Gamma }}`$.
###### Proof.
It is a standard fact that a vector in $`^n`$ is primitive if and only if it is the bottom row of some element of $`\overline{\mathrm{\Gamma }}`$. This implies $`𝔟`$ is surjective and that $`𝔟^1([0:\mathrm{}:0:0:1])=\mathrm{\Gamma }`$. The rest is clear. ∎
We can now describe the $`\mathrm{\Gamma }`$-orbits of cells in each $`W_T`$.
###### Proposition 3.5.
The $`\mathrm{\Gamma }`$-orbits of cells in $`W_T`$ are in one-to-one correspondence with the orbits $`O`$ of the right $`\overline{\mathrm{\Gamma }}_T`$-action on $`^{n1}`$.
###### Proof.
We have $`\mathrm{\Gamma }\backslash W_T=\mathrm{\Gamma }\backslash \overline{\mathrm{\Gamma }}/\overline{\mathrm{\Gamma }}_T`$ by Proposition 3.3, and this equals $`^{n1}/\overline{\mathrm{\Gamma }}_T`$ by Lemma 3.4. ∎
The first step of our computer program is to determine, for each type $`T`$, the decomposition of $`^{n1}`$ into right $`\overline{\mathrm{\Gamma }}_T`$-orbits. Since we are primarily studying $`H_\mathrm{\Gamma }^i(W;)`$ for $`i=5,6`$, we only need to work with the $`T`$ representing cells of dimensions 4, 5 and 6.
We note the following fact, whose proof is immediate.
###### Lemma 3.6.
Let $`\gamma _0\overline{\mathrm{\Gamma }}`$, and let $`𝐚=𝔟(\gamma _0)`$ in $`^{n1}`$. Then the stabilizer of $`𝐚`$ under the right action of $`\overline{\mathrm{\Gamma }}`$ is $`\gamma _0^1\mathrm{\Gamma }\gamma _0`$.
Let $`\sigma W`$ be a cell of type $`T`$, with $`\sigma =\gamma _0\sigma _T`$. Its stabilizer in $`\mathrm{\Gamma }`$, denoted $`\mathrm{\Gamma }_\sigma `$, is clearly
(3)
$$\mathrm{\Gamma }_\sigma =(\gamma _0\overline{\mathrm{\Gamma }}_T\gamma _0^1)\mathrm{\Gamma }.$$
Lemma 3.6 implies
###### Lemma 3.7.
The group $`\gamma _0^1\mathrm{\Gamma }_\sigma \gamma _0`$ is the subgroup of $`\overline{\mathrm{\Gamma }}_T`$ that preserves $`𝐚=𝔟(\gamma _0)`$ under the right action on $`^{n1}`$.
### 3.8.
In this subsection, we fix orientations on the cells $`\sigma W`$. It’s necessary to be extremely careful—mistakes in orientation are easy to make and will ruin the computations. The price to pay is to sort through the details of the action of $`\overline{\mathrm{\Gamma }}`$.
Recall that $`\overline{\mathrm{\Gamma }}_T`$ is the stabilizer of $`\sigma _T`$ in $`\overline{\mathrm{\Gamma }}`$. For each $`T`$, there is an orientation character $`\overline{\mathrm{\Gamma }}_T\{\pm 1\}`$ indicating whether or not $`\gamma \overline{\mathrm{\Gamma }}_T`$ preserves the orientation on $`\sigma _T`$. Our program stores the values of these characters along with $`\overline{\mathrm{\Gamma }}_T`$. Let $`\overline{\mathrm{\Gamma }}_T^+`$ be the subgroup of $`\overline{\mathrm{\Gamma }}_T`$ where the orientation is $`+1`$.
###### Remark 3.9.
To compute the value of the character at $`\gamma `$, we determine (i) how $`\gamma `$ acts on the orientation of the cone $`C`$, and divide by (ii) how $`\gamma `$ acts on the orientation of the Voronoǐ cone dual to $`\sigma _T`$. Dividing works because $`C`$ is locally the direct product of the cell and its dual Voronoǐ cone. As for (i), every element of $`\overline{\mathrm{\Gamma }}`$ acts by $`+1`$ on the orientation of $`C`$, since $`\overline{\mathrm{\Gamma }}`$ is a subgroup of the connected group $`\mathrm{SL}_n()`$. In this paper, where $`n=4`$ and $`dim\sigma >0`$, it turns out that all the dual Voronoǐ cones are simplicial; the sign in (ii) is the sign of the permutation that $`\gamma `$ effects on the bounding rays $`q(v)`$ of the cone, which is easily computed.
Let $`O`$ be a right $`\overline{\mathrm{\Gamma }}_T`$-orbit in $`^{n1}`$. We call $`O`$ non-orientable if for some (which implies every) $`𝐚O`$, there exists $`\gamma \overline{\mathrm{\Gamma }}_T\overline{\mathrm{\Gamma }}_T^+`$ with $`𝐚\gamma =𝐚`$. Otherwise, we call $`O`$ orientable. These notions depend on $`T`$, though we usually leave $`T`$ out of the notation.
If $`O`$ is orientable, fix some $`𝐚_0O`$. Define the orientation number of $`𝐚O`$ to be $`+1`$ (resp., $`1`$) according as $`𝐚=𝐚_0\gamma `$ for some $`\gamma \overline{\mathrm{\Gamma }}_T^+`$ (resp., $`\gamma \overline{\mathrm{\Gamma }}_T\overline{\mathrm{\Gamma }}_T^+`$). The orientation number is well-defined precisely because $`O`$ is orientable. Again, the notions depend on the choice of $`𝐚_0`$, though we leave $`𝐚_0`$ out of the notation.
Let $`\gamma \overline{\mathrm{\Gamma }}_T`$, and let $`\rho `$ be any cell of $`W`$ with any given orientation. Since $`\gamma `$ acts by diffeomorphisms on $`C`$, it carries the orientation on $`\rho `$ to some orientation on the cell $`\gamma \rho `$. We write
(4)
$$(\gamma )_{}(\rho )$$
to denote $`\gamma \rho `$ together with this orientation. Clearly $`(\gamma )_{}`$ is functorial, and preserves the relative orientation of $`\rho ,\tau `$ whenever $`\tau `$ is a codimension-one face of $`\rho `$.
Once and for all, fix orientations on the standard cells $`\sigma _T`$. We can now put orientations on all the cells of $`W`$.
###### Definition 3.10.
Let $`\sigma `$ be a cell in $`W_T`$ with $`\sigma =\gamma _0\sigma _T`$. Let $`O`$ be the right $`\overline{\mathrm{\Gamma }}_T`$-orbit in $`^{n1}`$ corresponding to $`\sigma `$ as in Proposition 3.5. If $`O`$ is orientable, we give $`\sigma `$ the orientation
(5)
$$(\text{orientation number of }𝐚)(\gamma _0)_{}(\sigma _T).$$
If $`O`$ is non-orientable, we give $`\sigma `$ an arbitrary orientation.
###### Proposition 3.11.
The quantity in (5) is well-defined.
###### Proof.
Let $`\sigma `$, $`\gamma _0`$, $`O`$, and $`𝐚_0`$ be as in Definition 3.10, with $`O`$ assumed orientable. Assume $`\sigma =\gamma _1\sigma _T`$ as well as $`\gamma _0\sigma _T`$. Let $`𝐚_1=𝔟(\gamma _1)`$. Then $`\gamma _1^1\gamma _0\overline{\mathrm{\Gamma }}_T`$. By definition of $`\overline{\mathrm{\Gamma }}_T^+`$, $`(\gamma _0)_{}(\sigma _T)=(\gamma _1)_{}(\sigma _T)`$ if and only if $`\gamma _1^1\gamma _0\overline{\mathrm{\Gamma }}_T^+`$. On the other hand, $`𝐚_1\gamma _1^1\gamma _0=[0:\mathrm{}:0:1]\gamma _0=𝐚`$, so $`𝐚`$ and $`𝐚_1`$ have the same orientation number if and only if $`\gamma _1^1\gamma _0\overline{\mathrm{\Gamma }}_T^+`$. ∎
We must understand how $`\overline{\mathrm{\Gamma }}`$ and $`\mathrm{\Gamma }`$ act on the orientations we have just chosen. The following lemma is immediate from (4) and Definition 3.10.
###### Lemma 3.12.
If $`\sigma =\gamma _0\sigma _T`$ corresponds to an orientable $`O`$, then $`\gamma _0`$ carries the chosen orientation of $`\sigma _T`$ to the chosen orientation of $`\sigma `$ times the orientation number of $`𝐚`$.
Here is a more general statement.
###### Proposition 3.13.
Let $`\sigma =\gamma _0\sigma _T`$ and $`\sigma _1=\gamma _1\sigma `$, for some $`\gamma _0,\gamma _1\overline{\mathrm{\Gamma }}`$. Let $`𝐚_0=𝔟(\gamma _0)`$ and $`𝐚_1=𝔟(\gamma _1\gamma _0)`$. Let $`O_1,O`$ be the right $`\overline{\mathrm{\Gamma }}_T`$-orbits in $`^{n1}`$ containing $`𝐚_1,𝐚_0`$; assume these orbits are both orientable. Then $`\gamma _1`$ carries $`\sigma `$ to $`\sigma _1`$ while multiplying the orientations by
(6)
$$(\text{orien. number of }𝐚)(\text{orien. number of }𝐚_1).$$
###### Proof.
Apply Lemma 3.12 twice. ∎
Fortunately, (6) becomes trivial when we consider $`\mathrm{\Gamma }`$ as opposed to $`\overline{\mathrm{\Gamma }}`$.
###### Proposition 3.14.
Let $`\sigma _1=\gamma _1\sigma `$ for $`\gamma _1\mathrm{\Gamma }`$. Let $`𝐚_1,𝐚,O_1,O`$ be as in Proposition 3.13, both orbits being assumed orientable. Then $`\gamma _1`$ carries $`\sigma `$ to $`\sigma _1`$ with orientations matching.
###### Proof.
We have $`[0:\mathrm{}:0:1]\gamma _1=[0:\mathrm{}:0:1]`$ by the definition of $`\mathrm{\Gamma }`$. Hence $`𝐚_1=([0:\mathrm{}:0:1]\gamma _1)\gamma _0=[0:\mathrm{}:0:1]\gamma _0=𝐚`$. Thus the expression in (6) is a square either of $`+1`$ or of $`1`$. ∎
### 3.15.
We compute the equivariant cohomology $`H_\mathrm{\Gamma }^{}(W;)`$ using a spectral sequence, following the exposition in \[13, VII.7–8\]. (The spectral sequence there is for equivariant homology; we make the appropriate modifications for cohomology.)
Let $`\sigma `$ be any cell in the well-rounded retract $`W`$. Recall that $`\mathrm{\Gamma }_\sigma `$ is the stabilizer of $`\sigma `$ in $`\mathrm{\Gamma }`$. Let $`_\sigma `$ be the $`\mathrm{\Gamma }_\sigma `$-module where $`\gamma \mathrm{\Gamma }_\sigma `$ acts by $`+1`$ if $`\gamma `$ preserves the orientation of $`\sigma `$ and by $`1`$ if it does not.
For each $`i`$, let $`W_{(i)}`$ be a fixed set of representatives of the $`\mathrm{\Gamma }`$-orbits of the cells in $`W`$ of dimension $`i`$. The $`E_1`$ term of the spectral sequence is
(7)
$$E_1^{i,j}=\underset{oW_{(i)}}{}H^j(\mathrm{\Gamma }_o;_o).$$
These cells $`o`$ (omicron) are in one-to-one correspondence with the $`O`$ of Proposition 3.5, as $`T`$ runs through the types of cells of dimension $`i`$. Because $``$ is a field of characteristic zero, all the terms in (7) vanish when $`j0`$. In particular, the spectral sequence collapses at $`E_2`$. The term $`H^0(\mathrm{\Gamma }_o;_o)`$ is the subset of $`\mathrm{\Gamma }_o`$-invariants in the module $`_o`$.
###### Proposition 3.16.
For any cell $`\sigma `$, let $`T`$ be the type of $`\sigma `$, and let $`O`$ be the right $`\overline{\mathrm{\Gamma }}_T`$-orbit in $`^{n1}`$ corresponding to $`\sigma `$ as in Proposition 3.5. Then $`H^0(\mathrm{\Gamma }_\sigma ;)`$ is $``$ if $`O`$ is orientable, and is 0 if $`O`$ is non-orientable.
###### Proof.
When $`O`$ is orientable, this follows from Proposition 3.14, merely because $`\mathrm{\Gamma }_\sigma \mathrm{\Gamma }`$. Now assume $`O`$ is non-orientable. Let $`\sigma =\gamma _0\sigma _T`$, with $`𝐚=𝔟(\gamma _0)`$. As we have said above, there is some $`\gamma _1\overline{\mathrm{\Gamma }}_T\overline{\mathrm{\Gamma }}_T^+`$ with $`𝐚\gamma _1=𝐚`$. By (2.2), $`\mathrm{\Gamma }_\sigma =(\gamma _0\overline{\mathrm{\Gamma }}_T\gamma _0^1)\mathrm{\Gamma }`$. The element $`\gamma _0\gamma _1\gamma _0^1`$ is in $`\mathrm{\Gamma }`$ by Lemma 2, so it is in $`\mathrm{\Gamma }_\sigma `$; clearly it carries $`\sigma `$ to itself while reversing orientation. Hence $`\mathrm{\Gamma }_\sigma `$ acts non-trivially on $`_\sigma `$, meaning $`H^0(\mathrm{\Gamma }_\sigma ;_\sigma )=0`$. ∎
###### Remark 3.17.
The proposition shows that the $`\mathrm{\Gamma }`$-orbits of cells coming from non-orientable $`O`$ contribute nothing to our spectral sequence. We ignore these objects for the rest of the computation, tacitly assuming that all $`O`$’s mentioned from now on are orientable.
To summarize:
###### Proposition 3.18.
The $`E_1^{i,0}`$ term of the equivariant cohomology spectral sequence for $`H_\mathrm{\Gamma }^i(W;)`$ is a direct sum $`_o`$, where $`o`$ runs through a set of $`i`$-cells in one-to-one correspondence with the orientable right $`\overline{\mathrm{\Gamma }}_T`$-orbits in $`^{n1}`$, for all types $`T`$ of cells of dimension $`i`$.
### 3.19.
In §3.21, we will describe the boundary maps $`d_1`$ of the spectral sequence. These are the only differentials we need consider, since the sequence collapses at $`E_2`$. In this subsection, we give some details concerning how the cells meet at their boundaries.
As usual, a *facet* of a cell $`\sigma `$ is any face of $`\sigma `$ of codimension one. Let $`_\sigma `$ be the set of facets of $`\sigma `$ in $`W`$.
We will need to understand how $`_\sigma `$ breaks up into orbits under the action of $`\mathrm{\Gamma }_\sigma `$ (and to choose a set $`_\sigma ^{}`$ of representatives of these orbits in $`_\sigma `$). We do this in Proposition 3.20 below. We will start by determining $`_{\sigma _T}`$ for the standard cells $`\sigma _T`$ in a form suited to our computation. We will then determine $`_\sigma `$ for any $`\sigma `$.
We make two conventions. (i) If types $`T`$ and $`T^{}`$ occur in the same discussion, it is assumed that a cell of type $`T`$ has at least some cells of type $`T^{}`$ as facets. (ii) If $`\sigma =x\sigma _T`$ for some $`x\overline{\mathrm{\Gamma }}`$, we identify $`\sigma `$ with the coset $`x\overline{\mathrm{\Gamma }}_T`$ as in Proposition 3.3. In the expressions of the form
$$\underset{T^{}}{}(),$$
the ($``$) will be a finite union of cosets, say $`x_1\overline{\mathrm{\Gamma }}_T^{}\mathrm{}x_k\overline{\mathrm{\Gamma }}_T^{}`$ for $`x_1,\mathrm{},x_k\overline{\mathrm{\Gamma }}`$. It is understood that ($``$) corresponds to the set of cells $`x_1\sigma _T^{},\mathrm{},x_k\sigma _T^{}`$, as in Proposition 3.3. Even when ($``$) is a more complicated object, like a double coset, its meaning is that one should decompose it into single cosets by choosing appropriate representatives (which will be the $`x_i`$).
The boundary of $`\sigma _T`$ is a union of (the closures of) various cells of type $`T^{}`$:
(8)
$$_{\sigma _T}=\underset{T^{}}{}\beta _{T^{},1}\overline{\mathrm{\Gamma }}_T^{}\mathrm{}\beta _{T^{},k}\overline{\mathrm{\Gamma }}_T^{}$$
for some $`\beta _{T^{},i}\overline{\mathrm{\Gamma }}`$. However, (8) is invariant under the left action of $`\sigma _T`$’s stabilizer $`\overline{\mathrm{\Gamma }}_T`$. Hence there must be finitely many $`\alpha _{(T,T^{},\iota )}\overline{\mathrm{\Gamma }}`$ such that
(9)
$$_{\sigma _T}=\underset{T^{}}{}\overline{\mathrm{\Gamma }}_T\alpha _{(T,T^{},1)}\overline{\mathrm{\Gamma }}_T^{}\mathrm{}\overline{\mathrm{\Gamma }}_T\alpha _{(T,T^{},k)}\overline{\mathrm{\Gamma }}_T^{}.$$
We have computed the $`\alpha _{(T,T^{},\iota )}`$ by hand for $`n4`$. In our cases of interest ($`n=4`$, $`dim\sigma 4`$), one finds there is only one $`\iota `$—that is, the right-hand side of (9) is actually just one double coset. In fact, we find that we may take $`\alpha _{(T,T^{},\iota )}`$ to be the identity except in one case, where $`(T,T^{})=(`$5b, 4b$`)`$ in the notation of and
$$\alpha _{\text{5b},\text{4b},1}=\left(\begin{array}{cccc}0& 0& 1& 0\\ 0& 0& 0& 1\\ 1& 0& 1& 0\\ 0& 1& 1& 0\end{array}\right).$$
From now on, we will write $`\alpha `$ for $`\alpha _{(T,T^{},1)}`$, the dependence on $`T`$ and $`T^{}`$ being understood.
To find an expression for $`_\sigma `$ for $`\sigma =\gamma _0\sigma _T`$, we multiply (9) by $`\gamma _0`$ to obtain
(10)
$$_\sigma =\underset{T^{}}{}\gamma _0\overline{\mathrm{\Gamma }}_T\alpha \overline{\mathrm{\Gamma }}_T^{}.$$
To exhibit this as a union of cosets of the form $`x_i\overline{\mathrm{\Gamma }}_T^{}`$, we rewrite it as
(11) $`_\sigma `$ $`={\displaystyle \underset{T^{}}{}}\gamma _0\alpha (\alpha ^1\overline{\mathrm{\Gamma }}_T\alpha )\overline{\mathrm{\Gamma }}_T^{}`$
(12) $`={\displaystyle \underset{T^{}}{}}\gamma _0\alpha (\alpha ^1\overline{\mathrm{\Gamma }}_T\alpha /(\alpha ^1\overline{\mathrm{\Gamma }}_T\alpha \overline{\mathrm{\Gamma }}_T^{}))\overline{\mathrm{\Gamma }}_T^{},`$
The last formula (for each $`T^{}`$) is a disjoint union of single cosets, in one-to-one correspondence with a set of representatives of the cosets in $`\alpha ^1\overline{\mathrm{\Gamma }}_T\alpha /(\alpha ^1\overline{\mathrm{\Gamma }}_T\alpha \overline{\mathrm{\Gamma }}_T^{})`$. It is a matter of formal manipulation to get an expression for $`\mathrm{\Gamma }_\sigma \backslash _\sigma `$:
(13) $`\mathrm{\Gamma }_\sigma \backslash _\sigma `$ $`={\displaystyle \underset{T^{}}{}}(\gamma _0\overline{\mathrm{\Gamma }}_T\gamma _0^1\mathrm{\Gamma })\backslash \gamma _0\alpha (\alpha ^1\overline{\mathrm{\Gamma }}_T\alpha /(\alpha ^1\overline{\mathrm{\Gamma }}_T\alpha \overline{\mathrm{\Gamma }}_T^{}))\overline{\mathrm{\Gamma }}_T^{}`$
(14) $`={\displaystyle \underset{T^{}}{}}\gamma _0(\overline{\mathrm{\Gamma }}_T\gamma _0^1\mathrm{\Gamma }\gamma _0)\backslash \alpha (\alpha ^1\overline{\mathrm{\Gamma }}_T\alpha /(\alpha ^1\overline{\mathrm{\Gamma }}_T\alpha \overline{\mathrm{\Gamma }}_T^{}))\overline{\mathrm{\Gamma }}_T^{}`$
(15) $`={\displaystyle \underset{T^{}}{}}\underset{B}{\underset{}{\gamma _0\alpha [\underset{A}{\underset{}{\alpha ^1(\overline{\mathrm{\Gamma }}_T\gamma _0^1\mathrm{\Gamma }\gamma _0)\alpha }}\backslash \alpha ^1\overline{\mathrm{\Gamma }}_T\alpha }}/\underset{C}{\underset{}{(\alpha ^1\overline{\mathrm{\Gamma }}_T\alpha \overline{\mathrm{\Gamma }}_T^{})}}]\overline{\mathrm{\Gamma }}_T^{}`$
where for each $`T^{}`$ the right-hand side of (15) is again expressed as a disjoint union of $`\overline{\mathrm{\Gamma }}_T^{}`$-cosets, in one-to-one correspondence with a set of representatives of the double coset expression in square brackets.
We must interpret (11)–(12) in terms of the $`\overline{\mathrm{\Gamma }}_T^{}`$-orbits in $`^{n1}`$. As usual, let $`O`$ be the right $`\overline{\mathrm{\Gamma }}_T`$-orbit corresponding to $`\sigma `$, with $`𝐚=𝔟(\gamma _0)`$. Let $`𝐜=𝔟(\gamma _0\alpha )`$, so that $`𝐜=𝐚\alpha `$. Now the quantity $`A`$ in (15) is exactly the subgroup of $`\alpha ^1\overline{\mathrm{\Gamma }}_T\alpha `$ that preserves $`𝐜`$. Hence any set of coset representatives for $`B`$ in (15) is a set of matrices whose bottom rows, in $`^{n1}`$, are the members of the right $`(\alpha ^1\overline{\mathrm{\Gamma }}_T\alpha )`$-orbit of $`𝐜`$. Equivalently, any set of representatives for $`B`$ is a set of matrices whose bottom rows are exactly the members of $`O\alpha `$. The group denoted $`C`$ in (15) acts on this orbit $`O\alpha `$, decomposing it into suborbits; for given $`T^{}`$, the disjoint $`\overline{\mathrm{\Gamma }}_T^{}`$-cosets in (11)–(12) are in one-to-one correspondence with these suborbits. We summarize this result as follows:
###### Proposition 3.20.
Let $`T`$, $`T^{}`$ and $`\alpha `$ be as introduced in this subsection. Let $`\sigma `$ be a cell of type $`T`$, represented by the $`\overline{\mathrm{\Gamma }}_T`$-orbit $`O`$ in the manner of §3.8. Decompose the orbit $`O\alpha `$ into its suborbits $`O_1,\mathrm{},O_k`$ under the group $`C=\alpha ^1\overline{\mathrm{\Gamma }}_T\alpha \overline{\mathrm{\Gamma }}_T^{}`$. Let $`𝐚_jO_j`$. Let $`\gamma _j\overline{\mathrm{\Gamma }}`$ be chosen so that $`𝔟(\gamma _j)=𝐚_j`$. (We may, in fact, take $`\gamma _j`$ to be of the form $`\gamma _0\widehat{\gamma }_j\alpha `$ for some $`\widehat{\gamma }_j\mathrm{\Gamma }_T`$.) Then the union over all $`T^{}`$ of the cells
$$\{\gamma _1\sigma _T^{},\mathrm{},\gamma _k\sigma _T^{}\}$$
is a set $`_\sigma ^{}`$ of representatives for $`\mathrm{\Gamma }_\sigma \backslash _\sigma `$.
For all pairs $`T,T^{}`$, we store the intersection $`C=\alpha ^1\overline{\mathrm{\Gamma }}_T\alpha \overline{\mathrm{\Gamma }}_T^{}`$ in our program.
### 3.21.
We can now determine the boundary maps $`d_1`$ in the spectral sequence. Recall that $`W_{(i)}`$ is a fixed set of representatives of the $`\mathrm{\Gamma }`$-orbits of the cells in $`W`$ of dimension $`i`$. Superseding the use of $`o`$ in §3.15, we let $`\sigma `$ run through $`W_{(i+1)}`$, and let $`\tau `$ run through the set $`_\sigma ^{}`$ of representatives of the facets of $`\sigma `$. To compute $`H_\mathrm{\Gamma }^i(W;)`$ for subgroups of $`\mathrm{SL}_4()`$ for $`i=5,6`$, we must work out the map $`d_1^{i,0}`$ for $`i=4`$ and 5.
We follow \[13, VII.8\], taking the dual to turn homology into cohomology. The map
(16)
$$d_1^{i,0}:\underset{\tau _0W_{(i)}}{}H^0(\mathrm{\Gamma }_{\tau _0};_{\tau _0})\underset{\sigma W_{(i+1)}}{}H^0(\mathrm{\Gamma }_\sigma ;_\sigma )$$
is a sum of terms, one for each pair $`(\sigma ,\tau )`$; here $`\tau _0`$ is the fixed representative in $`W_{(i)}`$ that is $`\mathrm{\Gamma }`$-equivalent to $`\tau `$. From now on, we focus on a single pair $`\sigma `$, $`\tau `$. Call their types $`T`$ and $`T^{}`$, respectively. Let $`O`$ be the right $`\overline{\mathrm{\Gamma }}_T`$-orbit in $`^{n1}`$ associated to $`\sigma `$ by Proposition 3.5. If $`\sigma =\gamma _0\sigma _T`$, let $`𝐚=𝔟(\gamma _0)O`$. Let $`\mathrm{\Gamma }_{\sigma \tau }=\mathrm{\Gamma }_\sigma \mathrm{\Gamma }_\tau `$.
As in \[13, p. 176\], the term for $`(\sigma ,\tau )`$ in $`d_1^{i,0}`$ is the composition
(17)
$$\text{}.$$
We now give the definition of these maps and show how to compute them. Note that all the coefficient modules are copies of $``$ on which the groups act trivially.
The map $`t_{\sigma \tau }`$ is the transfer map $``$ given by multiplication by the scalar $`[\mathrm{\Gamma }_\sigma :\mathrm{\Gamma }_{\sigma \tau }]`$. We evaluate this scalar as follows. First of all, $`\mathrm{\Gamma }_\sigma =(\gamma _0\overline{\mathrm{\Gamma }}_T\gamma _0^1)\mathrm{\Gamma }`$, which has the same cardinality as $`\overline{\mathrm{\Gamma }}_T(\gamma _0^1\mathrm{\Gamma }\gamma _0)`$. The latter is the subgroup of $`\overline{\mathrm{\Gamma }}_T`$ that fixes $`𝐚`$. Thus
(18)
$$\mathrm{\#}(\mathrm{\Gamma }_\sigma )=\frac{\mathrm{\#}(\overline{\mathrm{\Gamma }}_T)}{\mathrm{\#}(\overline{\mathrm{\Gamma }}_T\text{-orbit of }𝐚)}=\frac{\mathrm{\#}(\overline{\mathrm{\Gamma }}_T)}{\mathrm{\#}(O)}.$$
The numerator is known because we stored $`\overline{\mathrm{\Gamma }}_T`$. The denominator is easily recovered from the computer’s lists of orbits.
We now evaluate the order of $`\mathrm{\Gamma }_{\sigma \tau }`$. As before, $`\mathrm{\Gamma }_\sigma =\gamma _0\overline{\mathrm{\Gamma }}_T\gamma _0^1\mathrm{\Gamma }`$. By Proposition 3.20, $`\tau =\gamma _0\widehat{\gamma }\alpha \sigma _T^{}`$ for some $`\widehat{\gamma }\overline{\mathrm{\Gamma }}_T`$. Hence
$$\mathrm{\Gamma }_\tau =((\gamma _0\widehat{\gamma }\alpha )\overline{\mathrm{\Gamma }}_T^{}(\alpha ^1\widehat{\gamma }^1\gamma _0^1))\mathrm{\Gamma }.$$
Writing out $`\mathrm{\Gamma }_{\sigma \tau }=\mathrm{\Gamma }_\sigma \mathrm{\Gamma }_\tau `$ and conjugating by $`(\gamma _0\widehat{\gamma }\alpha )`$, we find that $`\mathrm{\Gamma }_{\sigma \tau }`$ has the same cardinality as
(19)
$$\underset{D}{\underset{}{(\alpha ^1\overline{\mathrm{\Gamma }}_T\alpha \overline{\mathrm{\Gamma }}_T^{})}}(\alpha ^1\widehat{\gamma }^1\gamma _0^1\mathrm{\Gamma }\gamma _0\widehat{\gamma }\alpha )$$
The group $`D`$ is the same as the group $`C`$ from (15), and is one of the groups we’ve stored. And the whole group in (19) is exactly the subgroup of $`D`$ that fixes the point $`𝐚_j`$ of Proposition 3.20. So
(20)
$$\mathrm{\#}(\mathrm{\Gamma }_{\sigma \tau })=\frac{\mathrm{\#}(\alpha ^1\overline{\mathrm{\Gamma }}_T\alpha \overline{\mathrm{\Gamma }}_T^{})}{\mathrm{\#}(O_j)}.$$
Again, the numerator is known from what’s stored, and the denominator is easy to evaluate.
We have
(21)
$$t_{\sigma \tau }=\frac{\text{(}\text{18}\text{)}}{\text{(}\text{20}\text{)}}.$$
As a consistency check, the program signals an error if the computed value of (21) is not an integer.
Brown’s $`u_{\sigma \tau }:`$ is the composition
$$H^0(\mathrm{\Gamma }_\tau ;_\tau )H^0(\mathrm{\Gamma }_{\sigma \tau };_\tau )H^0(\mathrm{\Gamma }_{\sigma \tau };_\sigma ).$$
The first arrow is induced from the inclusion $`\mathrm{\Gamma }_{\sigma \tau }\mathrm{\Gamma }_\tau `$, and is easily seen to be the identity. The second arrow is induced by the $`\mathrm{\Gamma }_{\sigma \tau }`$-map $`_{\sigma \tau }:_\sigma _\tau `$, namely the $`(\sigma ,\tau )`$-component of the cellular boundary operator on $`W`$. Let $`[\sigma :\tau ]`$ be $`\pm 1`$ depending on whether the orientation on $`\sigma `$ from Definition 3.10 does or does not induce the orientation on the facet $`\tau `$ from Definition 3.10. Then $`_{\sigma \tau }`$ is the map $``$ given by the scalar $`[\sigma :\tau ]`$. Thus $`u_{\sigma \tau }=[\sigma :\tau ]`$.
To evaluate $`[\sigma :\tau ]`$, we reduce the problem to evaluating the boundary operator on a small list of pairs of cells. Write $`\tau =\gamma _0\widehat{\gamma }\alpha \sigma _T^{}`$. Then $`\sigma =\gamma _0\widehat{\gamma }\sigma _T`$, since $`\widehat{\gamma }\overline{\mathrm{\Gamma }}_T`$. If $`𝐱^{n1}`$ is part of any oriented $`\overline{\mathrm{\Gamma }}_T`$-orbit (for any $`T`$), write $`\mathrm{sgn}_T(𝐱)`$ for the orientation number of $`𝐱`$ with respect to $`\overline{\mathrm{\Gamma }}_T`$. By Definition 3.10, $`\sigma `$ and $`\tau `$—together with their standard orientations—are given as follows:
(22) $`\sigma `$ $`=(\gamma _0\widehat{\gamma })_{}(\sigma _T)\mathrm{sgn}_T(𝔟(\gamma _0\widehat{\gamma }))`$
(23) $`\tau `$ $`=(\gamma _0\widehat{\gamma }\alpha )_{}(\sigma _T^{})\mathrm{sgn}_T^{}(𝔟(\gamma _0\widehat{\gamma }\alpha )).`$
Because $`(\mathrm{})_{}`$ is functorial,
$$\tau =(\gamma _0\widehat{\gamma })_{}(\alpha )_{}(\sigma _T^{})\mathrm{sgn}_T^{}(𝔟(\gamma _0\widehat{\gamma }\alpha )).$$
Since $`(\mathrm{})_{}`$ preserves the $`[\text{ }:\text{ }]`$ relation, we may cancel out $`(\gamma _0\widehat{\gamma })_{}`$’s, obtaining
(24)
$$[\sigma :\tau ]=\underset{E}{\underset{}{\mathrm{sgn}_T(𝔟(\gamma _0\widehat{\gamma }))}}\underset{F}{\underset{}{\mathrm{sgn}_T^{}(𝔟(\gamma _0\widehat{\gamma }\alpha ))}}[\sigma _T:(\alpha )_{}(\sigma _T^{})].$$
We evaluate each factor in (24) in turn. In $`F`$, $`𝔟(\gamma _0\widehat{\gamma }\alpha )`$ is the point $`𝐚_j`$ of Proposition 3.20, and $`\mathrm{sgn}_T^{}(𝐚_j)`$ is simply the orientation number of $`𝐚_j`$ in its $`\overline{\mathrm{\Gamma }}_T^{}`$-orbit. Similarly, $`E`$ is the orientation number of $`𝔟(\gamma _0\widehat{\gamma })`$, which is a $`\overline{\mathrm{\Gamma }}_T`$-translate of $`𝐚`$, in its $`\overline{\mathrm{\Gamma }}_T`$-orbit. In practice, though, we do not know $`\widehat{\gamma }`$ or $`𝔟(\gamma _0\widehat{\gamma })`$ explicitly, so it is easier to evaluate $`E`$ by another method. Note that $`𝔟(\gamma _0\widehat{\gamma })\alpha =𝔟(\gamma _0\widehat{\gamma }\alpha )=𝐚_j`$, an element of the $`(\alpha ^1\overline{\mathrm{\Gamma }}_T\alpha )`$-orbit $`O\alpha `$ of Proposition 3.20. It is easy to compute all the points in $`O\alpha `$, using the group $`\alpha ^1\overline{\mathrm{\Gamma }}_T\alpha `$ (which was stored). For some $`𝐱`$, we will have the equation $`𝐚_j=𝐱\alpha `$. Then $`E`$ will be the orientation number for this $`𝐱`$.
The quantities $`[\sigma _T:(\alpha )_{}(\sigma _T^{})]`$ are evaluated by hand, for all $`T,T^{}`$ occurring in (9). Several issues arise. First, we must find $`[\sigma _T:\sigma _T^{}]`$ whenever the standard cell $`\sigma _T^{}`$ is a facet of $`\sigma _T`$. No matter how we choose the orientations at the start, it is in general impossible to arrange our choices so that all the relative orientations are positive. As a general illustration, if a 0-cell, two 1-cells, and a 2-cell meet locally in a picture like the first quadrant of $`^2`$, and if both 1-cells are oriented to point away from the origin, then the pair (2-cell, $`x`$-axis) must have relative orientation opposite to that of the pair (2-cell, $`y`$-axis), no matter how we orient the 2-cell. In the $`\mathrm{SL}_4()`$ case, one can draw a schematic picture of how the standard cells meet and can read off all the relative orientations $`[\sigma _T:\sigma _T^{}]`$.
Second, we must compute $`(\alpha )_{}(\sigma _T^{})`$ when $`\alpha `$ isn’t the identity, knowing the orientation on $`\sigma _T^{}`$. This involves the same techniques as in Remark 3.9.
Third, we must find some facet $`\upsilon `$ of $`\sigma _T`$ whose orientation is known, and must compare $`[\sigma _T:\upsilon ]`$ to $`[\sigma _T:(\alpha )_{}(\sigma _T^{})]`$. In practice, if $`dim\sigma _T=k`$, this means finding a chain of $`(k1)`$-cells between $`\upsilon `$ and $`(\alpha )_{}(\sigma _T^{})`$ such that consecutive members of the chain meet in faces of dimension $`k2`$, and comparing the orientations of $`\upsilon `$ and $`(\alpha )_{}(\sigma _T^{})`$ across the $`(k2)`$-faces.
Finally, we must compute $`v_\tau `$, but this is easy. It is induced by the conjugation action of the element of $`\mathrm{\Gamma }`$ that carries $`\tau `$ to $`\tau _0`$. But by Proposition 3.14, any element of $`\mathrm{\Gamma }`$ preserves the orientations on the cells. Hence one finds $`v_\tau =+1`$.
## 4. Hecke operators
### 4.1.
We identify the cochain complex in (16) with the complex of cellular cochains on $`W`$ by identifying an $`i`$-cell $`\sigma `$ with a generator of $`H^0(\mathrm{\Gamma }_\sigma ;_\sigma )`$, taking care to make the signs match. Formulas like $`n(\sigma )\sigma `$ will denote the corresponding cocycles in either complex, and will be referred to as *$`W`$-cocycles*
Let $`\beta H^5(\mathrm{\Gamma };)`$ be a class, and let $`u=n(\sigma )\sigma `$ be a representative for $`\beta `$ in terms of the previous paragraph. Let $`T(l,k)`$ be a Hecke operator. To compute the action of $`T(l,k)`$ on $`u`$, we do the following:
1. Convert $`u`$ to a reduced $`1`$-sharbly cycle $`\xi `$, and then compute $`T(l,k)(\xi )`$ using (2) in §2.7.
2. Use the algorithm from to write $`T(l,k)(\xi )`$ as a sum of reduced $`1`$-sharbly cycles.
3. Convert these reduced $`1`$-sharbly cycles to $`W`$-cocycles.
Step 2 is described in detail in , and we refer the reader to that article. In this section, we focus on steps 1 and 3 in the context of §3.
We begin with a definition from :
###### Definition 4.2.
\[18, Definition 5.3\] Let $`𝐮`$ be a basis element of the $`k`$-sharblies $`S_k`$. Then a *lift* for $`𝐮`$ is an $`n\times (n+k)`$ integral matrix $`M`$ with primitive columns such that $`[M_1,\mathrm{},M_{n+k}]=𝐮`$, where $`M_i`$ is the $`i`$th column of $`M`$.
Modulo the action of $`\mathrm{GL}_4()`$, there is only one orbit of reduced basis $`0`$-sharblies and only three orbits of reduced basis $`1`$-sharblies. The identity matrix serves as a lift for a member of the first orbit, and lifts representing elements of the latter three orbits are
$$\left(\begin{array}{ccccc}1& 0& 0& 0& 1\\ 0& 1& 0& 0& 1\\ 0& 0& 1& 0& 1\\ 0& 0& 0& 1& 1\end{array}\right),\left(\begin{array}{ccccc}1& 0& 0& 0& 1\\ 0& 1& 0& 0& 1\\ 0& 0& 1& 0& 1\\ 0& 0& 0& 1& 0\end{array}\right),\text{and}\left(\begin{array}{ccccc}1& 0& 0& 0& 1\\ 0& 1& 0& 0& 1\\ 0& 0& 1& 0& 0\\ 0& 0& 0& 1& 0\end{array}\right).$$
We call these the *standard* $`0`$\- and $`1`$-sharblies. The sets of primitive vectors indexing the standard $`6`$ and $`5`$-cells in $`W`$ coincide with the sets of column vectors of these matrices. By abuse of language we will speak of the “standard sharbly of type $`T`$,” and will use the notation $`𝐮_T`$.
### 4.3.
Given a sharbly cycle $`\xi `$, we denote by $`\mathrm{supp}\xi `$ the support of $`\xi `$. Suppose that $`\mathrm{\Gamma }`$ is torsion-free. Then according to , a $`1`$-sharbly cycle $`\xi `$ mod $`\mathrm{\Gamma }`$ with coefficients in a ring $`R`$ can be encoded by a collection of $`4`$-tuples $`(𝐮,n(𝐮),\{𝐯\},\{L(𝐯)\})`$, where
1. $`𝐮\mathrm{supp}\xi `$,
2. $`n(𝐮)R`$,
3. $`\{𝐯\}=\mathrm{supp}𝐮`$, and
4. $`\{L(𝐯)\}`$ is a $`\mathrm{\Gamma }`$-equivariant set of lifts for $`\{𝐯\}`$.
The $`\mathrm{\Gamma }`$-equivariance condition in 4 is the following. Suppose that for $`𝐮,𝐮^{}\mathrm{supp}\xi `$ there exist $`𝐯\mathrm{supp}(𝐮)`$ and $`𝐯^{}\mathrm{supp}(𝐮^{})`$ such that $`𝐯=\gamma 𝐯^{}`$ for some $`\gamma \mathrm{\Gamma }`$. Then we require $`L(𝐯)=\gamma L(𝐯^{})`$.
In the case under study, $`\mathrm{\Gamma }`$ is *not* torsion-free, and the above data needs to be modified. Suppose that a $`0`$-sharbly $`𝐯\mathrm{supp}𝐮`$ has a nontrivial stabilizer $`\mathrm{\Gamma }(𝐯)\mathrm{\Gamma }`$, and let $`m`$ be any lift of $`𝐯`$. Then in the cycle $`\xi `$ we replace $`n(𝐮)𝐮`$ with
$$\underset{\gamma \mathrm{\Gamma }(𝐯)}{}\frac{n(𝐮)}{\mathrm{\#}\mathrm{\Gamma }(𝐯)}𝐮_\gamma ,$$
where $`𝐮_\gamma `$ has the same data as $`𝐮`$, except that we give $`𝐯`$ the lift $`\gamma m`$. (Note that this is possible in our case since the coefficient ring $`R=`$ is divisible.)
### 4.4.
Now we describe how to construct the data in §4.3 to produce a $`1`$-sharbly chain $`\xi `$ corresponding to the $`W`$-cocycle $`u`$. There are two steps.
First, choose $`\sigma `$ such that $`n(\sigma )0`$ in $`u`$. According to Proposition 3.3, the $`5`$-cell $`\sigma \mathrm{\Gamma }\backslash W`$ is encoded as a coset $`[\gamma \mathrm{\Gamma }_T]`$, where $`\gamma \mathrm{\Gamma }`$ and $`\mathrm{\Gamma }_T`$ is the stabilizer of the standard cell of type $`T`$. Moreover, the coset $`[\gamma \mathrm{\Gamma }_T]`$ is encoded as the orbit $`O`$, which can be regarded a set of triples $`\{(𝐚,\pm 1,\gamma _𝐚)\}`$, where $`𝐚^3(/N)`$, $`\gamma _𝐚\mathrm{\Gamma }_T`$, and $`\pm 1`$ is the orientation number (§3.8). From $`O`$ we arbitrarily choose a triple with orientation number $`1`$, and then using Hermite normal form construct a matrix $`\gamma \mathrm{SL}_4()`$ with bottom row equal to $`𝐚`$. Then the contribution of $`\sigma `$ to the $`1`$-sharbly chain is
$$n(\sigma )\gamma 𝐮_T.$$
We do this for all of $`\mathrm{supp}u`$ and sum to produce $`\xi `$. For each $`𝐮\mathrm{supp}\xi `$, we write $`𝐮(\gamma )`$ if we want to indicate the element $`\gamma `$ used in the construction of $`𝐮`$.
### 4.5.
At this stage, we have a $`1`$-sharbly chain, and we need to choose lifts to reflect the cycle structure of $`\xi `$. This we do as follows. In the spirit of §3.19, for each type $`T`$ we choose a set of matrices $`\mathrm{\Omega }_T\mathrm{SL}_4()`$ such that $`:S_1S_0`$ can be written as
$$:𝐮_T\underset{\omega \mathrm{\Omega }_T}{}\omega 𝐯,$$
where $`𝐯`$ is the standard basis $`0`$-sharbly. Note the absence of signs in this map—the signs in the boundary map in Definition 2.6 have been encoded in the $`\omega `$’s, which may nontrivially permute the column vectors of $`𝐯`$.
Form the $`0`$-sharbly chain
(25)
$$\underset{𝐮(\gamma ),\omega \mathrm{\Omega }_T}{}\gamma \omega 𝐯,$$
where $`𝐯`$ is the standard $`0`$-sharbly, we sum over all $`𝐮(\gamma )\mathrm{supp}\xi `$, and $`T`$ is the type of $`𝐮(\gamma )`$. (In (25) we have abbreviated $`\alpha _{(T,T^{},i)}`$ to $`\alpha _i`$, since $`T`$ is determined by $`𝐮(\gamma )`$, and there is only one type of reduced $`0`$-sharbly mod $`\mathrm{SL}_4()`$.) Note that this sum is in $`S_{}`$, not $`(S_{})_\mathrm{\Gamma }`$; the only relations we apply are those in the sharbly complex.
### 4.6.
After summing, we find that some $`0`$-sharblies cancel, and some remain. For those that canceled, we can choose any lifts we like, as long as we choose the same lifts for all terms that cancel each other.
The remaining $`0`$-sharblies form a chain $`\eta `$ that vanishes in $`(S_0)_\mathrm{\Gamma }`$, and we must choose nontrivial lifts for them. To do this, first arbitrarily choose lifts for each $`0`$-sharbly in $`\mathrm{supp}\eta `$. The data we computed in §3.2 allows us to easily compute the distinct orbits of $`\mathrm{\Gamma }`$ in $`\mathrm{supp}\eta `$. We do this and order each orbit.
Suppose $`𝐯_0`$ is the first $`0`$-sharbly in one of these orbits, and that it corresponds to the triple $`(𝐚_0,\pm 1,\gamma _{𝐚_0})`$. Let $`𝐯_0(\gamma )`$ be its lift. Then if $`𝐯`$ is any other $`0`$-sharbly in $`𝐯_0`$’s orbit, corresponding to the triple $`(𝐚,\pm 1,\gamma _𝐚)`$, we replace the lift $`𝐯(\gamma )`$ by
$$𝐯(\gamma )𝐯(\gamma )\gamma _𝐚^1\gamma _{𝐚_0}$$
After these lifts are constructed, the cycle $`\xi `$ is ready for input in the Hecke operator program.
### 4.7.
Upon completion, the Hecke operator program returns a reduced $`1`$-sharbly cycle, which we must convert to a $`W`$-cocycle. So let $`\xi =n(𝐮)𝐮`$ be a reduced $`1`$-sharbly cycle, and let $`𝐮\mathrm{supp}\xi `$. First we determine which of the three types $`T`$ of standard reduced $`1`$-sharblies $`𝐮`$ has. Then we must find a matrix $`\gamma (𝐮)\mathrm{SL}_4()`$ such that
$$\gamma (𝐮)𝐮_T=𝐮,$$
where $`𝐮_T`$ is the standard $`1`$-sharbly of type $`T`$. This is straightforward, although one must be careful to incorporate the orientation number of $`\gamma (𝐮)`$.
In practice, the main step is the following. Let $`𝐮`$ be a reduced $`1`$-sharbly basis element with lift $`M`$. We choose a nonsingular $`4\times 4`$ minor $`m`$ of $`M`$ and construct $`m^1`$. Then $`m^1𝐮`$ will be a $`1`$-sharbly with lift $`m^1M`$, and will be standard except possibly for one column vector. By multiplying $`m^1M`$ on the left by elements of the stabilizer of the standard $`0`$-sharbly, we can eventually produce the standard $`1`$-sharbly with same type as $`𝐮`$. This allows us to construct $`\gamma (𝐮)`$.
## 5. Numerical results
### 5.1.
In this section we present numerical data from our experiments. As mentioned in §1, to avoid floating point problems with $``$-coefficients we usually work with $`𝔽=𝔽_{31991}`$, the finite field with $`31991`$ elements, and in some cases work with $``$ or $``$. The computations were carried out on a variety of Unix machines at Columbia and Oklahoma State. The code for the cohomology of $`W`$3) was written in Common Lisp. The Hecke operator code (§4) was written in C++, and used the LiDIA library . Perl scripts patched together the outputs of the various programs and produced the tables in §5.5.
### 5.2.
We first describe how we performed linear algebra on the large sparse matrices that arise in our computations. Fix the level $`N`$. We use the notation of §3.21, working over $``$ at first, and letting $`V_{(i)}`$ be the domain of the map $`d_1^{i,0}`$ in (16). To compute $`H^5(\mathrm{\Gamma };)`$, we must find the kernel of $`d_1^{5,0}`$ modulo the image of $`d_1^{4,0}`$. We prefer to find the kernel of a single matrix. Regard the $`d_1^{i,0}`$ as matrices acting on the left on column vectors, and let $`𝔇_{}`$ be the matrix where $`d_1^{5,0}`$ is stacked on top of the transpose of $`d_1^{4,0}`$:
$$𝔇_{}=\left(\begin{array}{c}d_1^{5,0}\\ (d_1^{4,0})^{tr}\end{array}\right)$$
This matrix defines a map $`V_{(5)}V_{(6)}V_{(4)}`$, where we use the standard inner product to identify the transpose of $`d_1^{4,0}`$ with its adjoint. The kernel of this map is the space of *harmonic 5-cocycles*; it is isomorphic to $`H^5(\mathrm{\Gamma };)`$. Let $`s=dimV_{(5)}`$, the number of columns of $`𝔇_{}`$.
The matrix $`𝔇_{}`$ has coefficients in $``$. For any ring $`R`$, set $`𝔇_R=𝔇_{}R`$, and set $`𝔇=𝔇_𝔽`$. In the tables in §5.4, the value of “rank” we report in the rows labeled $`R`$ is a number almost certainly equal to $`dim\mathrm{ker}𝔇_R`$, whose computation is explained below. This number is also almost certainly equal to $`dimH^5(\mathrm{\Gamma };)`$.
### 5.3.
When $`𝔇`$ is very large, we could not have found its kernel without a sparse version of the Lanczos algorithm. This algorithm is usually used with real or complex matrices, particularly for eigenvalue problems. Following ideas in , we translated it into the mod-$`p`$ setting. Let $`𝔈=𝔇^{tr}𝔇`$, a symmetrized version of $`𝔇`$. We choose a random non-zero seed vector $`𝐯`$ with coefficients in $`𝔽`$ and consider the sequence $`𝐯,𝔈𝐯,𝔈^2𝐯,\mathrm{}`$. The Lanczos algorithm shows us how to compute not this sequence, but the sequence $`𝐯=𝐪_0,𝐪_1,\mathrm{}`$ resulting from it by the Gram-Schmidt orthogonalization process. We perform the Gram-Schmidt process mod $`p`$ in the naive way, using $`x_jy_j`$ for the inner product. This means we don’t have the usual guarantee that $`x_j^2`$ will be non-zero when $`(x_1,x_2,\mathrm{})(0,0,\mathrm{})`$. If the inner product is ever 0 for non-trivial $`(x_j)`$, we simply abort and choose another random seed $`𝐯`$; even for large $`𝔇`$, these aborts happen less than half the time.
The strength of the algorithm is that the RAM only has to hold the sparse matrix $`𝔇`$ and a few vectors of storage. It does not have to hold $`𝔈`$, which is dense in general. The $`𝐪`$’s form a dense matrix, but they may be stored on the disk, not in RAM. (In our implementation, $`𝔇^{tr}`$ was stored in RAM along with $`𝔇`$.)
Let $`k`$ be the largest value for which the set $`\{𝐪_0,\mathrm{},𝐪_{sk}\}`$ is linearly dependent. Reading the $`𝐪`$’s back in from disk, the algorithm allows us to backsolve for a non-zero vector $`𝐲\mathrm{ker}𝔈`$. What we want is an element of $`\mathrm{ker}𝔇`$. A priori, we only know $`\mathrm{ker}𝔈\mathrm{ker}𝔇`$, and we will see that the containment is not always an equality (though it would be over $``$ or $``$). However, by checking $`𝔇𝐲=0`$ directly, we always find in practice that the $`𝐲`$ we compute lie in $`\mathrm{ker}𝔇`$.
Thus each successful run of our algorithm produces one kernel vector for $`𝔇`$. It also produces $`k`$. One can easily show $`kdim\mathrm{ker}𝔈`$. We run Lanczos up to 30 or 40 times with different random seeds, and we find $`k`$ is independent of the random seed used.<sup>3</sup><sup>3</sup>3For one level, random seeds produced a certain value of $`k`$, while one random seed produced a value that was greater by $`1`$. The Lanczos method behaves this way when $`E=\{𝔈𝐯,𝔈^2𝐯,𝔈^3𝐯,\mathrm{}\}`$ does not span the image of $`𝔈`$, but only a proper subspace of the image. Since $`𝐯`$ is chosen randomly, it is extremely rare for $`E`$ to span less than the full image; our data bears this statement out. We conclude that $`dim\mathrm{ker}𝔈=k`$; though we have not proved this, the computational evidence seems conclusive. It is clear that $`dim\mathrm{ker}𝔈dim\mathrm{ker}𝔇`$.
To find a basis of $`\mathrm{ker}𝔇`$, we run Lanczos many times until we have a set $`S`$ of $`k+10`$ elements of $`\mathrm{ker}𝔇`$. We use mod $`p`$ Gram-Schmidt on subsets $`S^{}S`$ to find maximal linearly independent subsets of $`S`$. We start with several different $`S^{}`$’s. In all the cases we checked, we found that maximal linearly indepdent sets in $`S`$ had a common cardinality $`k^{}`$, that any $`k^{}`$-element subset of $`S`$ was linearly independent, and that any subset of $`S`$ with more than $`k^{}`$ elements was dependent. We conclude that $`k^{}`$ is the value of $`dim\mathrm{ker}𝔇`$; again, the computational evidence is convincing, though all we have proved is $`k^{}dim\mathrm{ker}𝔇`$. In the rows marked ”Lanczos” in Table 1, what we report as “rank” is $`k^{}`$. In the rows marked $``$ or $``$, what we report (namely $`k^{}`$) is provably the rank of $`H^5(\mathrm{\Gamma },)`$. We do find $`k^{}<k`$ sometimes in practice.
### 5.4.
In Table 1, we give the results of our Betti number computations. Gauß means that ordinary Gaussian elimination was used to find the kernel of $`𝔇`$, and Lanczos means the algorithm of §5.3 was used. The entries marked with $``$ are those for which $`k=k^{}`$.
### 5.5.
Next, we present the Hecke data we computed. These computations are much more arduous than computing Betti numbers, and grow in complexity very fast as a function of the number of cells of the retract $`W`$ mod $`\mathrm{\Gamma }`$. Hence we were able to compute only a few Hecke operators, usually only $`T(2,)`$ and $`T(3,)`$. Beyond level $`20`$, it becomes infeasible to compute $`T(3,)`$; hence most of our data at large levels is only for $`T(2,)`$. Happily this is usually sufficient to guess persuasively what is happening with the cohomology.
We use the following conventions and abbreviations. All polynomials should be considered as elements of $`\overline{𝔽}[X]`$, where the bar denotes algebraic closure. We denote the $`p`$-adic cyclotomic character of $`G_{}`$ by $`ϵ`$, so that $`ϵ(\mathrm{Frob}_l)=l`$ for any $`lp`$. The symbol IIa (resp. IIb, IV) denotes a Galois representation of the form $`ϵ^a\sigma _kϵ^bϵ^c`$, where $`(k,a,b,c)`$ is $`(2,0,2,3)`$ (resp. $`(2,2,0,1)`$, $`(4,0,1,2)`$), if $`\sigma _k`$ is the Galois representation associated to a weight $`k`$ classical holomorphic cuspidal newform $`fS_k^{\text{new}}(N^{})`$, where $`N^{}`$ divides $`N`$. The same symbol prefixed with an ”E” denotes a Galois representation of the same form except that $`\sigma _k`$ is the Galois representation attached to an Eisenstein series of weight $`k`$ for $`\mathrm{\Gamma }_1(N^{})`$.
The individual tables are organized as follows. For each level we give the rank, as defined at the end of §5.3. Then each block gives data for the Hecke eigenspaces. The first column gives the type Galois representation seemingly attached to this eigenspace. The second column gives the dimension of this eigenspace. The third column gives the index $`l`$ of the Hecke operator, and the fourth column gives the corresponding factored *Hecke polynomial*. This is the polynomial defined on the right of (1) in Definition 1.2; it succinctly encodes the Hecke action on any vector in the eigenspace. At the bottom of each table, we indicate the $``$-splitting of the spaces of newforms $`S_k=S_k^{\text{new}}(N)`$ under the action of the Hecke operators; this data is from . After each table, we comment of the eigenclasses.
| Level 11. $`\text{rank}=2`$. | | | |
| --- | --- | --- | --- |
| IIa | 1 | $`T_2`$ | $`(14X)(18X)(1+2X+2X^2)`$ |
| | | $`T_3`$ | $`(19X)(127X)(1+X+3X^2)`$ |
| | | $`T_5`$ | $`(125X)(1125X)(1X+5X^2)`$ |
| | | $`T_7`$ | $`(149X)(1343X)(1+2X+7X^2)`$ |
| IIb | 1 | $`T_2`$ | $`(1X)(12X)(1+8X+32X^2)`$ |
| | | $`T_3`$ | $`(1X)(13X)(1+9X+243X^2)`$ |
| | | $`T_5`$ | $`(1X)(15X)(125X+3125X^2)`$ |
| | | $`T_7`$ | $`(1X)(17X)(1+98X+16807X^2)`$ |
| $`dimS_2(11)=1`$, $`dimS_4(11)=2`$ | | | |
The weight 4 newform doesn’t lift.
| Level 13. $`\text{rank}=1`$. | | | |
| --- | --- | --- | --- |
| IV | 1 | $`T_2`$ | $`(12X)(14X)(1+5X+8X^2)`$ |
| | | $`T_3`$ | $`(13X)(19X)(1+7X+27X^2)`$ |
| | | $`T_5`$ | $`(15X)(125X)(1+7X+125X^2)`$ |
| | | $`T_7`$ | $`(17X)(149X)(1+13X+343X^2)`$ |
| $`dimS_2(13)=0`$, $`dimS_4(13)=1+2`$ | | | |
Only the rational weight 4 newform lifts.
| Level 14. $`\text{rank}=2`$. | | | |
| --- | --- | --- | --- |
| IIa | 1 | $`T_3`$ | $`(19X)(127X)(1+2X+3X^2)`$ |
| IIb | 1 | $`T_3`$ | $`(1X)(13X)(1+18X+243X^2)`$ |
| $`dimS_2(14)=1`$, $`dimS_4(14)=1+1`$ | | | |
No weight four newforms lift.
| Level 15. $`\text{rank}=2`$. | | | |
| --- | --- | --- | --- |
| IIa | 1 | $`T_2`$ | $`(14X)(18X)(1+X+2X^2)`$ |
| IIb | 1 | $`T_2`$ | $`(1X)(12X)(1+4X+32X^2)`$ |
| $`dimS_2(15)=1`$, $`dimS_4(15)=1+1`$ | | | |
No weight four newforms lift.
| Level 17. $`\text{rank}=2`$. | | | |
| --- | --- | --- | --- |
| IIa | 1 | $`T_2`$ | $`(14X)(18X)(1+X+2X^2)`$ |
| | | $`T_3`$ | $`(19X)(127X)(1+3X^2)`$ |
| | | $`T_5`$ | $`(125X)(1125X)(1+2X+5X^2)`$ |
| | | $`T_7`$ | $`(149X)(1343X)(14X+7X^2)`$ |
| IIb | 1 | $`T_2`$ | $`(1X)(12X)(1+4X+32X^2)`$ |
| | | $`T_3`$ | $`(1X)(13X)(1+243X^2)`$ |
| | | $`T_5`$ | $`(1X)(15X)(1+50X+3125X^2)`$ |
| | | $`T_7`$ | $`(1X)(17X)(1196X+16807X^2)`$ |
| $`dimS_2(17)=1`$, $`dimS_4(17)=1+3`$ | | | |
No weight four newforms lift.
| Level 19. $`\text{rank}=3`$. | | | |
| --- | --- | --- | --- |
| IIa | 1 | $`T_2`$ | $`(14X)(18X)(1+2X^2)`$ |
| | | $`T_3`$ | $`(19X)(127X)(1+2X+3X^2)`$ |
| | | $`T_5`$ | $`(125X)(1125X)(13X+5X^2)`$ |
| IIb | 1 | $`T_2`$ | $`(1X)(12X)(1+32X^2)`$ |
| | | $`T_3`$ | $`(1X)(13X)(1+18X+243X^2)`$ |
| | | $`T_5`$ | $`(1X)(15X)(175X+3125X^2)`$ |
| IV | 1 | $`T_2`$ | $`(12X)(14X)(1+3X+8X^2)`$ |
| | | $`T_3`$ | $`(13X)(19X)(1+5X+27X^2)`$ |
| | | $`T_5`$ | $`(15X)(125X)(1+12X+125X^2)`$ |
| $`dimS_2(19)=1`$, $`dimS_4(19)=1+3`$ | | | |
The rational weight four newform lifts.
| Level 20. $`\text{rank}=2`$. | | | |
| --- | --- | --- | --- |
| IIa | 1 | $`T_3`$ | $`(19X)(127X)(1+2X+3X^2)`$ |
| IIb | 1 | $`T_3`$ | $`(1X)(13X)(1+18X+243X^2)`$ |
| $`dimS_2(20)=1`$, $`dimS_4(20)=1`$ | | | |
The rational weight four newform doesn’t lift.
| Level 21. $`\text{rank}=3`$. | | | |
| --- | --- | --- | --- |
| IIa | 1 | $`T_2`$ | $`(14X)(18X)(1+X+2X^2)`$ |
| IIb | 1 | $`T_2`$ | $`(1X)(12X)(1+4X+32X^2)`$ |
| IV | 1 | $`T_2`$ | $`(12X)(14X)(1+3X+8X^2)`$ |
| $`dimS_2(21)=1`$, $`dimS_4(21)=1+1+2`$ | | | |
Of the weight four rational newforms, only one lifts.
| Level 23. $`\text{rank}=5`$. | | | |
| --- | --- | --- | --- |
| IIa | 2 | $`T_2`$ | $`(14X)(18X)(1\alpha X+2X^2)`$ |
| IIb | 2 | $`T_2`$ | $`(1X)(12X)(14\alpha X+32X^2)`$ |
| IV | 1 | $`T_2`$ | $`(12X)(14X)(1+2X+8X^2)`$ |
| $`dimS_2(23)=2`$, $`dimS_4(23)=1+4`$ | | | |
Here $`\alpha `$ satisfies $`\alpha ^2\alpha +1=0`$. The rational weight four newform lifts.
| Level 25. $`\text{rank}=7`$. | | | |
| --- | --- | --- | --- |
| IV | 1 | $`T_2`$ | $`(12X)(14X)(1+X+8X^2)`$ |
| EIIa | 2 | $`T_2`$ | $`(14X)(18X)(1+\alpha X+2X^2)`$ |
| EIV | 2 | $`T_2`$ | $`(12X)(14X)(1+\beta X+8X^2)`$ |
| EIIb | 2 | $`T_2`$ | $`(1X)(12X)(1+\gamma X+32X^2)`$ |
| $`dimS_2(25)=0`$, $`dimS_4(25)=1+1+1`$ | | | |
Of the three weight four rational newforms, only one lifts. Here $`\alpha `$ satisfies $`\alpha ^2+1=0`$, $`\beta `$ satisfies $`\beta ^2+49=0`$, and $`\gamma `$ satisfies $`\gamma ^2+16=0`$.
| Level 27. $`\text{rank}=12`$. | | | |
| --- | --- | --- | --- |
| IIa | 1 | $`T_2`$ | $`(14X)(18X)(1+2X^2)`$ |
| IIb | 1 | $`T_2`$ | $`(1X)(12X)(1+32X^2)`$ |
| IV | 1 | $`T_2`$ | $`(12X)(14X)(1+3X+8X^2)`$ |
| EIIa | 3 | $`T_2`$ | $`(14X)(18X)(1+X)(1+2X)`$ |
| EIV | 3 | $`T_2`$ | $`(12X)(14X)(1+X)(1+8X)`$ |
| EIIb | 3 | $`T_2`$ | $`(1X)(12X)(1+4X)(1+8X)`$ |
| $`dimS_2(27)=1`$, $`dimS_4(27)=1+1+2`$ | | | |
Of the two weight four rational newforms, only one lifts.
| Level 29. $`\text{rank}=6`$. | | | |
| --- | --- | --- | --- |
| IIa | 2 | $`T_2`$ | $`(14X)(18X)(1\alpha X+2X^2)`$ |
| | | $`T_3`$ | $`(19X)(127X)(1+\alpha X+3X^2)`$ |
| IIb | 2 | $`T_2`$ | $`(1X)(12X)(14\alpha X+32X^2)`$ |
| | | $`T_3`$ | $`(1X)(13X)(1+9\alpha X+243X^2)`$ |
| IV | 2 | $`T_2`$ | $`(12X)(14X)(1\alpha X+8X^2)`$ |
| | | $`T_3`$ | $`(13X)(19X)(1(3\alpha +8)X+27X^2)`$ |
| $`dimS_2(29)=2`$, $`dimS_4(29)=2+5`$ | | | |
Here $`\alpha `$ satisfies $`\alpha ^2+2\alpha 1=0`$. In this example, the weight two and weight four newforms that lift are defined over the same quadratic extension of $``$.
| Level 31. $`\text{rank}=6`$. | | | |
| --- | --- | --- | --- |
| IIa | 2 | $`T_2`$ | $`(14X)(18X)(1\alpha X+2X^2)`$ |
| | | $`T_3`$ | $`(19X)(127X)(1+2\alpha X+3X^2)`$ |
| IIb | 2 | $`T_2`$ | $`(1X)(12X)(14\alpha X+32X^2)`$ |
| | | $`T_3`$ | $`(1X)(13X)(1+18\alpha X+243X^2)`$ |
| IV | 2 | $`T_2`$ | $`(12X)(14X)(1\beta X+8X^2)`$ |
| | | $`T_3`$ | $`(13X)(19X)(1+(2\beta +6)X+27X^2)`$ |
| $`dimS_2(31)=2`$, $`dimS_4(31)=2+5`$ | | | |
Here $`\alpha `$ satisfies $`\alpha ^2\alpha 1=0`$, and $`\beta `$ satisfies $`\beta ^2+5\beta +2=0`$.
| Level 33. $`\text{rank}=10`$. | | | |
| --- | --- | --- | --- |
| IIa | 1 | $`T_2`$ | $`(14X)(18X)(1X+2X^2)`$ |
| IIb | 1 | $`T_2`$ | $`(1X)(12X)(14X+32X^2)`$ |
| IV | 1 | $`T_2`$ | $`(12X)(14X)(1+5X+8X^2)`$ |
| IV | 1 | $`T_2`$ | $`(12X)(14X)(1+X+8X^2)`$ |
| IIa | 3 | $`T_2`$ | $`(14X)(18X)(1+2X+2X^2)`$ |
| IIb | 3 | $`T_2`$ | $`(1X)(12X)(1+8X+32X^2)`$ |
| $`dimS_2(33)=1`$, $`dimS_4(33)=1+1+2+2`$ | | | |
The three dimensional eigenspaces are lifts of the weight two newform from level $`11`$.
| Level 35. $`\text{rank}=7`$. | | | |
| --- | --- | --- | --- |
| IIa | 1 | $`T_2`$ | $`(14X)(18X)(1+2X^2)`$ |
| IIb | 1 | $`T_2`$ | $`(1X)(12X)(1+32X^2)`$ |
| IIa | 2 | $`T_2`$ | $`(14X)(18X)(1\alpha X+2X^2)`$ |
| IIb | 2 | $`T_2`$ | $`(1X)(12X)(14\alpha X+32X^2)`$ |
| IV | 1 | $`T_2`$ | $`(12X)(14X)(1X+8X^2)`$ |
| $`dimS_2(35)=1+2`$, $`dimS_4(35)=1+2+3`$ | | | |
Here $`\alpha `$ satisfies $`\alpha ^2+\alpha 4=0`$.
| Level 37. $`\text{rank}=8`$. | | | |
| --- | --- | --- | --- |
| IIa | 1 | $`T_2`$ | $`(14X)(18X)(1+2X+2X^2)`$ |
| IIb | 1 | $`T_2`$ | $`(1X)(12X)(1+8X+32X^2)`$ |
| IIa | 1 | $`T_2`$ | $`(14X)(18X)(1+2X^2)`$ |
| IIb | 1 | $`T_2`$ | $`(1X)(12X)(1+32X^2)`$ |
| IV | 4 | $`T_2`$ | $`(12X)(14X)(1\alpha X+8X^2)`$ |
| $`dimS_2(37)=1+1`$, $`dimS_4(37)=4+5`$ | | | |
Here $`\alpha `$ satisfies $`\alpha ^4+6\alpha ^3\alpha ^216\alpha +6=0`$.
| Level 39. $`\text{rank}=10`$. | | | |
| --- | --- | --- | --- |
| IIa | 1 | $`T_2`$ | $`(14X)(18X)(1X+2X^2)`$ |
| IIb | 1 | $`T_2`$ | $`(1X)(12X)(14X+32X^2)`$ |
| IIa | 2 | $`T_2`$ | $`(14X)(18X)(1\alpha X+2X^2)`$ |
| IIb | 2 | $`T_2`$ | $`(1X)(12X)(14\alpha X+32X^2)`$ |
| IV | 1 | $`T_2`$ | $`(12X)(14X)(1+8X^2)`$ |
| IV | 3 | $`T_2`$ | $`(12X)(14X)(1+5X+8X^2)`$ |
| $`dimS_2(39)=1+2`$, $`dimS_4(39)=1+2+2`$ | | | |
Here $`\alpha `$ satisfies $`\alpha ^2+2\alpha 6=0`$. The three dimensional eigenspaces are lifts of the weight four newform from level $`13`$.
| Level 41. $`\text{rank}=9`$. | | | |
| --- | --- | --- | --- |
| IIa | 3 | $`T_2`$ | $`(14X)(18X)(1\alpha X+2X^2)`$ |
| IIb | 3 | $`T_2`$ | $`(1X)(12X)(14\alpha X+32X^2)`$ |
| IV | 3 | $`T_2`$ | $`(12X)(14X)(1\beta X+8X^2)`$ |
| $`dimS_2(41)=3`$, $`dimS_4(41)=3+7`$ | | | |
Here $`\alpha `$ satisfies $`\alpha ^3+\alpha ^25\alpha 1=0`$, and $`\beta `$ satisfies $`\beta ^3+3\beta ^25\beta 3=0`$.
| Level 43. $`\text{rank}=10`$. | | | |
| --- | --- | --- | --- |
| IIa | 1 | $`T_2`$ | $`(14X)(18X)(1+2X+2X^2)`$ |
| IIb | 1 | $`T_2`$ | $`(1X)(12X)(1+8X+32X^2)`$ |
| IIa | 2 | $`T_2`$ | $`(14X)(18X)(1\alpha X+2X^2)`$ |
| IIb | 2 | $`T_2`$ | $`(1X)(12X)(14\alpha X+32X^2)`$ |
| IV | 4 | $`T_2`$ | $`(12X)(14X)(1\beta X+8X^2)`$ |
| $`dimS_2(43)=1+2`$, $`dimS_4(43)=4+6`$ | | | |
Here $`\alpha `$ satisfies $`\alpha ^22=0`$, and $`\beta `$ satisfies $`\beta ^4+4\beta ^39\beta ^214\beta +2=0`$.
| Level 47. $`\text{rank}=11`$. | | | |
| --- | --- | --- | --- |
| IIa | 4 | $`T_2`$ | $`(14X)(18X)(1\alpha X+2X^2)`$ |
| IIb | 4 | $`T_2`$ | $`(1X)(12X)(14\alpha X+32X^2)`$ |
| IIa | 3 | $`T_2`$ | $`(14X)(18X)(1\beta X+2X^2)`$ |
| $`dimS_2(47)=4`$, $`dimS_4(47)=3+8`$ | | | |
Here $`\alpha `$ satisfies $`\alpha ^4\alpha ^35\alpha ^2+5\alpha 1=0`$, and $`\beta `$ satisfies $`\beta ^3+5\beta ^22\beta 12=0`$.
## 6. Interpretation of the numerical results
### 6.1.
The first step, for each Hecke eigenvector $`\beta H_5(\mathrm{\Gamma },)`$ and for each prime $`l`$, is to write down and factor the Hecke polynomial $`P_l(X)`$.
We then see that for the data computed so far, $`\beta `$ has one of the following Galois representations attached. As above, $`p`$ is a prime not dividing $`N`$ or any of the $`l`$’s we are looking at.
Let $`k=2`$ or $`4`$ and consider the continuous semisimple representation $`\sigma _k:G_{}\mathrm{GL}_2(_p)`$ unramified outside $`pN`$ attached to a classical Hecke eigenform $`f`$ of weight $`k`$ and level $`N^{}`$ dividing $`N`$. If $`N`$ is prime, we also assume $`f`$ has trivial nebentypus. Let $`\rho =ϵ^a\sigma _kϵ^bϵ^c`$ where $`(k,a,b,c)=(2,0,2,3)`$ or $`(2,2,0,1)`$, or $`(4,0,1,2)`$. If $`f`$ is an Eisenstein series, then $`\rho =\chi _0\chi _1ϵ\chi _2ϵ^2\chi _3ϵ^3`$, where the $`\chi _i`$ are Dirichlet characters of conductor dividing $`N`$ with values in a finite extension of $`𝔽`$, at least two of which are trivial.
Then for any $`\beta `$ there is some choice of such $`\rho `$ which is apparently attached to $`\beta `$, in the sense that the Hecke polynomial at $`l`$ equals the characteristic polynomial of $`\mathrm{Frob}_l`$ for all $`l`$ for which we computed the Hecke eigenvalues.
Thus it appears that none of our computed classes so far is cuspidal. Therefore we should be able to related them to cohomology of the boundary, either geometrically or in terms of Eisenstein series. We cannot do this thoroughly, because neither the cohomology of the Borel-Serre boundary nor the theory of Eisenstein cohomology has been sufficiently worked out for $`\mathrm{GL}_4/`$. This is not an easy task. We can give the following indications.
From results of Moeglin-Waldspurger , we don’t expect any of our classes to be residues of Eisenstein series. In the framework of we can guess that our classes lie either in the part of the cohomology indexed by the associate class of parabolic subgroups of $`\mathrm{GL}_4`$ of type $`(2,1,1)`$ or in the part indexed by the Borel subgroup. Here the cuspidal data on the $`\mathrm{GL}_2`$-factor of the Levi component of the first parabolic comes from the appropriate classical cuspform of weight 2 or 4, and we use the appropriate power of the determinant on the $`\mathrm{GL}_1`$ factors.
### 6.2.
Geometrically, we make the following comments. Let $`M`$ be the quotient of the symmetric space for $`\mathrm{SL}_4()`$ by $`\mathrm{\Gamma }`$ and let $`M`$ be the boundary of its Borel-Serre compactification. The covering of $`M`$ by its faces gives a spectral sequence for its cohomology. The $`E_2`$ page has for its $`(i,j)`$-th term $`H^i(\text{Tits building}/\mathrm{\Gamma },H^j(\text{Fiber}))`$. An element of that is an assignment: to every face $`e^{}(P)`$ of codimension $`i`$ we assign an element of the cohomology in degree $`j`$ of $`P\mathrm{\Gamma }`$. These assignments when restricted to a common face of codimension $`i+1`$ must add up to 0. Such an assignment gives a class in $`E_2^{i,j}(M)`$. If it persists in the spectral sequence to $`E_{\mathrm{}}`$, it will contribute to the cohomology $`H^{i+j}(M,)`$. There remains the question as to whether this contribution is the restriction of a class in $`H^{i+j}(M,)`$.
Note that if $`P=LU`$ is a Levi decomposition of $`P`$ then the spectral sequence of the fibration for the cohomology of $`P\mathrm{\Gamma }`$ corresponding to this decomposition is known to degenerate at $`E_2^{p,q}=H^p(\mathrm{\Gamma }_L,H^q(U\mathrm{\Gamma })`$ where $`\mathrm{\Gamma }_L`$ is the projection of $`\mathrm{\Gamma }`$ to $`L`$.
The classes we have computed so far we expect to be coming in this way from $`M`$. From the shape of the apparently associated Galois representations, here is what we believe is their origin. We only sketch the constructions, since a detailed description would require a thorough investigation of the cohomology of $`M`$ for arbitrary congruence subgroups of $`\mathrm{SL}_4()`$. First assume $`f`$ is a cuspform.
### 6.3.
The case where $`\sigma `$ has weight 2: By the Eichler-Shimura theorem, the cuspform $`f`$ that has $`\sigma `$ attached shows up as a class $`\alpha `$ in $`H^1(\mathrm{\Delta },)`$, where $`\mathrm{\Delta }`$ is the classical $`\mathrm{\Gamma }_0(N^{})\mathrm{SL}_2()`$. First suppose $`N^{}=N`$. Consider the following element of $`E_2^{0,5}(M)`$: On the standard parabolic subgroup of type $`(2,2)`$ which is the stabilizer of the span of $`(e_1,e_2)`$ in 4-space, we put the cohomology class $`\alpha \times 1`$ on the Levi component where we view the trivial coefficients of $`\alpha `$ as the module $`H^4(U\mathrm{\Gamma })`$. One sees that there is a unique class in the appropriate face corresponding to a $`(3,1)`$-type parabolic subgroup that has the same restriction to the type $`(2,1,1)`$ face they have in common and restricts to 0 on the other faces. Hence these two glue up to give a class in the boundary. The same construction with the transposed parabolic subgroup also gives a class, and these seem to account for all our $`\beta `$’s falling under this case. If $`N^{}N`$, we choose a $`(2,2)`$\- parabolic subgroup $`P`$ such that $`\mathrm{\Gamma }_L`$ has level $`N^{}`$. Then we imitate the construction above.
### 6.4.
The case where $`\sigma `$ has weight 4: Here not every class we construct on the boundary seems to lift to $`H^5(M)`$, but only some of them. We don’t know the reason for this. The construction here creates a class in $`E_2^{1,4}(M)`$, and if this correctly describes what we have computed, our computed classes of this type would be ghost classes. That is, the corresponding cohomology class in $`H^5(M)`$ restricts nontrivially to the boundary of $`M`$, but it restricts to 0 on each face of the boundary, since it is coming from a class in $`E_2^{i,j}(M)`$ with $`i>0`$.
For this construction, one first chooses a parabolic subgroup $`P`$ of type $`(2,1,1)`$. Note that $`H^3((U\mathrm{\Gamma }),)`$ contains a $`\mathrm{\Gamma }_L`$-submodule isomorphic to $`𝒱_2`$, the homogeneous complex polynomials of degree 2 on 2-space, after we identify $`\mathrm{\Gamma }_L`$ with a subgroup of $`\mathrm{GL}_2()\times \mathrm{GL}_1()\times \mathrm{GL}_1()`$. By the Eichler-Shimura theorem, the weight $`4`$ cuspform $`f`$ that has $`\sigma `$ attached shows up in $`H^1(\mathrm{\Delta },𝒱_2)`$. Thus we can view it as in $`H^1(\mathrm{\Gamma }_L,H^3(U\mathrm{\Gamma }))`$.
We can do this on three $`P`$’s which are not conjugate to each other in such a way that they, together with certain classes on $`(3,1)`$-type parabolic subgroups, all glue together to give a class in $`E_2^{1,4}(M)`$. The details are left for the reader.
### 6.5.
Finally, the case where $`\sigma `$ is an Eisenstein series is harder to understand. We haven’t worked out exactly how the gluing process goes in this case, so we’re not sure what stratum of the spectral sequence is occupied by the corresponding boundary classes.
We note that in every case, $`\rho `$ restricted to an inertia subgroup of $`G_{}`$ at $`p`$ has the form $`1ϵϵ^2ϵ^3`$, which is consistent with the conjecture of Ash-Sinnott , since the coefficient module of the cohomology classes we consider is the trivial module.
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# 1 Introduction
## 1 Introduction
There is now strong evidence in support of an oscillatory behavior of neutrinos. The results on atmospheric neutrinos from SuperKamiokande find a comprehensive explanation in terms of oscillations, indicating a non-zero neutrino mass. This result and the solar neutrino problem, which has been held as a signal for neutrino oscillations for long, have created much excitement in the community.
In this paper we examine the latest solar neutrino data from SuperKamiokande along with those from the radiochemical Chlorine and Gallium experiments assuming that MSW resonant flavor conversion is operative. The above experiments have presented the measured arrival rates which, in all the cases, are less than the predictions from the Standard Solar Model (SSM). In addition, SuperKamiokande (SK) has provided the observed electron energy distribution . We use these data to test the consistency of the MSW mechanism taken together with the SSM predictions in a two flavor scenario. There are several recent MSW analyses of the SK solar neutrino data \- . Although the data fitting method followed in all of them is to minimize a $`\chi ^2`$ function, the details of the statistical procedure used vary among the different groups. In this work, we indicate two different ways of performing the statistical analysis for the spectrum and the combined rate and spectrum data and check the consistency of the best-fit values of the mass squared differences and the mixing angles so obtained. For the analysis of the spectrum results we explore the possibility of using its moments as variables for fitting the data. The main advantage of these moments is that they probe the shape of the spectrum in a manner independent of the <sup>8</sup>B flux normalization uncertainties. For the combined analysis of the rates and the spectrum, apart from the standard procedure of varying the <sup>8</sup>B flux normalization and treating these two sets of data as independent, we also adopt a second method which takes into account the correlations among the rates and the spectrum data.
In addition to the best-fit values of the oscillation parameters, we present the 99% C.L. and 90% C.L. allowed regions and the goodness of fit (g.o.f) of a particular solution. By g.o.f. is meant the probability that the $`\chi ^2`$ will exceed $`\chi _{min}^2`$. When presenting the allowed region we take $`\chi _{min}^2`$ to be the value at the global minimum in that region<sup>1</sup><sup>1</sup>1The other approach is to present the allowed regions with respect to the local minimum.. For the neutrino fluxes and the neutrino production positions within the sun we use the BP98 solar model . We consider oscillation of $`\nu _e`$ to a sequential ($`\nu _\mu `$ or $`\nu _\tau `$) neutrino.
This paper is organized as follows. In the next section we present the formulae for oscillation of neutrinos with the inclusion of matter effects both in the sun and during their passage through the earth. In section 3 we use the data on the total solar neutrino rates as measured at the Chlorine, Gallium, and SuperKamiokande (1117-day data) detectors to obtain the best-fit values of the neutrino mass splitting and the mixing angles. In section 4 we consider the electron energy spectrum observed at SuperKamiokande. Using the MSW predictions, we obtain the best-fit values from a direct fit to the data as well as from a fit to the normalized moments. In section 5 we use both the total rates data and the SK electron energy spectrum data to make a combined fit. As noted earlier, here we allow the normalization of the <sup>8</sup>B spectrum to vary and compare the results with those obtained when the SSM prediction for this normalization is used, allowing the inclusion of correlations between the total rates and the observed spectrum via astrophysical uncertainties. We end in section 7 with a summary, discussions, and conclusions.
## 2 Oscillation Probability
In this work we restrict ourselves to the simplest case of mixing between two neutrino flavors. Assuming the neutrino mass eigenstates, $`\nu _i`$, reaching the earth to be incoherent , the survival probabaility of an electron neutrino can be written as
$$P_{ee}=\underset{i}{}P_{ei}^SP_{ie}^E,$$
(1)
where $`P_{ei}^S`$ is the probability of an electron neutrino state to transform into the $`i`$-th mass state at the solar surface and $`P_{ie}^E`$ is the conversion probability of the $`i`$-th mass state to the $`\nu _e`$ state after traversing the earth. One can express $`P_{ee}`$ in terms of the day-time (i.e. no earth matter effect) probability $`P_{ee}^D`$ as<sup>2</sup><sup>2</sup>2This expression is not applicable for maximal mixing ($`\mathrm{cos}2\theta `$ =0) and for this case we use eq. (1) directly.
$$P_{ee}=P_{ee}^D+\frac{(2P_{ee}^D1)(\mathrm{sin}^2\theta P_{2e}^E)}{\mathrm{cos}2\theta },$$
(2)
where
$$P_{ee}^D=0.5+[0.5\mathrm{\Theta }(E_\nu E_A)X]\mathrm{cos}2\theta _M\mathrm{cos}2\theta ,$$
(3)
$`\mathrm{\Theta }`$ being the Heaviside function. $`\theta `$ is the mixing angle in vacuum and $`\theta _M`$ is the mixing angle in matter given by
$$\mathrm{tan}2\theta _M=\frac{\mathrm{\Delta }m^2\text{ }\mathrm{sin}2\theta }{\mathrm{\Delta }m^2\text{ }\mathrm{cos}2\theta 2\sqrt{2}G_Fn_eE_\nu }.$$
(4)
Here $`n_e`$ is the ambient electron density, $`E_\nu `$ the neutrino energy, and $`\mathrm{\Delta }m^2`$ (= $`m_2^2m_1^2`$) the mass squared difference in vacuum.
$$E_A=\mathrm{\Delta }m^2\mathrm{cos}2\theta /2\sqrt{2}G_Fn_e|_{pr},$$
(5)
gives the minimum $`\nu _e`$ energy that can encounter a resonance inside the sun, $`n_e|_{pr}`$ being the electron density at the point of production. $`X`$ is the jump-probability between the mass eigenstates and for an exponential density profile, as is approximately the case in the sun, it is given by
$$X=\frac{\mathrm{exp}[\pi \gamma _R(1\mathrm{cos}2\theta )]\mathrm{exp}[2\pi \gamma _R]}{1\mathrm{exp}[2\pi \gamma _R]},$$
(6)
where $`\gamma _R=\gamma \mathrm{cos}2\theta /\mathrm{sin}^22\theta `$ and
$$\gamma =\frac{\pi }{4}\frac{\mathrm{\Delta }m^2}{E_\nu }\frac{\mathrm{sin}^22\theta }{\mathrm{cos}2\theta }\frac{1}{|\frac{d\mathrm{ln}n_e}{dr}|_{res}}.$$
(7)
For calculating $`P_{2e}^E`$ in eq. (2) we treat the earth as a slab of constant density (4.5 gm/cc). We have verified that this is a reasonable approximation since the location of the SK detector ensures that only in a rather small fraction of the time does the neutrino pass through the denser core.
The definition of $`\chi ^2`$ used by us in the following sections is,
$$\chi ^2=\underset{i,j}{}\left(F_i^{th}F_i^{exp}\right)(\sigma _{ij}^2)\left(F_j^{th}F_j^{exp}\right).$$
(8)
Here $`F_i^\xi =T_i^\xi /T_i^{BP98}`$ where $`\xi `$ is $`th`$ (for the theoretical prediction with oscillations) or $`exp`$ (for the experimental value) and $`T_i`$ stand for the quantities being fit (total rates from different experiments, electron energy spectrum for different energy bins, etc.). The error matrix $`\sigma _{ij}`$ contains the experimental errors, the theoretical errors and their correlations.
## 3 Total rates
In this section we perform an analysis of the total rates as measured at the various experiments. The data that we use for the total rates are given in Table 1. In particular, we use the 1117-day SK data. For the Ga experiments we take the weighted average of the SAGE and Gallex results. Because SK has better statistics we do not include the Kamiokande results.
To get the best-fit values of the parameters, we minimise a $`\chi ^2`$ function defined as in eq. (8). For evaluating the error matrix, $`\sigma _{ij}`$, we use the procedure described in . Our best-fit results for the total rates are summarized in Table 2. These fits have 1 degree of freedom (3 experimental data points – 2 parameters). It is clear from this Table that the small mixing angle solution for the sequential neutrino case gives by far the best fit to the total rates data.
It is seen from Table 2 that the best-fit values obtained for the 1117-day data are not markedly different from those obtained for the 825-day data. This indicates that the best-fit points obtained from the analysis of total rates are quite robust and are not expected to change drastically with more data from SK. However, the quality of fit is quite sensitive to these small changes: the g.o.f. has become a little poorer for the SMA solution while for both the LMA and LOW solutions it has improved with the accumulation of more SK data.
In Fig. 1 we show the 99% C.L. and 90% C.L. contours in the $`\mathrm{sin}^2\theta \mathrm{\Delta }m^2`$ plane for the MSW solution. There is no solution in the dark side ($`\mathrm{sin}^2\theta >0.5`$) from the analysis of total rates. Note that in allowed regions were obtained in the dark side from the total rates analysis while in no allowed region was found for $`\theta >\pi /4`$. We agree with ; a possible origin of this is that the same numerical density profile of the sun from BP98 has been used in these analyses. Our results for the total rates for the 1117 day data are in agreement with the analysis given in .
## 4 Observed spectrum and its moments
In addition to the total rates, SK has provided the number of events (normalized to the SSM prediction) in 18 electron recoil energy bins of width 0.5 MeV in the range 5.0 MeV to 14 MeV and a 19th bin which covers the events in the range 14 to 20 MeV . The systematics of the first bin are still under study and for our analysis we do not use it.
### 4.1 Observed spectrum
In this subsection we present the results obtained by directly fitting the SK spectral data. The theoretical predictions are calculated bin by bin and in the fitting procedure, in addition to the neutrino mixing parameters $`\mathrm{\Delta }m^2`$ and $`\mathrm{sin}^2\theta `$, we also allow the absolute normalization of the <sup>8</sup>B flux, $`X_B`$, to vary<sup>3</sup><sup>3</sup>3$`X_B`$ = 1, for the SSM.. The error matrix $`\sigma _{ij}`$ used by us (see eq. (8)) is
$$(\sigma _{ij}^2)_{sp}=\delta _{ij}(\sigma _{i,stat}^2+\sigma _{i,uncorr}^2)+\sigma _{i,exp}\sigma _{j,exp}+\sigma _{i,cal}\sigma _{j,cal},$$
(9)
where we have included the statistical error, the uncorrelated systematic errors and the energy-bin-correlated experimental errors as well as those from the calculation of the shape of the expected spectrum . Since we vary the normalization of the <sup>8</sup>B flux we do not include its astrophysical uncertainties separately.
The best-fit point from this analysis is found to be
* $`\mathrm{\Delta }m^2=2.29\times 10^6`$ eV<sup>2</sup>, $`\mathrm{sin}^2\theta =0.009`$, $`X_B=1.4`$, $`\chi _{min}^2=9.46`$, g.o.f. = 85.23%.
For these values of $`\mathrm{\Delta }m^2`$ and $`\mathrm{sin}^2\theta `$, choosing $`X_B=1`$, the data on total rates give a bad fit ($`\chi ^2`$ = 73.82) as the high energy <sup>8</sup>B neutrinos get suppressed more than observed. For the spectrum fit this problem can be avoided by a high $`X_B`$ $`>`$ 1. Since the <sup>8</sup>B flux normalization is allowed to vary, a large range of $`\mathrm{\Delta }m^2`$ and $`\mathrm{sin}^2\theta `$ including the $`\theta >\pi /4`$ region remains allowed by the spectral data as is shown in Fig. 2. Only the area inside the contours is disallowed at 90% C.L. from spectral data analysis.
To further examine this mismatch between the fits to the total rates and those to the SK spectral data, in Fig. 3 we present the behavior of $`\chi ^2`$ as obtained from the spectral analysis if we keep the parameters $`\mathrm{sin}^2\theta `$ and $`\mathrm{\Delta }m^2`$ in the SMA, LMA, and LOW regions of the fit to the total rates. In Figs. 3(a), 3(b), and 3(c) are shown the variation of the $`\chi ^2`$ from the spectral data with $`\mathrm{sin}^2\theta `$ lying in the SMA, LMA, and LOW regions respectively. For any chosen $`\mathrm{sin}^2\theta `$, we let $`\mathrm{\Delta }m^2`$ vary over the corresponding range permitted by the fit to the total rates at 99% C.L. (from Fig. 1) and plot the minimum value found. We consider two cases,
* the <sup>8</sup>B normalization is held fixed at its SSM value (solid curves)
* the <sup>8</sup>B normalization is permitted to vary (broken curves).
For Figs. 3(d), 3(e), and 3(f), the roles of $`\mathrm{sin}^2\theta `$ and $`\mathrm{\Delta }m^2`$ are interchanged. Fig. 3 indicates that if we allow the <sup>8</sup>B normalization to vary then the $`\mathrm{\Delta }m^2`$ and $`\mathrm{sin}^2\theta `$ allowed at 99% C.L. from the total rates are allowed at 90% C.L. from the spectral analysis.
Figs. 2 and 3 lead us to the conclusion that, at the present moment, the SK electron spectral data do not provide tight controls over the allowed parameter range.
### 4.2 Moments of spectrum
Since the absolute normalization of the <sup>8</sup>B flux is not precisely known, it is of interest to look for variables which probe neutrino oscillation effects in the data in a manner immune to this uncertainty. Normalized moments of the observed electron spectrum can be useful as one such set of variables . In practice, to compare with the data, it is convenient to standardize with respect to the SSM predictions by using
$$M_n=\frac{_i\left[\frac{N(E_i)}{\{N(E_i)\}_{SSM}}\right]E_i^n}{_i\left[\frac{N(E_i)}{\{N(E_i)\}_{SSM}}\right]},$$
(10)
where $`E_i`$ is the mean energy of the $`i`$-th bin and $`N(E_i)`$ is the number of events in this bin. Depending on whether the experimental or the theoretically predicted value of the variable is under consideration, $`N(E_i)`$ is obtained either from experiments or from the theoretical model under test. It is clear that these variables carry information about the shape of the neutrino spectrum which, if oscillations are operative, undergoes modification from the SSM prediction due to the energy dependence of the survival probability. It is obvious that the above moments are independent of the absolute normalization of the <sup>8</sup>B flux.
We have calculated the moments of the 1117-day data on the electron energy spectrum presented by SuperKamiokande . These are presented in Table 3. The error in the higher moments increases rapidly with the order and the ones beyond the sixth are not of much use.
Using these variables in a $`\chi ^2`$-analysis we find that fitting the first four moments results in the same best-fit values of $`\mathrm{\Delta }m^2`$ and $`\mathrm{sin}^2\theta `$ as those obtained using the first five or six moments. The best-fit region<sup>4</sup><sup>4</sup>4The $`\chi ^2`$ remains unchanged when $`\mathrm{sin}^2\theta `$ is varied over the indicated range is found to be
* $`\mathrm{\Delta }m^2=6.94\times 10^6`$ eV<sup>2</sup>, $`\mathrm{sin}^2\theta =0.0070.010`$, $`\chi _{min}^2=0.0001`$, g.o.f = 99.995%.
The very small value of $`\chi _{min}^2`$ for this fit should not be regarded as a major success of the theory but rather reflects the large errors associated with the moments as obtained from the present data. It is gratifying that the best-fit values obtained by this method are in the same broad region as those from fitting the recoil electron energy spectrum.
## 5 Combined fits to rates and spectrum
In this section we present the results of the combined fit to the total rates and the spectrum data. We have performed this global fit by the following two methods:
* We treat the rates and the electron spectrum data as independent. In this approach we vary the <sup>8</sup>B flux normalization as a free parameter.
* We fix the <sup>8</sup>B flux normalization at the SSM value (=1) and include the correlations of the <sup>8</sup>B flux uncertainty between the rates and spectrum data. To our knowledge, this approach has not been pursued in any previous analysis.
### 5.1 Fits using the <sup>8</sup>B flux normalization as a free parameter
For this case the definition of $`\chi ^2`$ is,
$`\chi ^2`$ $`=`$ $`{\displaystyle \underset{i,j=1,3}{}}\left(F_i^{th}F_i^{exp}\right)(\sigma _{ij}^2)\left(F_j^{th}F_j^{exp}\right)`$ (11)
$`+{\displaystyle \underset{i,j=1,18}{}}\left(X_BR_i^{th}R_i^{exp}\right)(\sigma _{ij}^2)_{sp}\left(X_BR_j^{th}R_j^{exp}\right),`$
where the first term on the r.h.s is from the fit to the total rates and the second from that to the spectral data. As we allow the normalization of the <sup>8</sup>B flux to vary as a free parameter we switch off the SSM astrophysical uncertainties arising because of this component. Since it is the <sup>8</sup>B flux that enters the rates as well as the spectrum data, in this manner of fitting the data the correlations between the rates and the spectrum are absent; the error matrix is block diagonal and one can treat $`\chi _{rate}^2`$ and $`\chi _{spectrum}^2`$ as independent. There are 18 (= 21 – 3) degrees of freedom in this case. The best-fit values we obtain are presented in Table 4.
In Fig. 4a we show the 99% and 90% C.L. allowed regions for the combined analysis of total rates and the observed electron spectrum. The best-fit points in these plots are obtained by varying $`X_B`$ in addition to $`\mathrm{\Delta }m^2`$ and $`\mathrm{sin}^2\theta `$.
### 5.2 Fits including correlations between rates and spectrum via <sup>8</sup>B flux
For this case we include the correlations between the theory errors in the rate and the spectrum data. This comes through the <sup>8</sup>B flux, as it enters both. Since we include the SSM astrophysical uncertainties in the <sup>8</sup>B flux the normalization factor for it is held fixed at the SSM value. Now the individual $`\chi ^2`$ due to the spectrum and the rates cannot be summed independently and the combined $`\chi ^2`$ is defined as,
$$\chi ^2=\underset{i,j=1,21}{}\left(F_i^{th}F_i^{exp}\right)(\sigma _{ij}^2)\left(F_j^{th}F_j^{exp}\right),$$
(12)
where the $`\sigma _{ij}`$ is now a 21 $`\times `$ 21 matrix defined in the following way,
* For $`i,j`$ = 1 …3
$$\sigma _{ij}^2=(\sigma _{ij}^2)_{th}+(\sigma _{ij}^2)_{exp},$$
(13)
where
$$(\sigma _{ij}^2)_{th}=\delta _{ij}\underset{\alpha =1}{\overset{8}{}}R_{\alpha i}^2(\mathrm{\Delta }C_{\alpha i})^2+\underset{\alpha ,\beta =1}{\overset{8}{}}R_{\alpha i}R_{\beta j}\underset{k=1}{\overset{e11}{}}a_{\alpha k}a_{\beta k}(\mathrm{\Delta }\mathrm{ln}X_k)^2.$$
(14)
where the first term is due to the cross-section uncertainties and the second term ($`\sigma _{ap}`$) is due to the astrophysical uncertainties . The off-diagonal elements in the error matrix come through $`\sigma _{ap}`$. $`R_{\alpha i}`$ denotes the contribution of the $`\alpha `$-th source to the rate of the $`i`$-th experiment. $`a_{\alpha k}=\delta \mathrm{ln}\varphi _\alpha /\delta \mathrm{ln}X_k`$, where $`\delta \mathrm{ln}\varphi _\alpha `$ is the error in the $`\alpha `$-th component of the spectrum due to the input parameter $`X_k`$ .
* For $`i`$ = 4 …21 and $`j`$ = 1 …3
$$\sigma _{ij}^2=\underset{\alpha =1}{\overset{8}{}}R_{{}_{}{}^{8}Bi}R_{\alpha j}\underset{k=1}{\overset{11}{}}a_{{}_{}{}^{8}Bk}a_{\alpha k}(\mathrm{\Delta }\mathrm{ln}X_k)^2.$$
(15)
* For $`i`$ = 1 …3 and $`j`$ = 4 …21
$$\sigma _{ij}^2=\underset{\alpha =1}{\overset{8}{}}R_{\alpha i}R_{{}_{}{}^{8}Bj}\underset{k=1}{\overset{11}{}}a_{\alpha k}a_{{}_{}{}^{8}Bk}(\mathrm{\Delta }\mathrm{ln}X_k)^2.$$
(16)
* for $`i`$ = 4 …21 and $`j`$ = 4 …21
$$\sigma _{ij}^2=(\sigma _{ij}^2)_{sp}+R_{{}_{}{}^{8}Bi}R_{{}_{}{}^{8}Bj}\underset{k=1}{\overset{11}{}}a_{{}_{}{}^{8}Bk}a_{{}_{}{}^{8}Bk}(\mathrm{\Delta }\mathrm{ln}X_k)^2.$$
(17)
In this case the number of degrees of freedom is 19 (= 21 – 2). The $`\chi _{min}^2`$ and the best-fit values we obtain are shown in Table 5.
We find that the fits are of poorer quality in this case as compared to the previous one.
In Fig. 4b we show the 99% and 90% C.L. allowed regions for the combined analysis of total rates and the observed electron spectrum for MSW conversion to sequential neutrinos including the correlations between the rates and the spectrum due to the astrophysical uncertainties of the <sup>8</sup>B flux normalization.
## 6 Summary, Discussions, and Conclusions
In this paper we have performed a detailed $`\chi ^2`$-analysis of the latest SK solar neutrino data together with the results from the Cl and Ga experiments in terms of two-generation MSW conversions of $`\nu _e`$ to sequential ($`\nu _\mu `$, $`\nu _\tau `$) neutrinos.
Compared to the recent analyses in the literature - there are two new features in our analysis.
* We fit the observed electron energy spectrum data in two different ways, exploring for the first time, the use of moments of the energy spectrum in a $`\chi ^2`$-analysis.
* The combined fits to the total rates and spectral data are also performed in two different manners. In the first, the <sup>8</sup>B flux normalization is used as a free parameter while in the other the SSM normalization is chosen for it and correlations between the rates and spectral data due to astrophysical uncertainties of the <sup>8</sup>B flux are included.
We find that the two-generation MSW scenario can well explain the data on total rates. The solution in the SMA (Small Mixing Angle) region is preferred over the other possibilities although the quality of the fit is poorer as compared to the one obtained using the 825-day SK data.
The best-fit from the spectrum data comes in a region disallowed from the total rates. In this region the <sup>8</sup>B neutrinos are suppressed much more than required by the rates data. For the analysis of the spectrum, the absolute normalization of the <sup>8</sup>B flux, $`X_B`$, has been permitted to be greater than unity, thus effectively compensating the shortfall. We have explored the use of normalized moments of the observed electron energy spectrum to signal MSW resonant flavour conversion. These variables are independent of the absolute normalization of the <sup>8</sup>B flux and probe the effect of oscillations on the spectral shape. This procedure is somewhat handicapped by the large errors on the moments calculated from the present data. However, the best-fit values obtained by the two methods are more or less in agreement.
Similarly, for the two methods followed in the combined $`\chi ^2`$ analysis of the rates and time averaged spectrum data, the best-fit values are not much different. The first approach gives a better fit because we utilise the freedom of varying the <sup>8</sup>B flux normalization. We remark that in the combined analysis, where the <sup>8</sup>B normalization is held fixed at the SSM value, the correlations between the rates and the spectrum data are found to be important and thus one should use caution regarding results obtained treating these as independent. For both methods, the best-fit from the combined analysis falls in the LMA region. Compared to the rates analysis the goodness of fit of the LOW(SMA) region increases(decreases). With the inclusion of the day-night dependence of the data the goodness of fit in the SMA region worsens further .
In this work we have not included the new GNO result which is consistent with the Gallex and SAGE data. Thus its inclusion is not expected to affect the conclusions drastically. For illustration we give below the results of the global analysis of rates and spectrum including the GNO data. We take the weighted average of Gallex and GNO and treat SAGE as a separate experiment. The best-fit values and $`\chi _{min}^2`$ that we get are:
* $`\mathrm{sin}^2\theta =5.26\times 10^4,\mathrm{\Delta }m^2=5.28\times 10^6`$ eV$`{}_{}{}^{2},X_B=0.61,\chi _{min}^2=12.73`$, g.o.f = 85.21% (SMA)
* $`\mathrm{sin}^2\theta `$ = 0.18, $`\mathrm{\Delta }m^2=2.48\times 10^5`$ eV<sup>2</sup>, $`X_B`$ = 1.39, $`\chi _{min}^2`$ = 11.55, g.o.f = 90.39% (LMA)
* $`\mathrm{sin}^2\theta `$ = 0.41, $`\mathrm{\Delta }m^2=9.39\times 10^8`$ eV<sup>2</sup>, $`X_B`$ = 0.89, $`\chi _{min}^2`$ = 19.85, g.o.f = 40.34% (LOW)
Thus the global best-fit continues to be in the LMA region.
In conclusion, we have probed the most recent solar neutrino data on total rates and the observed electron energy spectrum at SK from various angles within the framework of MSW flavour conversion. We find good fits in some situations but a degree of uncertainty still remains since different fits do not prefer the same values of the parameters. More data from the running and new experiments, it is hoped, will further sharpen the results in the near future.
Acknowledgements
D.M. and A.R. are partially supported by the Eastern Centre for Research in Astrophysics, India. A.R. also acknowledges a research grant from the Council of Scientific and Industrial Research, India. We would like to thank Sandhya Choubey for pointing out an error in one of our computer codes and J.W.F. Valle for drawing our attention to the updated analysis in . S.G. would like to thank Plamen Krastev for many helpful correspondences.
Figure Captions
Fig 1. The 99% and 90% C.L. allowed regions in the $`\mathrm{\Delta }m^2`$ \- $`\mathrm{sin}^2\theta `$ plane from the analysis of total rates for the Chlorine and Gallium detectors and the 1117-day data from SK. The best-fit points are also indicated. The dark side ($`\mathrm{sin}^2\theta >0.5`$) is indicated by the dashed line.
Fig 2. The 90% C.L. allowed region in the $`\mathrm{\Delta }m^2`$ \- $`\mathrm{sin}^2\theta `$ plane from the 1117-day SK recoil electron spectrum data. The regions enclosed by the contours are disallowed. The best-fit point is indicated. The dark side is to the right of the dashed line.
Fig. 3. The minimum $`\chi ^2`$ for fits to the SK 1117-day recoil electron spectrum as a function of $`\mathrm{sin}^2\theta `$ ($`\mathrm{\Delta }m^2`$) are shown in (a), (b), and (c) ((d), (e), and (f)) when the parameter ranges are determined by the 99% C.L. allowed regions in the SMA, LMA, and LOW fits respectively to the total rates data. The solid (broken) curves are obtained when $`X_B`$ is held fixed at its SSM value (allowed to vary). The dash-dotted line indicates the 90% C.L. limit for 3 parameters. See text for more details.
Fig. 4. The 99% and 90% C.L. allowed region in the $`\mathrm{\Delta }m^2`$ \- $`\mathrm{sin}^2\theta `$ plane from an analysis of the total rates from the Chlorine and Gallium detectors and the 1117-day SK data taken together with the 1117-day SK recoil electron spectrum. The normalization of the <sup>8</sup>B flux is chosen as a free parameter in (a) and held fixed at the SSM value in (b). In (b) the correlations between the rates and spectrum data are included. The best-fit points are also indicated. The dark side corresponds to the right of the dashed line.
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# 1 Introduction
## 1 Introduction
Our understanding of the non–perturbative sector of field and string theories has greatly progressed in recent times. In , for the first time, the entire non–perturbative contribution to the holomorphic part of the Wilsonian effective action was computed for $`N=2`$ globally supersymmetric (SUSY) theories with gauge group $`SU(2)`$, using ansätze dictated by physical intuitions. A few years later, a better understanding of non–perturbative configurations in string theory led to the conjecture that certain IIB string theory correlators on an $`AdS_5\times S^5`$ background are related to Green’s functions of composite operators of an $`N=4`$ $`SU(N_c)`$ Super Yang–Mills (SYM) theory in four dimensions in the large $`N_c`$ limit . Although supported by many arguments, these remarkable results remain conjectures and a clear mathematical proof seems to be out of reach at the moment. In our opinion this state of affairs is mainly due to the lack of adequate computational tools in the non–perturbative region. To the extent of our knowledge, the only way to perform computations in this regime in SUSY theories and from first principles is via multi–instanton calculus. Using this tool, many partial checks have been performed on these conjectures, both in $`N=2`$ and $`N=4`$ SUSY gauge theories . The limits on these computations come from the exploding amount of algebraic manipulations to be performed and from the lack of an explicit parametrization of instantons of winding number greater than two . In order to develop new computational tools that might allow an extension to arbitrary winding number, we revisit instanton computations for $`N=2`$ in the light of the topological theory built out of $`N=2`$ SYM, i.e. the so–called Topological Yang–Mills theory (TYM) .<sup>1</sup><sup>1</sup>1This is somewhat different from previous work in which the spectral curves which describe the moduli space of vacua for $`N=2`$ theories with various gauge groups were put in relation with integrable systems which, in turn, are related to 2–dimensional topological field theories. The study of the relationship between this approach and the one we present here goes beyond the scope of this paper.
That the TYM might play an important rôle in instanton computations became apparent with the results of . The agreement of these computations with the results of pointed out that instantons saturate all the non–perturbative contributions to the $`N=2`$ SYM low–energy effective action, and that the saddle point expansion around the classical solution does not receive any perturbative correction. This situation seems to be related to a sort of localization theorem . In this respect $`N=2`$ stands as an isolated case. Its effective action can be separated into a holomorphic part (which encodes the geometry of the quantum moduli space of vacua) and a non–holomorphic one. In a powerful non–renormalization theorem for the former was found. Subsequently, the quantum holomorphic piece was exactly determined , and its non–perturbative part checked against instanton calculations . It is this kind of contributions that we claim can also be computed from a closely related topological field theory. The story for the non–holomorphic terms is completely different. Indeed, the leading instanton contribution to the higher–derivative terms in the effective action gets perturbative corrections , in a way similar to what occurs in the $`N=4`$ SUSY theory, in which the non–perturbative contributions to all the relevant correlators get also perturbatively corrected, as a consequence of the vanishing of the $`R`$–symmetry anomaly.
From now on we will focus on the $`N=2`$ case. First of all, we generalize the standard approach to TYM to encompass the case in which the scalar field acquires a non–vanishing vacuum expectation value (v.e.v.), and define a BRST operator on field space, which on the moduli space $`^+`$ ($`^{}`$) of (anti–)self–dual gauge connections acts as the exterior derivative. This allows us to show that the computation of the relevant Green’s functions boils down to integrating differential forms on $`^+`$ ($`^{}`$). More precisely, after integrating out the quantum fluctuations one is left with a theory living on the instanton moduli space. To describe this space we will make use of the ADHM construction .
The TYM framework allows for a better understanding of the geometry underlying the computation of correlators of observables, and casts also new light on old problems concerning instanton calculus. On one hand we will learn that the BRST operator built with a zero v.e.v. for the scalar and that obtained with a non–zero v.e.v. cannot be smoothly deformed one into the other. Thus, it does not make sense trying to match computations performed in these two different regimes. On the other hand, in the case of non–zero scalar v.e.v., the twisted formulation naturally leads to the construction of the so–called “constrained instanton” field configurations thus giving a firmer basis to this approach. The Ward identities associated to the scalar supersymmetry transformations show that the constrained instanton computational method actually gives the correct result for the Green’s functions of physical observables. This in turn implies that, as argued before, instantons saturate the non–perturbative contribution to the relevant Green’s functions, and shows how the non–renormalization theorem of explicitly works in the context of instanton calculus.
It is worth remarking that in the geometrical approach outlined here, the instanton measure for the $`N=2`$ SYM theory (i.e. the integration measure over the moduli, or collective coordinates) emerges in a very natural way, without resorting to any intricate zero–mode calculation as in previous approaches. In the same vein, we derive with purely algebraic methods an explicit realization of the BRST algebra on instanton moduli space. This derivation is deeply related to the construction of this space as a hyperkähler quotient , which we study in the last section.
In the case of non–zero scalar v.e.v., the instanton action, i.e. the $`N=2`$ SYM action functional computed on the zero–modes (in our picture these are the field configurations onto which the action functional of TYM projects) can be interpreted as the commutator of the BRST charge with an appropriate function. This leads to the possibility of writing the correlators of physical observables as integrals of total differentials on $`^+`$. The circumstance that these Green’s functions can be computed in principle on the boundary of $`^+`$ may greatly help in computations, since instantons on $`^+`$ obey a kind of dilute gas approximation, as we will explain in subsec. 5.3. This might lead to recursion relations of the type found in , and simplify instanton calculations. We also explore how the geometrical approach described here works in the case of vanishing v.e.v., and apply these ideas to the computation of correlators in Witten’s topological field theory.
To avoid misunderstandings, we stress that the fact that certain correlators of $`N=2`$ SYM can be calculated using the formalism of the TYM does not mean that the former is a topological theory: $`N=2`$ SYM is a “physical” theory with its own degrees of freedom and a running coupling constant. In fact, TYM is formally derived from $`N=2`$ SYM by the twisting procedure which, in flat space, turns out to be just a variable redefinition. However, promoting the scalar supersymmetry generator present in the twisted $`N=2`$ algebra to a BRST charge implies great changes in the physical interpretation of the theory; some fields become ghosts and their engineering dimensions change . TYM theory deserves its name topological because it is a theory with zero degrees of freedom, whose correlators can be related to topological invariants of the four dimensional manifold on which the theory lives. Also in SUSY gauge theories a class of position–independent correlators exists . One realizes that a subset of correlators of $`N=2`$ SYM coincides with a subset of the observables defined in TYM over $`\text{IR}^4`$: as a consequence, these Green’s functions can be computed in either theory, according to one’s preferences.
Summarizing, we believe this approach provides us with a natural and simplifying framework for investigating the non–perturbative dynamics of SUSY gauge theories. An abbreviated account of part of the results described here was presented in .
This paper is organized as follows. In subsec. 2.1 we recall some basics of topological field theories with vanishing v.e.v. for the scalar field. We derive the set of identities which define a BRST operator and show that the functional integration projects the fields onto the subspace of the zero–modes of the relevant kinetic operators in the instanton background. In subsec. 2.2 we generalize this discussion to the case of a non–vanishing v.e.v. and clarify the relationship between our approach and the constrained instanton computational method. In subsec. 3.1, after an introduction to the ADHM construction of instantons, we write the solutions to the equations of motion derived from the TYM action. In subsec. 3.2 we use the results of the previous subsection and the identities associated to the BRST symmetry to find how the BRST charge acts on the relevant quantities defined in the ADHM construction (e.g. the instanton moduli) in the absence of a v.e.v. for the scalar field, and in subsec. 3.3. in the case of a non–vanishing v.e.v. We present in subsec. 3.4 a purely algebraic (and independent) derivation of the BRST algebra on instanton moduli space and of the solutions to the equations of motion (which were obtained in the previous subsections). In sec. 4 we discuss how to compute instanton–dominated Green’s functions using the formalism we have developed. It is important to understand how in our approach the instanton measure arises. This crucial issue is discussed in subsec. 4.1, where we also study in detail the cases of winding number one and two. In subsec. 4.2 we compute the multi–instanton action (which is non–zero when the scalar field acquires a non–vanishing v.e.v.) and show that it can be written as a BRST–exact quantity. Sec. 5 is devoted to the calculation of $`u=\text{Tr}\varphi ^2`$, the gauge invariant quantity which parametrizes the moduli space of quantum vacua of the $`N=2`$ SYM theory (from which one can obtain the Seiberg–Witten low–energy Wilsonian action using Matone’s relation ). First, we find a general expression in our framework for the $`k`$–instanton contribution to $`u`$. Then, on one hand, in subsec. 5.1 and 5.2, we compute $`\text{Tr}\varphi ^2`$ in the bulk of $`^+`$ for winding numbers $`k=1,2`$; on the other hand, in subsec. 5.3, using the observation of subsec. 4.2, we show that the contribution to $`u`$ can be written as a total derivative integrated on the moduli space of instantons. This suggests the interesting possibility of computing it directly on the boundary of $`^+`$; we explicitly check this in a $`k=1`$ computation, getting the correct result. In sec. 6 we consider the case of a vanishing v.e.v. (to which our formalism also applies), and compute $`\text{Tr}\varphi ^2\text{Tr}\varphi ^2`$ for winding number $`k=1`$ both in the bulk and on the boundary of the instanton moduli space. Finally, in sec. 7 we construct the metric on the $`8k`$–dimensional moduli space of self–dual gauge connections for winding number $`k=2`$ following the aforementioned hyperkähler quotient procedure.
## 2 Topological Yang–Mills Theory
It is well known that, if the generators of the rotation group of $`\text{IR}^4`$ are redefined in a suitably twisted fashion, the $`N=2`$ SYM theory gives rise to the TYM theory considered in . A key feature of the twisted theory is the presence of a scalar fermionic symmetry $`Q`$, which is still an invariance of the theory when this is formulated on a generic (differentiable) four–manifold $`M`$. This scalar symmetry will play a major rôle, for the following reasons. First, the correlation functions of the physical observables of the theory (the cohomology classes of $`Q`$) are independent of the metric on $`M`$ by virtue of the Ward identities associated to $`Q`$ . This also implies that these functions must be independent of the positions of the operatorial insertions.<sup>2</sup><sup>2</sup>2That in some SUSY gauge theories there exists a class of position–independent correlators was observed in . Moreover, the same Ward identities entail that certain Green’s functions can be computed exactly in the semiclassical limit; this is why instantons come into play. Finally, when one modifies the scalar supersymmetry charge $`Q`$ to make it nilpotent, the resulting (BRST) operator acts as the exterior derivative on the anti–instanton moduli space $`^{}`$ . As we will see, functional integration reduces to integrating differential forms on $`^{}`$ (this is what we call the localization procedure). We will later show that the Green’s functions of the observables can be written as integrals of total derivatives on $`^{}`$.
Let us first recall how the twisting operation works . The global symmetry group of the $`N=2`$ SUSY theory in flat space is
$$SU(2)_L\times SU(2)_R\times SU(2)_A\times U(1)_R,$$
(2.1)
where the first two factors represent the Euclidean Lorentz group (i.e. the rotation group of $`\text{IR}^4`$), while $`SU(2)_A`$ is the automorphism group of the $`N=2`$ supersymmetry algebra and $`U(1)_R`$ is the usual $`R`$–symmetry. The twist consists in replacing one of the factors of the rotation group, say for definiteness $`SU(2)_R`$, with a diagonal subgroup $`SU(2)_R^{}`$ of $`SU(2)_R\times SU(2)_A`$. The symmetry group of the twisted theory is then
$$SU(2)_L\times SU(2)_R^{}\times U(1)_R.$$
(2.2)
With respect to the twisted group, the SUSY charges decompose as a scalar $`Q`$, a self–dual antisymmetric tensor $`Q_{\mu \nu }`$ and a vector $`Q_\mu `$:
$`\overline{Q}_{\dot{\alpha }}^{\dot{A}}QQ_{\mu \nu },`$
$`Q_\alpha ^{\dot{A}}Q_\mu .`$ (2.3)
In particular, the charge $`Q`$ belongs to the $`(0,0)^1`$ representation of (2.2), while the charges $`Q_\mu ,Q_{\mu \nu }`$ belong respectively to the $`(\frac{1}{2},\frac{1}{2})^1`$ and $`(0,1)^1`$ representations.<sup>3</sup><sup>3</sup>3With the upper index we denote the $`R`$–symmetry charge. In the twisted theory it is natural to redefine the fields of $`N=2`$ SYM as
$`A_\mu A_\mu ,`$
$`\overline{\lambda }_{\dot{\alpha }}^{\dot{A}}\eta \chi _{\mu \nu },`$
$`\lambda _\alpha ^{\dot{A}}\psi _\mu ,`$
$`\varphi \varphi .`$ (2.4)
The anticommuting fields $`\eta ,\chi _{\mu \nu },\psi _\mu `$ are respectively a scalar, a self–dual antisymmetric tensor and a vector, belonging to the $`(0,0)^1`$, $`(0,1)^1`$ and $`(\frac{1}{2},\frac{1}{2})^1`$ representations of the twisted group; the gauge field $`A_\mu `$ and the scalar field $`\varphi `$ belong respectively to the $`(\frac{1}{2},\frac{1}{2})^0`$ and $`(0,0)^2`$ representation.
In the following we will be mainly interested in the multiplet of fields $`(A_\mu ,\psi _\mu ,\varphi )`$, whose transformations under the action of $`Q`$ read
$`QA_\mu =\psi _\mu ,`$
$`Q\psi _\mu =D_\mu \varphi ,`$
$`Q\varphi =0.`$ (2.5)
These equations imply that $`Q`$ is nilpotent modulo gauge transformations with parameter $`\varphi `$, since
$`Q^2A_\mu =D_\mu \varphi ,`$
$`Q^2\psi _\mu =[\varphi ,\psi _\mu ],`$
$`Q^2\varphi =0`$ (2.6)
(this is analogous to what one comes across in studying the supersymmetry algebra in the Wess–Zumino gauge). (2.6) allows us to regard $`Q`$ as a BRST–like charge. To this end, let us assign $`Q`$ a ghost number equal to 1; this can be done by simply identifying its $`R`$–charge with the ghost number. Accordingly, the fields of the twisted $`N=2`$ vector multiplet (2.4) acquire a ghost number equal to their respective $`R`$–charge. Since the canonical dimension of the BRST operator is usually taken to be zero, we redefine the canonical dimension of $`Q`$ to zero. The resulting canonical dimensions and ghost numbers of the fields are summarized in the table below.
| Fields | $`A`$ | $`\psi `$ | $`\chi `$ | $`\eta `$ | $`\varphi `$ | $`\overline{\varphi }`$ |
| --- | --- | --- | --- | --- | --- | --- |
| dimension | 1 | 1 | 2 | 2 | 0 | 2 |
| ghost # | 0 | 1 | -1 | -1 | 2 | -2 |
At this point we need to distinguish the case in which the scalar v.e.v. vanishes from that in which it is non–vanishing. In the next subsection we will focus on the former situation; the latter requires a more detailed discussion, and will be studied separately in subsec. 2.2.
### 2.1 Case I: Zero Vacuum Expectation Value for the Scalar Field
As (2.6) shows, $`Q`$ is not nilpotent. A strictly nilpotent BRST charge can be obtained from $`Q`$ by introducing a generalized BRST operator $`s`$ including both the gauge symmetry and the scalar supersymmetry of the theory ,
$$s=s_g+Q.$$
(2.7)
$`s_g`$ is the usual BRST operator associated to the gauge symmetry,
$`s_gA=Dc,`$
$`s_g\psi =[c,\psi ],`$
$`s_g\varphi =[c,\varphi ],`$
$`s_gc={\displaystyle \frac{1}{2}}[c,c],`$ (2.8)
the ghost number and canonical dimension of the ghost field $`c`$ being respectively one and zero. The 1–form $`A=A_\mu dx^\mu `$ is the gauge connection, with curvature $`F=1/2F_{\mu \nu }dx^\mu dx^\nu =dA+AA`$; $`\psi =\psi _\mu dx^\mu `$ is an anticommuting 1–form, and $`D`$ is the covariant exterior derivative on the manifold $`M`$. The action of $`Q`$ on the ghost field $`c`$ is obtained by requiring that $`s^2=0`$, and turns out to be simply
$$Qc=\varphi .$$
(2.9)
(2.5), (2.8) and (2.9) thus lead to the following BRST identities :
$`sA=\psi Dc,`$
$`s\psi =[c,\psi ]D\varphi ,`$
$`s\varphi =[c,\varphi ],`$
$`sc={\displaystyle \frac{1}{2}}[c,c]+\varphi ,`$ (2.10)
The algebra (2.1) can be read as the definition and the Bianchi identities for the curvature
$$\widehat{F}=F+\psi +\varphi $$
(2.11)
of the connection
$$\widehat{A}=A+c$$
(2.12)
of the universal bundle $`P\times 𝒜/𝒢`$, where $`P,𝒜,𝒢`$ are respectively the principal bundle over $`M`$, the space of connections and the group of gauge transformations. The exterior derivative on the manifold $`M\times 𝒜/𝒢`$ is given by
$$\widehat{d}=d+s.$$
(2.13)
Notice that from the last of (2.1) we learn that the scalar field $`\varphi `$ can be seen as the curvature of the connection $`c`$.
We now come to define the observables of the TYM theory; these are given by the elements of the equivariant cohomology of $`s`$ , which satisfy the descent equations
$`s{\displaystyle \frac{1}{2}}\text{Tr}F^2=d\text{Tr}F\psi ,`$
$`s\text{Tr}F\psi =d\text{Tr}\left(\varphi F+{\displaystyle \frac{1}{2}}\psi ^2\right),`$
$`s\text{Tr}\left(\varphi F+{\displaystyle \frac{1}{2}}\psi ^2\right)=d\text{Tr}\varphi \psi ,`$
$`s\text{Tr}\varphi \psi ={\displaystyle \frac{1}{2}}d\text{Tr}\varphi ^2,`$
$`s{\displaystyle \frac{1}{2}}\text{Tr}\varphi ^2=0.`$ (2.14)
(2.14) allows one to build local functions of the fields which are BRST invariant modulo $`d`$–exact terms; the simplest example of a physical observable is the gauge invariant polynomial $`\text{Tr}\varphi ^2`$, as the last of (2.14) shows. We will see in sec. 3 that this geometrical approach provides us with an operative tool which allows us to compute the Green’s functions of observables starting only from the knowledge of the universal connection (2.12), in particular without solving any equation of motion.
As shown in , a TYM action can be interpreted as a pure gauge–fixing term,
$$S_{\mathrm{TYM}}=2d^4xs\text{Tr}\mathrm{\Psi },$$
(2.15)
where the gauge–fixing fermion is chosen to be
$$\mathrm{\Psi }=\chi ^{\mu \nu }F_{\mu \nu }^+D^\mu \overline{\varphi }\psi _\mu +\overline{c}^\mu A_\mu ,$$
(2.16)
and
$$F_{\mu \nu }^+=\frac{1}{2}\left(F_{\mu \nu }+\frac{1}{2}ϵ_{\mu \nu \rho \sigma }F^{\rho \sigma }\right)$$
is the self–dual component of the field strength $`F_{\mu \nu }`$. The anti–fields $`\chi _{\mu \nu }`$, $`\overline{\varphi }`$ and $`\overline{c}`$ transform under the BRST symmetry as
$`s\chi _{\mu \nu }=B_{\mu \nu },`$
$`s\overline{\varphi }=\eta ,`$
$`s\overline{c}=b,`$ (2.17)
while the Lagrange multipliers $`B_{\mu \nu }`$, $`\eta `$ and $`b`$ as
$`sB_{\mu \nu }=0,`$
$`s\eta =0,`$
$`sb=0.`$ (2.18)
The anti–field $`\overline{c}`$ has ghost number $`1`$ and dimension $`2`$, while the Lagrange multipliers ($`B_{\mu \nu }`$, $`b`$) have ghost number $`0`$ and dimension $`2`$. Acting with the BRST operator $`s`$ in (2.15), we obtain the following explicit form for the TYM action
$`S_{\mathrm{TYM}}`$ $`=`$ $`2{\displaystyle }d^4x\text{Tr}[B^{\mu \nu }F_{\mu \nu }^+\chi ^{\mu \nu }(D_{[\mu }\psi _{\nu ]})^++\eta D^\mu \psi _\mu +`$ (2.19)
$`\overline{\varphi }(D^2\varphi [\psi ^\mu ,\psi _\mu ])+b^\mu A_\mu +`$
$`+\chi ^{\mu \nu }[c,F_{\mu \nu }^+]\overline{\varphi }[c,D^\mu \psi _\mu ]\overline{c}s(^\mu A_\mu )],`$
where
$$(D_{[\mu }\psi _{\nu ]})^+=\frac{1}{4}\left(D_\mu \psi _\nu D_\nu \psi _\mu +ϵ_{\mu \nu \rho \sigma }D^\rho \psi ^\sigma \right)$$
(2.20)
is the self–dual component of the tensor $`D_{[\mu }\psi _{\nu ]}`$. (2.19) is obtained integrating by parts the term in $`\overline{\varphi }`$ of (2.16); the corresponding surface term vanishes, since in this subsection we limit ourselves to study the case in which all the fields have trivial boundary conditions.
The main property of the action (2.19) is that it localizes the fields in the algebra (2.1) onto certain sections of the universal bundle. In particular, functional integration over the fields $`B^{\mu \nu }`$ and ($`\chi ^{\mu \nu }`$, $`\eta `$) in the first line of (2.19) leads respectively to
$$F_{\mu \nu }^+=0,$$
(2.21)
which implies that $`A`$ is an anti–self–dual gauge connection, and
$`(D_{[\mu }\psi _{\nu ]})^+=0,`$
$`D^\mu \psi _\mu =0,`$ (2.22)
which entails that $`\psi `$ is an element of the tangent bundle $`T_A^{}`$. In turn, functional integration over the anti–field $`\overline{\varphi }`$ leads to the equation
$$D^2\varphi =[\psi ^\mu ,\psi _\mu ]$$
(2.23)
for the field $`\varphi `$, while integration on the Lagrange multiplier $`b`$ imposes the usual transversality condition
$$^\mu A_\mu =0.$$
(2.24)
By plugging the expression for $`\psi `$ deduced from the first equation in (2.1) into the transversality condition $`D^\mu \psi _\mu =0`$, we get
$$D^2c=D^\mu (sA_\mu ),$$
(2.25)
which determines the ghost field $`c`$. Two observations are in order. First, let us remark that the first equation in (2.1) is an old acquaintance . It is in fact very well known that differentiating the gauge connection with respect to collective coordinates ($`sA`$) fails to give a transverse zero–mode ($`\psi `$). To ensure the correct gauge condition, the addition of a gauge transformation ($`Dc`$) is needed: (2.25) just determines this gauge transformation. A suitable framework for multi–instanton calculations is given by the ADHM construction. It is interesting to discover that this construction, together with the TYM formalism, allows one to write the universal connection (2.12) (and consequently the ghost $`c`$ and the gauge transformation $`Dc`$) in a very natural and straightforward way . We will focus on this aspect in sec. 3. Finally, the terms in the last line of (2.19) vanish due to the conditions (2.21), (2.22) and (2.24).
Summarizing, we have seen that after functional integration on the anti–fields and the Lagrange multipliers, we are left with an integration on the space of anti–self–dual gauge connections $`^{}`$ and its tangent bundle $`T_A^{}`$, with a functional measure equal to 1, since the action $`S_{\mathrm{TYM}}`$ vanishes on the field subspace identified by the (zero–mode) equations (2.21)–(2.25).
Notice that in this approach the functional integral is performed exactly, since the gauge–fixing fermion (2.16) is linear in the antifields, and there are no perturbative corrections. It is important to remark that the action obtained by twisting the $`N=2`$ SYM theory (i.e. the action of Witten’s topological field theory ) actually differs from (2.19) by some extra terms, which spoil the linearity of $`S_{\mathrm{TYM}}`$ in the anti–fields. However, as we will show below, they are BRST–exact terms corresponding to a continuous deformation of the gauge–fixing
$$S_{N=2}=S_{\mathrm{TYM}}+s𝒱;$$
(2.26)
the v.e.v. of an $`s`$–closed operator $`𝒪`$, i.e. such that $`s𝒪=0`$, is controlled by the Ward identity (in the following $`[\delta \phi ]`$ is shorthand for the integration measure)
$`<𝒪>_{S+s𝒱}`$ $``$ $`{\displaystyle [\delta \phi ]e^{(S+s𝒱)}𝒪}={\displaystyle [\delta \phi ]e^S𝒪(1s𝒱+\mathrm{})}`$ (2.27)
$`=`$ $`<𝒪>_S<𝒪s𝒱>_S+\mathrm{}=<𝒪>_S<s(𝒪𝒱)>_S+\mathrm{}=`$
$`=`$ $`<𝒪>_S,`$
the last equality following from the fact that the v.e.v. of an $`s`$–exact operator $`𝒫=s𝒬`$ vanishes if $`𝒬`$ is globally defined.<sup>4</sup><sup>4</sup>4In this respect we remark that $`\text{Tr}\varphi ^2`$ is not globally defined. This fact plays an important rôle in breaking $`N=2`$ SUSY into $`N=1`$ . This means that the action (2.19) and the twisted $`N=2`$ SYM action are completely equivalent, in the sense that the Green’s functions of $`s`$–closed operators can be computed using any one of them obtaining the same result. We now show that the Lagrangian obtained by twisting the $`N=2`$ SYM theory can be derived by modifying the gauge–fixing fermion (2.16), thus proving that the twisted version of $`N=2`$ SYM action introduced in and the TYM action (2.19) differ only by BRST–exact terms . To this end, let us consider the modified gauge–fixing fermion
$$\mathrm{\Psi }^{(\alpha )}=\chi ^{\mu \nu }\left(F_{\mu \nu }^+\frac{\alpha }{2}B_{\mu \nu }\right)D^\mu \overline{\varphi }\psi _\mu +\overline{c}^\mu A_\mu ,$$
(2.28)
where $`\alpha `$ is a gauge–fixing parameter; upon functional integration over the Lagrange multiplier $`B_{\mu \nu }`$, one gets
$`S_{\mathrm{TYM}}^{(\alpha )}`$ $`=`$ $`2{\displaystyle }d^4x\text{Tr}[{\displaystyle \frac{1}{2\alpha }}F^{+\mu \nu }F_{\mu \nu }^+\chi ^{\mu \nu }(D_{[\mu }\psi _{\nu ]})^++\eta D^\mu \psi _\mu +`$ (2.29)
$`\overline{\varphi }(D^2\varphi [\psi ^\mu ,\psi _\mu ])+b^\mu A_\mu +`$
$`+\chi ^{\mu \nu }[c,F_{\mu \nu }^+]\overline{\varphi }[c,D^\mu \psi _\mu ]\overline{c}s(^\mu A_\mu )].`$
This functional leads to the same equations of motion of the $`N=2`$ SYM theory.
Three observations are in order. First, notice that the partition function (or more generally the Green’s functions) defined by (2.29) has to be studied using the usual saddle point techniques due to the presence of the kinetic term for the gauge field. As discussed in , supersymmetry ensures the cancellation of bosonic and fermionic determinants arising when integrating out the non–zero modes and the resulting integration is over $`^{}`$, which is the same result obtained from the use of the action (2.19). Second, the renormalization group invariant scale (which multiplies the Green’s functions computed in the conventional supersymmetric theory) does not appear. In fact the divergent term of the one–loop effective action is proportional to $`(F^+)^2`$ which is obviously vanishing on $`^{}`$ .
Last, and most importantly, let us observe that the equivalence between $`S_{\mathrm{TYM}}^{(\alpha )}`$, Eq. (2.29), and $`S_{\mathrm{TYM}}`$, Eq. (2.19), in the computation of correlators of observables is not surprising, since they cannot depend on the choice of the gauge parameter $`\alpha `$; therefore nothing prevents us from choosing directly $`\alpha =0`$. As we will discuss in sec. 2.2, in the presence of a non–trivial v.e.v. for the scalar field, the equivalence of (2.19) to the $`N=2`$ SYM action (in the sense previously specified) makes the configurations of the constrained instanton method emerge from functional integration, without any approximation procedure.
Let us now sketch the geometrical interpretation of the instanton calculus suggested by the topological formulation. To begin with, the first equation in (2.1) together with (2.22) imply that, as announced, the BRST operator $`s`$ has an intriguing explicit realization on the moduli space as the exterior derivative . Furthermore, once the universal gauge connection (2.12) is given, the other field configurations $`\psi ,\varphi `$ are in turn immediately determined as components of the universal curvature (2.11), as it will be worked out in detail in sec. 3. Topological correlators are then built up as differential forms on the moduli space, where the form degree of the fields equals their ghost number. For example, for winding number $`k=1`$ the top form on the (8–dimensional) instanton moduli space is given by $`\text{Tr}\varphi ^2(x_1)\text{Tr}\varphi ^2(x_2)`$. We will explicitly compute the corresponding Green’s function in sec. 6, both with a bulk calculation and with a calculation on the boundary of the instanton moduli space, finding the same result.
### 2.2 Case II: Non–Zero Vacuum Expectation Value for the Scalar Field
In this subsection we will extend the construction of the TYM to encompass the presence of a non–vanishing v.e.v. for the scalar field. To this end, observe first that a non–zero v.e.v. for $`\varphi `$,
$$\underset{|x|\mathrm{}}{lim}\varphi =v\frac{\sigma _3}{2i},$$
(2.30)
implies the existence of a (non–zero) central charge $`Z`$ in the SUSY algebra. Then the operator defined in (2.7) is no longer nilpotent; instead, it closes on a $`U(1)`$ central charge transformation
$`(s_g+Q)^2A=ZAD\varphi _Z,`$
$`(s_g+Q)^2\psi =Z\psi [\varphi _Z,\psi ],`$
$`(s_g+Q)^2\varphi _Z=Z\varphi _Z0,`$ (2.31)
where the scalar field $`\varphi _Z`$ plays the rôle of a gauge parameter and satisfies the equation
$`D^2\varphi _Z=0,`$
$`\underset{|x|\mathrm{}}{lim}\varphi _Z=v{\displaystyle \frac{\sigma _3}{2i}}.`$ (2.32)
Notice that from (2.31) it follows that the central charge $`Z`$ has ghost number two and canonical dimension zero. We also remind that the scalar charge $`Q`$ commutes with $`s_g`$, while $`Z`$ commutes with all the other charges by definition.
From (2.31) it follows that, in order to ensure the nilpotency property, the operator (2.7) has to be properly extended by including the central charge. We then define an extended BRST operator as
$$s=s_g+Q\lambda Z+\frac{}{\lambda },$$
(2.33)
where $`\lambda `$ is a fermionic parameter with ghost number $`1`$ and canonical dimension zero, such that $`\lambda Z`$ has the usual quantum numbers of a BRST operator. It is easy to see that the last term of (2.33) is needed to ensure the nilpotency of $`s`$; in fact, by using (2.31) and the property $`\lambda ^2=0`$, we get
$$s^2=(s_g+Q)^2\frac{}{\lambda }(\lambda Z)=0.$$
(2.34)
If we define the ghost field
$$\mathrm{\Lambda }\lambda \varphi _Z,$$
(2.35)
with the transformation
$$s\mathrm{\Lambda }=\varphi _Z[c,\mathrm{\Lambda }],$$
(2.36)
we can finally write the resulting BRST algebra as
$`sA=\psi D(c+\mathrm{\Lambda }),`$
$`s\psi =[c+\mathrm{\Lambda },\psi ]D\varphi ,`$
$`s\varphi =[c+\mathrm{\Lambda },\varphi ],`$
$`s(c+\mathrm{\Lambda })={\displaystyle \frac{1}{2}}[c+\mathrm{\Lambda },c+\mathrm{\Lambda }]+\varphi .`$ (2.37)
Notice that now the curvature $`\varphi `$ of the ghost $`c+\mathrm{\Lambda }`$ is decomposed as the sum
$$\varphi =\varphi _Z+\varphi _g,$$
(2.38)
where $`\varphi _g`$ is related to the usual gauge transformations
$`sc=\varphi _g{\displaystyle \frac{1}{2}}[c,c],`$
$`\underset{|x|\mathrm{}}{lim}\varphi _g=0.`$ (2.39)
(2.35) entails that the field $`\mathrm{\Lambda }`$, which has ghost number one and dimension zero, can be seen as the ghost related to the central charge symmetry. This fact can be analyzed in more detail by considering the equations (2.22) for the $`\psi `$ field: in the presence of a non–zero central charge, they no longer have a unique solution, since any field configuration of the form $`\psi D\stackrel{~}{\mathrm{\Lambda }}`$ satisfies the same equations on the anti–instanton background $`F_{\mu \nu }^+=0`$
$`(D_{[\mu }\psi _{\nu ]})^+[F_{\mu \nu }^+,\stackrel{~}{\mathrm{\Lambda }}]=(D_{[\mu }\psi _{\nu ]})^+=0,`$
$`D^\mu \psi _\mu D^2\stackrel{~}{\mathrm{\Lambda }}=D^\mu \psi _\mu =0,`$ (2.40)
provided that
$`D^2\stackrel{~}{\mathrm{\Lambda }}=0,`$
$`\underset{|x|\mathrm{}}{lim}\stackrel{~}{\mathrm{\Lambda }}0;`$ (2.41)
according to (2.32), this identifies $`\stackrel{~}{\mathrm{\Lambda }}`$ as the parameter of a central charge transformation. This degeneracy will be removed once the boundary conditions for the ghost $`\stackrel{~}{\mathrm{\Lambda }}`$ are fixed. Notice that if the central charge were zero, the equation $`D^2\stackrel{~}{\mathrm{\Lambda }}=0`$ would have trivial boundary conditions; its solution would be $`\stackrel{~}{\mathrm{\Lambda }}=0`$, and the degeneracy would disappear, as expected. In our case instead, from (2.32) and (2.35) it follows that $`\mathrm{\Lambda }`$ is a solution of (2.41) satisfying the boundary condition
$$\underset{|x|\mathrm{}}{lim}\mathrm{\Lambda }=\lambda v\frac{\sigma _3}{2i}$$
(2.42)
induced by the asymptotic behavior of the scalar field $`\varphi `$. The new BRST algebra could then be derived from (2.1) by just shifting the ghost $`c`$ to $`c+\mathrm{\Lambda }`$; $`c`$ is related to the usual gauge transformations, whereas $`\mathrm{\Lambda }`$ takes into account the new $`U(1)`$ transformations generated by the central charge.
Let us now evaluate the action (2.15). This can be done by noticing that, according to the algebra (2.2), its explicit form can be obtained simply by substituting the ghost $`c`$ in (2.19) with its shifted version $`c+\mathrm{\Lambda }`$. Unlike the case studied in the previous subsection, the action gets now a contribution from integrating by parts the term $`s\text{Tr}(D^\mu \overline{\varphi }\psi _\mu )`$ which does not vanish due to the non–trivial boundary conditions (2.30). Explicitly
$$sd^4x2\text{Tr}[(D^\mu \overline{\varphi })\psi _\mu ]=sd^4x2^\mu \text{Tr}(\overline{\varphi }\psi _\mu )sd^4x2\text{Tr}(\overline{\varphi }D^\mu \psi _\mu ).$$
(2.43)
We then get
$`S_{\mathrm{TYM}}`$ $`=`$ $`2{\displaystyle }d^4x\text{Tr}[B^{\mu \nu }F_{\mu \nu }^+\chi ^{\mu \nu }(D_{[\mu }\psi _{\nu ]})^++\eta D^\mu \psi _\mu +`$ (2.44)
$`\overline{\varphi }(D^2\varphi [\psi ^\mu ,\psi _\mu ])+b^\mu A_\mu +`$
$`+\chi ^{\mu \nu }[c+\mathrm{\Lambda },F_{\mu \nu }^+]\overline{\varphi }[c+\mathrm{\Lambda },D^\mu \psi _\mu ]\overline{c}s(^\mu A_\mu )]+`$
$`2s{\displaystyle d^4x2\text{Tr}(\overline{\varphi }D^\mu \psi _\mu )}+2{\displaystyle d^4x^\mu s\text{Tr}(\overline{\varphi }\psi _\mu )}.`$
The functional integration over the anti–fields and the Lagrange multipliers goes as in subsec. 2.1, and leads to the same set of (zero–mode) equations
$`F_{\mu \nu }^+=0,`$ (2.45)
$`(D_{[\mu }\psi _{\nu ]})^+=0,`$ (2.46)
$`D^\mu \psi _\mu =0,`$ (2.47)
$`D^2\varphi =[\psi ^\mu ,\psi _\mu ];`$ (2.48)
the key difference with respect to sec. 2.2 is that now the scalar field $`\varphi `$ has non–trivial boundary conditions as per (2.30). It is worth remarking that the (zero–mode subspace) configurations dictated by (2.45)–(2.48), together with the boundary condition (2.30), are exactly those which are exploited in the context of the constrained instanton method as approximate solutions to the saddle point equations. Instead, as we have explained, in our approach such field configurations naturally come into play after functional integrating the anti–fields and the Lagrange multipliers; no approximation is involved. The Ward identity (2.27) of the previous subsection applies also here, implying the equivalence of the action (2.44) to that of the $`N=2`$ SYM theory in computing the Green’s functions of the physical observables. This explains why the constrained instanton method gives the correct result for the calculation of these correlators.
Once the connection $`\widehat{A}=A+c+\mathrm{\Lambda }`$ is known, the BRST identities (2.2) provide us with the field configurations of $`A`$, $`\psi `$, $`\varphi `$ which solve the equations of motion (2.45)–(2.48). This possibility, and the circumstance that the BRST operator acts on instanton moduli space as the exterior derivative conspire to make it possible to explicitly work out (in the ADHM formalism) the aforementioned field configurations without solving their equations of motion. In the next section we will first guess the ADHM expression for $`\widehat{A}`$, and then show how the procedure outlined here works.
As in the zero v.e.v. case, functional integration is performed exactly, and we are left with an integration over $`^{}`$ and $`T_A^{}`$. No perturbative renormalization of the physical correlators calculated with TYM is allowed, in agreement with the non–renormalization theorem of . Green’s functions are built up as differential forms on the moduli space also in this case. However, the functional measure is now crucially different from 1, since the action computed on the zero–mode subspace gets a non–vanishing contribution $`S_{\mathrm{inst}}`$ from the last term of (2.44), which reads
$$S_{\mathrm{inst}}=2d^4x^\mu s\text{Tr}(\overline{\varphi }\psi _\mu ).$$
(2.49)
This in turn implies that $`\mathrm{exp}(S_{\mathrm{inst}})`$ acts as a generating functional for differential forms on the moduli space. This gives rise to non–trivial correlation functions which take contribution from topological sectors of any winding number $`k`$. The most interesting example is the v.e.v. of the gauge invariant, $`s`$–exact operator $`\text{Tr}\varphi ^2`$, i.e. $`u(v)=\text{Tr}\varphi ^2`$, which plays a prominent rôle in the context of the Seiberg–Witten model. In sec. 5 we will focus on this particular Green’s function. First, in (5.9) we will give the general expression for the contribution to $`u(v)`$ coming from the topological sector of winding number $`k`$; furthermore in subsec. 5.1 and 5.2 we will perform the computation for instanton number equal to one and two respectively. Last, in subsec. 5.3 we will illustrate the possibility of computing $`u(v)`$ with a calculation on the boundary of instanton moduli space. To support this idea we will work out the $`k=1`$ computation explicitly.
## 3 The BRST Algebra on Instanton Moduli Space
When restricted to configurations which obey the equations of motion dictated by the TYM action, the BRST algebra gets realized on instanton moduli space. In the following we will construct this realization explicitly. To this end we will start by briefly recalling some basic elements of the ADHM construction of instantons, which provides us with a parametrization of this moduli space. This description is given in terms of a redundant set of parameters; we will then focus on its reparametrization symmetries, which will play a major rôle in the following. Our first goal will be the construction of the BRST algebra on instanton moduli space starting from the knowledge of the solutions to (2.21), (2.22), (2.23), (2.25), for a generic winding number, which were found in . In our set–up we will also need a new ingredient, i.e. the solution to (2.25) for the ghost field $`c`$, which we will obtain ex novo.
However, a completely different path could be followed: indeed, we will show that it is possible to construct the algebra directly on instanton moduli space, in particular without solving any field equation. This is an important remark, since in this way the construction of the algebra acquires a geometrical meaning and stands on its own. This approach is further developed in sec. 7, where we show its close relationship with the hyperkähler quotient construction of the instanton moduli space.
### 3.1 Construction of the Solutions to the Equation of Motion in the ADHM Formalism
In sec. 2 we saw that the TYM action localizes the fields relevant to the BRST transformations in (2.1) onto a set of configurations dictated by a system of coupled differential equations. Here we will review how to construct explicit solutions to these equations of motion.
To begin with, recall that gauge fields $`A`$ are projected onto instanton configurations<sup>5</sup><sup>5</sup>5In the previous section we adopted the standard convention in topological field theories of taking the gauge curvature to be anti–self–dual. Unfortunately the literature on instanton calculus adopts the opposite convention (self–dual), to which we will conform from now on.. As it is well known, self–dual $`SU(2)`$ connections on $`S^4`$ can be put into one to one correspondence with holomorphic vector bundles of rank $`2`$ over $`\text{}\text{IP}^3`$ admitting a reduction of the structure group to its compact real form. The ADHM construction is an algorithm which gives all these holomorphic bundles and consequently all $`SU(2)`$ connections on $`S^4`$ (this $`S^4`$ should be thought of as the conformal compactification of $`\text{IR}^4`$. For the construction of instantons on $`\text{IR}^4`$, see for example ).
The construction is purely algebraic and we find it more convenient to use quaternionic notations. The points, $`x`$, of the quaternionic space $`\text{}\text{}^2\text{IR}^4`$ can be conveniently represented in the form $`x=x^\mu \sigma _\mu `$, with $`\sigma _\mu =(i\sigma _c,\text{1 l}_{2\times 2}),c=1,2,3.`$ The $`\sigma _c`$’s are the usual Pauli matrices, and $`\text{1 l}_{2\times 2}`$ is the 2–dimensional identity matrix. The conjugate of $`x`$ is $`x^{}=x^\mu \overline{\sigma }_\mu `$. A quaternion is said to be real if it is proportional to $`\text{1 l}_{2\times 2}`$ and imaginary if it has vanishing real part.
The prescription to find an instanton of winding number $`k`$ is the following: introduce a $`(k+1)\times k`$ quaternionic matrix linear in $`x`$
$$\mathrm{\Delta }=a+bx,$$
(3.1)
where $`a`$ has the generic form
$$a=\left(\begin{array}{ccc}w_1& \mathrm{}& w_k\\ & & \\ & a^{}& \end{array}\right);$$
(3.2)
$`a^{}`$ is a $`k\times k`$ quaternionic matrix. The (anti–hermitean) gauge connection is then written as
$$A=U^{}dU,$$
(3.3)
where $`U`$ is a $`(k+1)\times 1`$ matrix of quaternions providing an orthonormal frame of $`\text{Ker}\mathrm{\Delta }^{}`$, i.e.
$`\mathrm{\Delta }^{}U`$ $`=`$ $`0,`$ (3.4)
$`U^{}U`$ $`=`$ $`\text{1 l}_{2\times 2}.`$ (3.5)
The constraint (3.5) ensures that $`A`$ is an element of the Lie algebra of the $`SU(2)`$ gauge group. The condition of self–duality
$${}_{}{}^{}F=F$$
(3.6)
on the field strength of (3.3) is imposed by restricting the matrix $`\mathrm{\Delta }`$ to obey
$$\mathrm{\Delta }^{}\mathrm{\Delta }=(\mathrm{\Delta }^{}\mathrm{\Delta })^T,$$
(3.7)
where the superscript $`T`$ stands for transposition of the quaternionic elements of the matrix (without transposing the quaternions themselves). (3.7) in turn implies $`\mathrm{\Delta }^{}\mathrm{\Delta }=f^1\text{1 l}_{2\times 2}`$, where $`f`$ is an invertible hermitean $`k\times k`$ matrix (of real numbers). From (3.3), the field strength of the gauge field can be computed and it is
$$F=U^{}d\mathrm{\Delta }fd\mathrm{\Delta }^{}U.$$
(3.8)
From this one can derive the following remarkable expressions for $`\text{Tr}(FF)`$ (see also )<sup>6</sup><sup>6</sup>6The conventions used in this paper imply that the Pontryagin index is given by $`1/(8\pi ^2)\text{Tr}(FF)`$, and it is positive (negative) on (anti–)instanton configurations.:
$`\text{Tr}(FF)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{}\mathrm{}\mathrm{ln}\mathrm{det}\mathrm{\Delta }^{}\mathrm{\Delta }d^4x`$ (3.9)
$`=`$ $`d\text{Tr}\left[PdD(dD)^{}D(dD)^{}+{\displaystyle \frac{1}{3}}(D^{}dD)(D^{}dD)(D^{}dD)\right],`$ (3.10)
where
$$P=UU^{}=1\mathrm{\Delta }f\mathrm{\Delta }^{}$$
(3.11)
is the projector on the kernel of $`\mathrm{\Delta }^{}`$, and according to the columns of $`\mathrm{\Delta }`$, which are independent, have been orthonormalized and collected into a matrix we have called $`D`$.
Gauge transformations are implemented in this formalism as right multiplication of $`U`$ by a unitary (possibly $`x`$–dependent) quaternion. Moreover, $`A`$ is invariant under reparametrizations of the ADHM data as follows:
$$\mathrm{\Delta }Q\mathrm{\Delta }R,$$
(3.12)
with $`QSp(k+1),RGL(k,\text{})`$. It is straightforward to see that (3.12) preserves the bosonic constraint (3.7). These symmetries can be used to simplify the expressions of $`a`$ and $`b`$. Exploiting this fact, in the following we will choose the matrix $`b`$ to be
$$b=\left(\begin{array}{c}0_{1\times k}\\ \text{1 l}_{k\times k}\end{array}\right).$$
(3.13)
Choosing the canonical form (3.13) for $`b`$, the bosonic constraint (3.7) becomes
$`a^{}=a_{}^{}{}_{}{}^{T},`$ (3.14)
$`a^{}a=(a^{}a)^T.`$ (3.15)
Moreover, in this case there still exist left–over $`O(k)\times SU(2)`$ reparametrizations of the form (3.12), where now $`RO(k)`$,
$$Q=\left(\begin{array}{cccc}q& 0& \mathrm{}& 0\\ 0& & & \\ \mathrm{}& & R^T& \\ 0& & & \end{array}\right),$$
(3.16)
and $`qSU(2)`$. These transformations act non–trivially on the matrix $`a`$ and leave $`b`$ invariant. After imposing the constraint (3.7), the number of independent degrees of freedom contained in $`\mathrm{\Delta }`$ (that is the number of independent collective coordinates that the ADHM formalism uses to describe an instanton of winding number $`k`$) is $`8k+k(k1)/2`$; modding out the $`O(k)\times SU(2)`$ reparametrization transformations, we would get $`8k3`$ truly independent degrees of freedom. However (3.4) and (3.5) do not determine $`U_0/|U_0|`$, where $`U_0`$ is the first component of $`U`$; this adds three extra degrees of freedom, so that in conclusion we end up with a moduli space of dimension $`8k`$ (the instanton moduli space $`^+`$). It is easy to convince oneself that the arbitrariness in $`U_0/|U_0|`$ can be traded for the $`SU(2)`$ reparametrizations; in other words, one can forget to mod out the $`SU(2)`$ factor of the reparametrization group $`O(k)\times SU(2)`$ but fix the phase of the quaternion $`U_0`$ (setting for example $`U_0=|U_0|\text{1 l}_{2\times 2}`$). This is what we will actually do in the following.
We now focus our attention on the other fields involved in the BRST algebra (2.1). To begin with, the TYM action projects the anti–commuting 1–form $`\psi _\mu `$ onto the solutions to
$${}_{}{}^{}(D_{[\mu }\psi _{\nu ]})=D_{[\mu }\psi _{\nu ]},D_\mu \psi _\mu =0,$$
(3.17)
where $`D`$ is the covariant derivative in the instanton background, Eq.(3.3). The solution to (3.17) can be written as
$$\psi =U^{}f(d\mathrm{\Delta }^{})U+U^{}(d\mathrm{\Delta })f^{}U,$$
(3.18)
where $``$ is a $`(k+1)\times k`$ matrix of quaternions, whose elements are Grassmann variables; moreover, in order for (3.18) to satisfy (3.17), $``$ must obey the constraint
$$\mathrm{\Delta }^{}=(\mathrm{\Delta }^{})^T.$$
(3.19)
(3.17) tell us that the $`\psi `$ zero–modes are the tangent vectors to the instanton moduli space $`^+`$; as it is well known, the number of independent zero–modes is $`8k`$ (the dimension of $`^+`$), and we would like to see how this is implemented in the formalism of the ADHM construction. To this end, note that $``$ has $`k(k+1)`$ quaternionic elements ($`4k(k+1)`$ real degrees of freedom) which are subject to the $`4k(k1)`$ constraints given by (3.19). The number of independent $``$’s satisfying (3.19) is thus $`8k`$, as desired.
If we work in the gauge in which $`b`$ has the canonical form (3.13), then (3.19) can be conveniently elaborated as follows. We put $``$ in a form which parallels the one for $`a`$ in (3.2), i.e.
$$=\left(\begin{array}{ccc}\mu _1& \mathrm{}& \mu _k\\ & & \\ & ^{}& \end{array}\right),$$
(3.20)
$`^{}`$ being a $`k\times k`$ quaternionic matrix. Plugging (3.2), (3.13), (3.20) into (3.19) we get
$`^{}=_{}^{}{}_{}{}^{T},`$ (3.21)
$`a^{}=(a^{})^T.`$ (3.22)
When $`\mathrm{\Delta }`$ is transformed according to (3.12), the $``$’s must also be reparametrized in such a way to keep the constraint (3.19) unchanged. This implies that the $``$’s undergo the same formal reparametrization of $`\mathrm{\Delta }`$, that is
$$QR.$$
(3.23)
We now turn to the scalar field configuration. This is dictated by (2.23), which should be supplemented by some boundary condition at infinity. Without loss of generality, we will set
$$\underset{|x|\mathrm{}}{lim}\varphi =𝒜_{00}=v\sigma ^3/2i,$$
(3.24)
where $`v\text{}`$. The solution to (2.23) and (3.24) was found in and reads
$$\varphi =U^{}f^{}U+U^{}𝒜U,$$
(3.25)
where
$$𝒜=\left(\begin{array}{cccc}𝒜_{00}& 0& \mathrm{}& 0\\ 0& & & \\ \mathrm{}& & 𝒜^{}& \\ 0& & & \end{array}\right).$$
(3.26)
Here $`𝒜^{}`$ is a $`k\times k`$ real antisymmetric matrix, and the condition for (3.25) to satisfy (2.23) is
$$\mathrm{\Delta }^{}𝒜\mathrm{\Delta }(\mathrm{\Delta }^{}𝒜\mathrm{\Delta })^T=\mathrm{\Lambda }_f,$$
(3.27)
where we set
$$\mathrm{\Lambda }_f=^{}(^{})^T;$$
(3.28)
note that $`\mathrm{\Lambda }_f=\mathrm{\Lambda }_f^T\text{1 l}_{2\times 2}`$. Hereafter, by $`𝒜\mathrm{\Delta }`$ we intend the $`x`$–independent $`(k+1)\times k`$ matrix
$$𝒜\mathrm{\Delta }=\left(\begin{array}{ccc}𝒜_{00}w_1w_m𝒜_{m1}^{}& \mathrm{}& 𝒜_{00}w_kw_m𝒜_{mk}^{}\\ \\ & [𝒜^{},a^{}]& \end{array}\right).$$
(3.29)
(3.29) is obtained by exploiting the fact that, in the calculations of the observables, expressions like $`𝒜\mathrm{\Delta }`$ are always multiplied from the left by $`U^{}`$; therefore, recalling that $`U^{}a=U^{}bx`$, we can eliminate the $`x`$–dependence in $`𝒜\mathrm{\Delta }`$. Using (3.2) and (3.29), the $`x`$–dependence disappears also from the l.h.s. of (3.27), which becomes
$$\mathrm{\Delta }^{}𝒜\mathrm{\Delta }(\mathrm{\Delta }^{}𝒜\mathrm{\Delta })^TL𝒜^{}+\mathrm{\Lambda }_b(𝒜_{00});$$
(3.30)
according to , the action of $`L`$ on $`k\times k`$ matrices $`\mathrm{\Omega }^{}`$ is given by
$$L\mathrm{\Omega }^{}=\frac{1}{2}\{\mathrm{\Omega }^{},W\}+\frac{1}{2}\text{Tr}\left([\overline{a}^{},\mathrm{\Omega }^{}]a^{}\overline{a}^{}[a^{},\mathrm{\Omega }^{}]\right),$$
(3.31)
where $`W_{kl}=\overline{w}_kw_l+\overline{w}_lw_k`$, and
$$[\mathrm{\Lambda }_b]_{ij}(\mathrm{\Omega }_0)=\overline{w}_i\mathrm{\Omega }_0w_j\overline{w}_j\mathrm{\Omega }_0w_i.$$
(3.32)
Note that $`[\mathrm{\Lambda }_b]_{ij}(\mathrm{\Omega }_0)`$ are $`c`$–numbers when $`\mathrm{\Omega }_0^{}=\mathrm{\Omega }_0`$. (3.27) can be now more compactly written as
$$L𝒜^{}=\mathrm{\Lambda }_b(𝒜_{00})\mathrm{\Lambda }_f.$$
(3.33)
The structure of (3.33) suggests setting
$$𝒜^{}=𝒜_b^{}+𝒜_f^{},$$
(3.34)
where
$`L𝒜_b^{}`$ $`=`$ $`\mathrm{\Lambda }_b(𝒜_{00}),`$ (3.35)
$`L𝒜_f^{}`$ $`=`$ $`\mathrm{\Lambda }_f.`$ (3.36)
This decomposition is useful since the solution $`\varphi _{\mathrm{hom}}`$ to the homogeneous equation
$$D^2\varphi =0$$
(3.37)
with the non–trivial boundary condition (3.24) has the form
$$\varphi _{\mathrm{hom}}=U^{}𝒜_bU,$$
(3.38)
where we set
$$𝒜_b=\left(\begin{array}{cccc}𝒜_{00}& 0& \mathrm{}& 0\\ 0& & & \\ \mathrm{}& & 𝒜_b^{}& \\ 0& & & \end{array}\right);$$
(3.39)
using (3.30), (3.35) can be written as
$$\mathrm{\Delta }^{}𝒜_b\mathrm{\Delta }(\mathrm{\Delta }^{}𝒜_b\mathrm{\Delta })^T=0,$$
(3.40)
which is the homogeneous equation associated to (3.27). Moreover, $`L`$ is a generally invertible operator acting on $`k\times k`$ matrices . As a consequence, $`𝒜_b^{}0`$ if and only if $`𝒜_{00}0`$. This is because non–trivial solutions to the homogeneous equation (3.37) exist only when non–trivial boundary conditions on $`\varphi `$ are imposed. On the other hand, the solution $`\varphi _{\mathrm{inh}}`$ to (2.23) supplemented by trivial boundary conditions
$$\underset{|x|\mathrm{}}{lim}\varphi _{\mathrm{inh}}=0,$$
(3.41)
reads as
$$\varphi _{\mathrm{inh}}=U^{}f^{}U+U^{}𝒜_fU,$$
(3.42)
with $`𝒜_f`$ given by
$$𝒜_f=\left(\begin{array}{cccc}0& 0& \mathrm{}& 0\\ 0& & & \\ \mathrm{}& & 𝒜_f^{}& \\ 0& & & \end{array}\right).$$
(3.43)
As before, the reparametrization invariance (3.12) induces a transformation on the matrix $`𝒜`$, which can be found by requiring that the new matrix still satisfies (3.27) when $`\mathrm{\Delta }`$ is replaced by its transformed expression; to this end one must have
$$𝒜Q𝒜Q^{}.$$
(3.44)
The last field relevant to our discussion is the ghost field $`c`$, which in principle should be determined by solving (2.25). However, the definition for the universal connection $`\widehat{A}`$ given in (2.12), the expression (2.13) for $`\widehat{d}`$ together with the explicit form (3.3) for $`A`$ suggest a simple guess for its ADHM expression; we write
$$c=U^{}(s+𝒞)U,$$
(3.45)
where $`𝒞`$ is the connection associated with the reparametrizations of the ADHM construction, which are shown in (3.12). Therefore, under these symmetries it transforms as
$$𝒞Q(𝒞+s)Q^{}.$$
(3.46)
In sec. 2 we observed that the first equation in (2.1) together with (2.22) imply that the BRST operator $`s`$ has an explicit realization on instanton moduli space as the exterior derivative. Since we are describing this space in terms of a redundant parametrization, every expression should be covariant with respect to the reparametrization symmetry group. This implies that ordinary derivatives on the instanton moduli space (which is described by the redundant set of $`8k+k(k1)/2`$ ADHM collective coordinates) have to be replaced by covariant ones, and $`s`$ by its covariant counterpart
$$𝒮=s+𝒞,$$
(3.47)
which is exactly what appears in (3.45). The criterion to fix $`𝒞`$ is clear: one has simply to plug (3.45) into (2.25), and solve for $`𝒞`$. In the next section we will illustrate an alternative (and quicker) way to construct $`𝒞`$. At this point we must alert the reader that the situation in which $`𝒜_{00}0`$ requires a more detailed analysis. As we will discuss in subsec. 3.3, in this case the correct guess for $`𝒞`$ is
$$𝒞=\left(\begin{array}{cccc}𝒞_{00}& 0& \mathrm{}& 0\\ 0& & & \\ \mathrm{}& & 𝒞^{}& \\ 0& & & \end{array}\right),$$
(3.48)
where $`𝒞^{}=(𝒞^{})^T`$ and $`𝒞_{00}`$ is non–vanishing when $`𝒜_{00}0`$. For the moment, let us only observe that (at least when $`𝒜_{00}=0`$), $`𝒞`$ is from its very definition a moduli–dependent $`(k+1)\times (k+1)`$ matrix antisymmetric in its lowest $`k\times k`$ block and zero elsewhere.
In summary, we are left with four ADHM matrices:
* $`\mathrm{\Delta }`$, which collects the ADHM data of the instanton configuration ($`8k+k(k1)/2`$ degrees of freedom),
* $``$, which parametrizes the tangent vectors to the instanton moduli space ($`8k`$ degrees of freedom),
* $`𝒜`$, which is the solution to (3.27), and
* $`𝒞`$, which is the matrix in (3.48).
Under the action of the group of reparametrization of the ADHM construction, $`\mathrm{\Delta }`$ and $``$ transform in the fundamental representation, whereas $`𝒞`$ transforms as a connection and $`𝒜`$ as the curvature of a connection. We want to warn the reader that only $`\mathrm{\Delta }`$ and $`𝒞`$ will emerge as independent quantities. Once they are given, all the other quantities (i.e. $``$ and $`𝒜`$) will be completely determined, as we will show in the next subsection.
### 3.2 The BRST Algebra in the ADHM Formalism: the Zero Vacuum Expectation Value Case
As we have seen, the TYM action projects the field $`A`$ onto the solutions of the self–duality equations (3.6), and the anti–commuting 1–form $`\psi `$ onto the tangent vectors to the instanton moduli space (the solutions to (3.17)); moreover, the scalar field $`\varphi `$ must satisfy (2.23), possibly with its boundary condition (3.24), and the ghost field $`c`$ satisfies (2.25), which is induced by the transversality condition of $`\psi `$ in the instanton background. In the last section we used the ADHM formalism to write the solutions to these coupled equations, that we collect here for the sake of clarity:
$`A`$ $`=`$ $`U^{}dU,`$
$`c`$ $`=`$ $`U^{}(s+𝒞)U,`$
$`\psi `$ $`=`$ $`U^{}f(d\mathrm{\Delta })^{}U+U^{}(d\mathrm{\Delta })f^{}U,`$
$`\varphi `$ $`=`$ $`U^{}f^{}U+U^{}𝒜U.`$ (3.49)
The BRST transformations of these fields are written in (2.1). If we now plug (3.2) into (2.1), we end with a set of equations which will provide us with explicit expressions for the variations $`s\mathrm{\Delta }`$, $`s`$, $`s𝒞`$, $`s𝒜`$ in terms of $`\mathrm{\Delta }`$, $``$, $`𝒞`$, $`𝒜`$. At the same time we will also show how to determine the explicit form of $`𝒞`$.
A preliminary ingredient which is necessary for this computation is the knowledge of $`sU`$; we would like to express it in terms of $`s\mathrm{\Delta }`$, otherwise we would be forced to solve the highly non–trivial set of algebraic equations (3.4), (3.5) for $`U`$. The following trick is then useful. Perform the BRST variation of (3.4),
$$(s\mathrm{\Delta })^{}U+\mathrm{\Delta }^{}sU=0;$$
(3.50)
this can be read as an equation for $`s\mathrm{\Delta }`$, whose solution is<sup>7</sup><sup>7</sup>7The following expression for $`sU`$ would still be valid if $`s`$ would represent a generic variation.
$$sU=\mathrm{\Delta }f(s\mathrm{\Delta })^{}U+U(U^{}sU).$$
(3.51)
We are now in a position to start computing $`sA`$ by varying the first equation of (3.2). This way we get
$$sA=U^{}\left[s\mathrm{\Delta }f(d\mathrm{\Delta })^{}+d\mathrm{\Delta }f(s\mathrm{\Delta })^{}\right]U[D,U^{}sU],$$
(3.52)
where, for a generic 1–form $`K`$, we put $`[D,K]=dK+AK+KA`$. Here and in the following we repeatedly use the fact that
$$\mathrm{\Delta }^{}dU=(d\mathrm{\Delta })^{}U,$$
(3.53)
which is a consequence of (3.4). We now substitute the explicit expressions found for $`\psi `$ and $`c`$ into the r.h.s. of the first of (2.1), thus getting
$`\psi Dc`$ $`=`$ $`U^{}(fd\mathrm{\Delta }^{}+d\mathrm{\Delta }f^{})U+`$ (3.54)
$`[D,U^{}sU][D,U^{}𝒞U].`$
If we equate the r.h.s. of (3.54) to the r.h.s. of (3.52) we obtain, after a little algebra,
$$U^{}(fd\mathrm{\Delta }^{}+d\mathrm{\Delta }f^{})U=U^{}\left[(s\mathrm{\Delta }+𝒞\mathrm{\Delta })fd\mathrm{\Delta }^{}+d\mathrm{\Delta }f(s\mathrm{\Delta }+𝒞\mathrm{\Delta })^{}\right]U;$$
(3.55)
from here we conclude that
$$=s\mathrm{\Delta }+𝒞\mathrm{\Delta }$$
(3.56)
modulo “irrelevant” terms, that is terms which vanish when right (left) multiplied by $`U`$ ($`U^{}`$). The same strategy as before can be repeatedly applied to the remaining equations in (2.1), thus obtaining the complete action of the BRST operator on $`\mathrm{\Delta }`$, $``$, $`𝒞`$, $`𝒜`$. The result of this exercise is
$`s\mathrm{\Delta }`$ $`=`$ $`𝒞\mathrm{\Delta },`$
$`s`$ $`=`$ $`𝒜\mathrm{\Delta }𝒞,`$
$`s𝒜`$ $`=`$ $`[𝒞,𝒜],`$ (3.57)
$`s𝒞`$ $`=`$ $`𝒜𝒞𝒞,`$
which is the realization of the BRST algebra on the instanton moduli space.<sup>8</sup><sup>8</sup>8Using (3.29) and similar expressions for $`𝒜`$, $`𝒞\mathrm{\Delta }`$ , $`𝒞`$, the $`x`$–dependence completely disappears from (3.2).
Three observations are in order. First, it is straightforward to show that $`s^2`$ is nilpotent as it should. This can be simply done by applying once again $`s`$ to each equation in (3.2). Therefore, on instanton moduli space $`s`$ is the exterior derivative, as we announced in the previous sections. Second, the last two equations in (3.2) and the nilpotency of $`s`$ suggest that $`𝒜`$ can be interpreted as the curvature of the connection $`𝒞`$ (these equations then becoming the Bianchi identity for $`𝒜`$ and its definition in terms of $`𝒞`$). Last, using the covariant derivative defined in (3.47), we can rewrite the BRST algebra on instanton moduli space in a more compact form as
$`𝒮\mathrm{\Delta }`$ $`=`$ $`,`$
$`𝒮`$ $`=`$ $`𝒜\mathrm{\Delta },`$
$`𝒮𝒜`$ $`=`$ $`0,`$ (3.58)
$`s𝒞+𝒞𝒞`$ $`=`$ $`𝒜.`$
We now discuss the important point of how to compute the connection $`𝒞`$. This can be done by plugging the first equation of (3.2) into the fermionic constraint (3.19), thus getting
$$\mathrm{\Delta }^{}𝒞\mathrm{\Delta }\left(\mathrm{\Delta }^{}𝒞\mathrm{\Delta }\right)^T=(\mathrm{\Delta }^{}s\mathrm{\Delta })^T\mathrm{\Delta }^{}s\mathrm{\Delta }.$$
(3.59)
It can also be shown that this equation is equivalent to the ADHM transcription of (2.25). In the following considerations we restrict our attention to the case in which $`𝒞_{00}=0`$ (recall (3.48)); as we said in the last section, this is true if and only if $`𝒜_{00}=0`$. The case $`(𝒞_{00},𝒜_{00})(0,0)`$ is crucially different and deeply related to the fact that when the scalar field acquires a non–zero v.e.v., the theory has a new invariance (the $`U(1)`$ central charge symmetry). For these reasons it will be separately analysed in sec.3.3. In this section we limit ourselves to $`𝒞`$’s of the form
$$𝒞_f=\left(\begin{array}{cccc}0& 0& \mathrm{}& 0\\ 0& & & \\ \mathrm{}& & 𝒞_f^{}& \\ 0& & & \end{array}\right),$$
(3.60)
where $`𝒞_f^{}=(𝒞_f^{})^T`$. Using the expression for the operator $`L`$ introduced in (3.31), we can rewrite (3.59) more compactly as
$$L𝒞_f^{}=\mathrm{\Lambda }_𝒞,$$
(3.61)
where we have defined
$$\mathrm{\Lambda }_𝒞\mathrm{\Delta }^{}s\mathrm{\Delta }(\mathrm{\Delta }^{}s\mathrm{\Delta })^T.$$
(3.62)
The solution to (3.61) is then formally written as
$$𝒞_f^{}=L^1\mathrm{\Lambda }_𝒞,$$
(3.63)
due to the invertibility of $`L`$.
The rôle of the connection $`𝒞`$ can be conveniently elucidated by setting it to zero in (3.2); in this case the BRST algebra would read<sup>9</sup><sup>9</sup>9The following transformations are in close relationship to the supersymmetry transformations of the ADHM matrices given in .
$`s\mathrm{\Delta }`$ $`=`$ $`,`$
$`s`$ $`=`$ $`𝒜\mathrm{\Delta },`$ (3.64)
$`s𝒜`$ $`=`$ $`0.`$
It can be immediately shown that in this case the operator $`s`$ would fail to be nilpotent. Indeed, the action of $`s^2`$ on the ADHM matrices would become
$`s^2\mathrm{\Delta }`$ $`=`$ $`𝒜\mathrm{\Delta },`$
$`s^2`$ $`=`$ $`𝒜,`$ (3.65)
$`s^2𝒜`$ $`=`$ $`0.`$
$`s^2`$ would then be nilpotent only up to transformations generated by $`k\times k`$ moduli-dependent antisymmetric matrices, i.e. local reparametrizations in the moduli space. (3.2) are the transcription of (2.6) on the moduli space.
In summary, the universal connection $`\widehat{A}`$ is given by
$$\widehat{A}=U^{}(d+s+𝒞)U.$$
(3.66)
We want now to comment on the interpretation of the results obtained in this section. The crucial observation is that, once (3.66) is given, the ADHM matrices $``$ and $`𝒜`$ are in turn determined by (3.2) as the covariant derivative of $`\mathrm{\Delta }`$ and the curvature of the connection $`𝒞`$ respectively; the only independent variables are the collective coordinates contained in $`\mathrm{\Delta }`$ (the instanton moduli and other moduli possibly associated with redundancies of the ADHM parametrization) and their differentials (the entries of the matrix $`s\mathrm{\Delta }`$). Once the reparametrization invariance has been gauged away (by giving some convenient prescription; see the explicit examples in sec. 4.1), physical quantities become, through their ADHM expression, differential forms<sup>10</sup><sup>10</sup>10When scalar fields have non–zero v.e.v.’s this picture is slightly modified; we postpone this discussion to sec. 3.3. on the $`8k`$–dimensional (anti–)instanton moduli space $`^+`$ ($`^{}`$). Operatively, this amounts to first identifying a correct parametrization for the instanton configuration (in term of the ADHM matrix $`\mathrm{\Delta }`$ introduced in (3.7)), and then to computing the explicit expression for the 1–form $`𝒞`$ using (3.19), in which $``$ is substituted by its expression (3.56). Finally, $`𝒜`$ is determined by the last equation in (3.2).
### 3.3 The BRST Algebra in the ADHM Formalism: the Non–Zero Vacuum Expectation Value Case
The realization of the BRST algebra (2.2) on instanton moduli space in the case in which scalar fields have non–vanishing v.e.v. closely parallels that of sec. 3.2. In particular, the universal connection $`\widehat{A}=A+c+\mathrm{\Lambda }`$ is again expressed as
$$\widehat{A}=U^{}(d+s+𝒞)U,$$
(3.67)
and its curvature equation and Bianchi identities give rise to the same algebra (3.2) for the ADHM matrices. Notice however that in this case $`𝒞`$ is given by
$$𝒞=\left(\begin{array}{cccc}𝒞_{00}& 0& \mathrm{}& 0\\ 0& & & \\ \mathrm{}& & 𝒞^{}& \\ 0& & & \end{array}\right),$$
(3.68)
the element $`𝒞_{00}`$ of (3.68) being related to the asymptotic behavior of the ghost $`c+\mathrm{\Lambda }`$ at $`|x|\mathrm{}`$,
$`\underset{|x|\mathrm{}}{lim}(c+\mathrm{\Lambda })\underset{|x|\mathrm{}}{lim}U^{}(s+𝒞)U=𝒞_{00}.`$ (3.69)
From (2.42) it then follows
$$𝒞_{00}=\lambda 𝒜_{00}.$$
(3.70)
The ADHM connection $`𝒞`$ can be calculated by solving (3.59); since
$$\mathrm{\Delta }^{}𝒞\mathrm{\Delta }(\mathrm{\Delta }^{}𝒞\mathrm{\Delta })^TL𝒞^{}+\mathrm{\Lambda }_b(𝒞_{00}),$$
(3.71)
then (3.59) can be written as
$$L𝒞^{}=\mathrm{\Lambda }_b(𝒞_{00})\mathrm{\Lambda }_𝒞,$$
(3.72)
where $`\mathrm{\Lambda }_b`$ has been defined in (3.32), and $`\mathrm{\Lambda }_𝒞`$ is given by (3.62). (3.72) and (3.59) are formally identical to (3.33), (3.27) respectively; as in sec. 3.1, they suggest us to set
$$𝒞=𝒞_b+𝒞_f,$$
(3.73)
where $`𝒞_b`$ satisfies the associated homogeneous equation
$$\mathrm{\Delta }^{}𝒞_b\mathrm{\Delta }(\mathrm{\Delta }^{}𝒞_b\mathrm{\Delta })^T=0.$$
(3.74)
The solution to (3.74) is (unique and) completely specified only after imposing boundary conditions, as in (3.70). This reflects, in the ADHM language, the degeneracy (2.40) in the definition of tangent vector to the instanton moduli space, which is due to the existence of central charge transformations; we know in fact that the fermionic constraint (3.19), from which (3.59) directly follows, is just the ADHM transcription of the fermionic zero–mode equations (3.17). As in that case, the solution is unique once the non–trivial boundary condition (3.70) is imposed.<sup>11</sup><sup>11</sup>11See also the discussion after (3.84). Let us now set
$$𝒞^{}=𝒞_f^{}+𝒞_b^{};$$
(3.75)
then (3.72) gives
$`L𝒞_b^{}`$ $`=`$ $`\mathrm{\Lambda }_b(𝒞_{00}),`$ (3.76)
$`L𝒞_f^{}`$ $`=`$ $`\mathrm{\Lambda }_𝒞,`$
whose solution is unique once $`𝒞_{00}`$ has been specified by means of the boundary condition (3.70). Note that, if $`𝒞_{00}`$ were zero, then also $`\mathrm{\Lambda }_b(𝒞_{00})`$ would vanish; therefore, due to the invertibility of $`L`$, the equation for $`𝒞_b^{}`$ would only admit the trivial solution $`𝒞_b^{}=0`$.
The matrices $``$ and $`𝒜`$ are in turn determined by means of the ADHM algebra to be
$`=𝒮\mathrm{\Delta },`$
$`𝒜=s𝒞+𝒞𝒞;`$ (3.77)
in particular, for the $`(00)`$ element of $`𝒜`$, we have
$`𝒜_{00}=s𝒞_{00}={\displaystyle \frac{}{\lambda }}\left(\lambda v{\displaystyle \frac{\sigma _3}{2i}}\right)=v{\displaystyle \frac{\sigma _3}{2i}},`$
$`s𝒜_{00}={\displaystyle \frac{}{\lambda }}\left(v{\displaystyle \frac{\sigma _3}{2i}}\right)=0,`$ (3.78)
from which it follows the expected asymptotic behavior (2.30) for the scalar field $`\varphi `$
$$\underset{|x|\mathrm{}}{lim}\varphi \underset{|x|\mathrm{}}{lim}U^{}𝒜U=𝒜_{00}=v\frac{\sigma _3}{2i}.$$
(3.79)
Note that the nilpotent BRST operator $`s`$ acts on the external parameters $`𝒞_{00},𝒜_{00}`$, given respectively by (3.70) and (3.78), just as the partial derivative with respect to $`\lambda `$, while its restriction on the other elements of the ADHM matrices would act as the usual exterior derivative on the moduli space.
### 3.4 Algebraic Construction of the BRST Transformations
In this section we will derive the realization of the BRST algebra on the instanton moduli space in a direct way, i.e. using neither the BRST algebra in field space (2.1) nor the expressions for the field configuration (3.2) onto which the TYM action projects. The only ingredient we need will be a parametrization for the moduli space of instantons; in terms of the ADHM construction, this is equivalent to determine a matrix $`\mathrm{\Delta }`$ which satisfies (3.7). As discussed in sec. 3.1, the ADHM space of parameters is acted upon by an $`O(k)`$ reparametrization symmetry. The gauging of this symmetry will turn out to be what is required to make the BRST variations of the ADHM data $`\mathrm{\Delta }`$ consistent with the algebraic constraints (3.7) which determine them. The BRST algebra on instanton moduli space, (3.2), will thus emerge as the most general set of deformations of the ADHM data compatible with (3.7).
To show this, let us now start by performing an infinitesimal scalar variation (that we call $`s`$ for obvious reasons) of the bosonic constraint (3.7). We get
$$(s\mathrm{\Delta })^{}\mathrm{\Delta }+\mathrm{\Delta }^{}s\mathrm{\Delta }=[(s\mathrm{\Delta })^{}\mathrm{\Delta }]^T+(\mathrm{\Delta }^{}s\mathrm{\Delta })^T.$$
(3.80)
This relation should be read as an equation for $`s\mathrm{\Delta }`$, and we want to guess its solution. We write it as
$$s\mathrm{\Delta }=𝒞\mathrm{\Delta },$$
(3.81)
where $``$ is defined as the matrix which satisfies (3.19). $`𝒞`$ is constrained by the structure of $`\mathrm{\Delta }`$ (which satisfies (3.14), (3.15) in the gauge defined by (3.13)) and $``$ (which is is fixed by (3.21), (3.22)); in conclusion, the most general expression of $`𝒞`$ consistent with (3.81) is
$$𝒞=\left(\begin{array}{cccc}𝒞_{00}& 0& \mathrm{}& 0\\ 0& & & \\ \mathrm{}& & 𝒞^{}& \\ 0& & & \end{array}\right),$$
(3.82)
where $`𝒞^{}`$ is a real antisymmetric $`k\times k`$ matrix, $`(𝒞^{})^{}=𝒞^{}`$. If we plug (3.81) into (3.80), the terms containing $``$ exactly cancel out thanks to (3.19), whereas $`𝒞`$ is fixed by the equation
$$\mathrm{\Delta }^{}(𝒞+𝒞^{})\mathrm{\Delta }=\left[\mathrm{\Delta }^{}(𝒞+𝒞^{})\mathrm{\Delta }\right]^T,$$
(3.83)
which becomes
$$L(𝒞^{}+𝒞_{}^{}{}_{}{}^{})=\mathrm{\Lambda }_b(𝒞_{00}+𝒞_{00}^{}),$$
(3.84)
where $`L`$ is a generally invertible operator. As a consequence one must have $`𝒞_{00}=𝒞_{00}^{}`$. One is thus led to an expression of $`𝒞`$ which coincides with the one previously suggested in (3.48).
Let us now pause for a moment and count the number of degrees of freedom in (3.81). On one hand, the ADHM matrix $`\mathrm{\Delta }`$ and its variation $`s\mathrm{\Delta }`$ both contain $`8k+k(k1)/2`$ (unconstrained) degrees of freedom<sup>12</sup><sup>12</sup>12We recall the reader that we always work with a canonical choice of $`b`$.; $``$ contains instead $`8k`$ degrees of freedom, after solving (3.19). On the other hand, $`𝒞`$ contains $`3+k(k1)/2`$ parameters, and we would be led to an apparent mismatch in counting the number of degrees of freedom in (3.81). Actually the three degrees of freedom introduced by $`𝒞_{00}`$ are not new; instead they are already included in the number of independent solutions to (3.19). This can be understood once we decompose $`𝒞`$ as in (3.73), where
$$𝒞_b=\left(\begin{array}{cccc}𝒞_{00}& 0& \mathrm{}& 0\\ 0& & & \\ \mathrm{}& & 𝒞_b^{}& \\ 0& & & \end{array}\right)$$
(3.85)
is defined in such a way to satisfy the homogeneous equation (3.74). On one hand this equation is equivalent to
$$L𝒞_b^{}=\mathrm{\Lambda }_b(𝒞_{00});$$
(3.86)
therefore it has non–trivial solutions when $`𝒞_{00}0`$. On the other hand, it is formally identical to (3.19) with $``$ replaced by $`𝒞_b\mathrm{\Delta }`$. In order to avoid double counting, its three independent solutions should then not be considered as “new”. Finally, $`𝒞_f`$ is now constrained by (3.73), (3.81), (3.19) to satisfy an equation identical to (3.59); thus, it just takes into account the genuinely new $`k(k1)/2`$ parameters which are related to the $`O(k)`$ reparametrization invariance. We then conclude that there is a complete balance in the number of degrees of freedom in (3.81).
If we perform the $`s`$–variation of the fermionic constraint (3.19), we get
$$(s\mathrm{\Delta })^{}+\mathrm{\Delta }^{}s=[(s\mathrm{\Delta })^{}]^T+(\mathrm{\Delta }^{}s)^T.$$
(3.87)
Taking into account (3.81), we find
$$(^{}+\mathrm{\Delta }^{}𝒞)[(^{}+\mathrm{\Delta }^{}𝒞)]^T=(\mathrm{\Delta }^{}s)^T\mathrm{\Delta }^{}s;$$
(3.88)
(3.88) should be thought of as an equation for $`s`$. Analogously to the previous case, its most general solution can be cast into the form
$$s=𝒜\mathrm{\Delta }𝒞,$$
(3.89)
where $`𝒜`$ has the same form as in (3.26). If we plug (3.89) into (3.88), we obtain that $`𝒜`$ must satisfy the following relation:
$$\mathrm{\Delta }^{}𝒜\mathrm{\Delta }(\mathrm{\Delta }^{}𝒜\mathrm{\Delta })^T=(^{})^T^{}.$$
(3.90)
Then (3.90) is identical to (3.27), which was obtained from a completely different point of view (the equations of motion for the scalar field $`\varphi `$), and its solution is given by (3.34), (3.35), (3.36).
We want now to clarify the relation between $`𝒜`$ and $`𝒞`$ as defined in this section. To this end, let us perform one more $`s`$–variation of (3.81) and (3.89); after a little algebra we get
$`s^2\mathrm{\Delta }`$ $`=`$ $`\left(𝒜s𝒞𝒞𝒞\right)\mathrm{\Delta },`$
$`s^2`$ $`=`$ $`\left(𝒜s𝒞𝒞𝒞\right)+\left(s𝒜+[𝒞,𝒜]\right)\mathrm{\Delta }.`$ (3.91)
Once one requires the nilpotency of the BRST operator $`s`$, then
$`𝒜s𝒞𝒞𝒞=0,`$ (3.92)
$`s𝒜+[𝒞,𝒜]=0.`$ (3.93)
Therefore it is possible to interpret (3.92) as the definition of $`𝒜`$ as the field strength of $`𝒞`$ and (3.93) as its Bianchi identity. This completely clarifies the relation between $`𝒜`$ and $`𝒞`$.
In order to check the consistency of the super–constraints with the BRST variations, we still have to perform the $`s`$–variation of (3.90). If we do this, we get
$$\mathrm{\Delta }^{}(s𝒜+[𝒞,𝒜])\mathrm{\Delta }[\mathrm{\Delta }^{}(s𝒜+[𝒞,𝒜])\mathrm{\Delta }]^T=0,$$
(3.94)
which is trivially satisfied thanks to (3.93).
Summarizing, we have found that consistency between the BRST variation of the bosonic ADHM matrix $`\mathrm{\Delta }`$ and the constraint (3.7) it obeys, yields
$`s\mathrm{\Delta }`$ $`=`$ $`𝒞\mathrm{\Delta },`$
$`s`$ $`=`$ $`𝒜\mathrm{\Delta }𝒞,`$
$`s𝒜`$ $`=`$ $`[𝒞,𝒜],`$ (3.95)
$`s𝒞`$ $`=`$ $`𝒜𝒞𝒞,`$
where $``$ satisfies (3.19). As anticipated, this set of equations gives an explicit realization of the BRST algebra on instanton moduli space, and it coincides with that found in (3.2) with completely different methods.
## 4 The Set–Up of the Calculation of Instanton Green’s Functions
In this section we explain how to perform instanton calculations in our picture. As an application of our techniques, we will then focus on computing correlators in the case in which the relevant instanton configurations have winding number $`k=1,2`$. In $`N=2`$ SYM with non–vanishing v.e.v. for the scalar field, we will be interested in evaluating the correlator $`<\text{Tr}\varphi ^2>`$. These computations will show the main features of the formalism developed in the previous sections.
To make these characteristics more evident, let us now summarize the “standard” strategy to perform instanton calculations in SUSY theories .
* The action is expanded around the saddle point up to quadratic fluctuations.
* The fields are expanded in eigenmodes and the functional measure is replaced by an integration over the coefficients of the mode expansion. The contribution of the zero–modes and that of the non–zero modes are now clearly identified.
* The fields in the correlator are also expanded in modes and the part containing the non–zero modes is discarded since it represents higher order quantum corrections.
* The non–zero modes are then integrated out. This integration gives a ratio of determinants which is one thanks to SUSY .
* The last step consists in performing the integration over the zero–modes. In order to deal with the zero–mode sector, one has to trade integrations over the bosonic zero–modes for integrations over collective coordinates; this gives rise to a bosonic Jacobian. Moreover, one has to keep into account chiral selection rules which single out the non–vanishing Green’s functions. Operatively, these selection rules amount to say that all the Grassmann integrations over the fermionic collective coordinates have to be soaked up by explicitly inserting the appropriate number (say $`n`$) of zero–modes; thus, the only non–zero amplitudes will be those which admit an expansion in terms of fermion zero–modes such that the coefficient multiplying the term with $`n`$ fermionic collective coordinates does not vanish. This gives rise to a fermionic Jacobian, which is the determinant of the matrix whose entries are the overlaps of the fermionic zero–mode wave functions.
Our starting point will be the last step which, in the formalism of the previous sections, amounts to integrating the Lagrange multipliers in the gauge fixed TYM action (2.44). This integration naturally projects the fields $`A`$, $`\psi `$, $`\varphi `$ onto the zero–modes subspace, which is identified by (3.6), (3.17) and (2.23) (supplemented by appropriate boundary conditions on $`\varphi `$)<sup>13</sup><sup>13</sup>13If one wanted to work with anti–instantons, then (3.6) and (3.17) should be replaced by (2.21) and (2.22), respectively.; the configurations which solve these equations were written in (3.2). Through these expressions, physical amplitudes will depend on $`\mathrm{\Delta }`$, $`𝒞`$, $``$ and $`𝒜`$. The ADHM equations (3.7) fix the number of independent (bosonic) collective coordinates to be $`8k+k(k1)/2`$; gauge–fixing the left–over $`O(k)`$ symmetry further reduce this number to $`8k`$. Moreover the first relation in (3.2) together with (3.19) allows one to compute the connection $`𝒞^{}`$ as a 1–form expanded on a basis of differentials of the bosonic moduli. If one substitutes back the computed expression for $`𝒞`$ into the first equation in (3.2), then the $``$’s become in turn differential 1–forms on instanton moduli space $`^+`$.<sup>14</sup><sup>14</sup>14Recall that the the number of independent $``$’s is $`8k`$ (as the number of bosonic moduli) by virtue of (3.19). Finally (3.27) gives $`𝒜`$ as a function of $`\mathrm{\Delta }`$ and $``$. We then conclude that any polynomial in the fields becomes, after projection onto the zero–mode subspace, a well–defined differential form on $`^+`$ . We can then symbolically write
$$fields=_^+\left[(fields)e^{S_{\mathrm{TYM}}}\right]_{zeromodesubspace}.$$
(4.1)
Let us now call $`\{\widehat{\mathrm{\Delta }}_i\}`$ ($`\{\widehat{}_i\}`$), $`i=1,\mathrm{},p`$, where $`p=8k`$, a basis of (ADHM) coordinates on $`^+`$ ($`T_A^+`$). (3.81) thus yields $`\widehat{}_i=s\widehat{\mathrm{\Delta }}_i+(\widehat{𝒞\mathrm{\Delta }})_i`$. A generic function on the zero–mode subspace will then have the expansion
$`g(\widehat{\mathrm{\Delta }},\widehat{})`$ $`=`$ $`g_0(\widehat{\mathrm{\Delta }})+g_{i_1}(\widehat{\mathrm{\Delta }})\widehat{}_{i_1}+{\displaystyle \frac{1}{2!}}g_{i_1i_2}(\widehat{\mathrm{\Delta }})\widehat{}_{i_1}\widehat{}_{i_2}+\mathrm{}`$ (4.2)
$`+`$ $`{\displaystyle \frac{1}{p!}}g_{i_1i_2\mathrm{}i_p}(\widehat{\mathrm{\Delta }})\widehat{}_{i_1}\widehat{}_{i_2}\mathrm{}\widehat{}_{i_p},`$
the coefficients of the expansion being totally antisymmetric in their indices. Now the first of (3.2) implies that the $`\widehat{}_i`$’s and the $`s\widehat{\mathrm{\Delta }}_i`$’s are related by a (moduli–dependent) linear transformation $`K_{ij}`$, which is completely known once the explicit expression for $`𝒞`$ is plugged into the $`\widehat{}_i`$’s:
$$\widehat{}_i=K_{ij}(\widehat{\mathrm{\Delta }})s\widehat{\mathrm{\Delta }}_j.$$
(4.3)
It then follows that
$`\widehat{}_{i_1}\widehat{}_{i_2}\mathrm{}\widehat{}_{i_p}`$ $`=`$ $`K_{i_1j_1}K_{i_2j_2}\mathrm{}K_{i_pj_p}s\widehat{\mathrm{\Delta }}_{j_1}s\widehat{\mathrm{\Delta }}_{j_2}\mathrm{}s\widehat{\mathrm{\Delta }}_{j_p}=`$ (4.4)
$`=`$ $`ϵ_{j_1\mathrm{}j_p}K_{i_1j_1}K_{i_2j_2}\mathrm{}K_{i_pj_p}s^p\widehat{\mathrm{\Delta }}=`$
$`=`$ $`ϵ_{i_1\mathrm{}i_p}(\mathrm{det}K)s^p\widehat{\mathrm{\Delta }},`$
where $`s^p\widehat{\mathrm{\Delta }}s\widehat{\mathrm{\Delta }}_1\mathrm{}s\widehat{\mathrm{\Delta }}_p`$. From (4.2), (4.3) we conclude that
$`{\displaystyle _^+}g(\widehat{\mathrm{\Delta }},\widehat{})`$ $`=`$ $`{\displaystyle \frac{1}{p!}}{\displaystyle _^+}g_{i_1i_2\mathrm{}i_p}(\widehat{\mathrm{\Delta }})\widehat{}_{i_1}\widehat{}_{i_2}\mathrm{}\widehat{}_{i_p}=`$ (4.5)
$`=`$ $`{\displaystyle _^+}s^p\widehat{\mathrm{\Delta }}|\mathrm{det}K|g_{12\mathrm{}p}(\widehat{\mathrm{\Delta }}).`$
This formula is an operative tool to calculate physical amplitudes. Here the determinant of $`K`$ naturally stands out as the instanton integration measure for $`N=2`$ SYM theories. This important ingredient of the calculation is obtained in standard instanton calculations as a ratio of bosonic and fermionic zero–mode Jacobians. Instead, in our approach it emerges in a geometrical and very direct way, without the need of any computation of ratios of determinants, nor of any knowledge of the explicit expressions of bosonic and fermionic zero–modes. The only ingredient is the connection $`𝒞`$. As an instructive exercise, in the following subsection we will compute $`K`$ and its determinant (i.e. the instanton measure) in the cases of winding number equal to one and two. We anticipate that the results we get will agree with previously known formulae; however they are obtained here in a very quick and straightforward way.
A last remark concerns what happens to the action of the theory, $`S_{\mathrm{TYM}}`$, when it is restricted to the zero–mode subspace (we called the corresponding expression $`S_{\mathrm{inst}}`$ for obvious reasons). In sec. 2.2 we saw that for $`v0`$ it is non–vanishing; its expression was given in (2.49). In the following we will need to explicitly compute $`S_{\mathrm{inst}}`$ as a function of the instanton moduli; this will be done in sec. 4.2, where we will also be able to write it as a total BRST derivative.
### 4.1 The Instanton Measure for Winding Number $`k=1,2`$
The ADHM bosonic and fermionic matrices can be written as
$$\mathrm{\Delta }=\left(\begin{array}{c}w\\ x_0x\end{array}\right),=\left(\begin{array}{c}\mu \\ \xi \end{array}\right).$$
(4.6)
For $`k=1`$ there are no constraints over the collective coordinates; therefore the left–over reparametrization group introduced in (3.12) is trivial. As a result we simply have
$$=\left(\begin{array}{c}sw\\ sx_0\end{array}\right),$$
(4.7)
and $`\mathrm{det}K=1`$. The instanton measure is then given by $`s^4x_0s^4w`$, which is the well–known ’t Hooft measure . We now move on to the more interesting case of $`k=2`$.
The ADHM bosonic matrix reads
$$\mathrm{\Delta }=\left(\begin{array}{cc}w_1& w_2\\ x_1x& a_1\\ a_1& x_2x\end{array}\right)=\left(\begin{array}{cc}w_1& w_2\\ a_3& a_1\\ a_1& a_3\end{array}\right)+b(xx_0),$$
(4.8)
where $`x_0=(x_1+x_2)/2`$, $`a_3=(x_1x_2)/2`$. We also need the expression of the matrix $``$ which is defined in (3.19). Since this constraint is very similar to (3.7) (to get convinced of this fact just think that two solutions of (3.19) are given by $``$ proportional to $`a`$ and $`b`$) it is convenient to choose a form of $``$ which parallels (4.8), i.e.
$$=\left(\begin{array}{cc}\mu _1& \mu _2\\ \xi +_3& _1\\ _1& \xi _3\end{array}\right)=\left(\begin{array}{cc}\mu _1& \mu _2\\ _3& _1\\ _1& _3\end{array}\right)b\xi .$$
(4.9)
The solution to the bosonic constraint (3.7) is simply given by
$$a_1=\frac{1}{4|a_3|^2}a_3(\overline{w}_2w_1\overline{w}_1w_2+\mathrm{\Sigma }),$$
(4.10)
where $`\mathrm{\Sigma }`$ is an arbitrary real parameter related to the left–over $`O(2)`$ symmetry. In the following we will exploit this $`O(2)`$ gauge freedom to put $`\mathrm{\Sigma }`$ to zero. The constraint (3.19) is satisfied imposing
$$_1=\frac{a_3}{2|a_3|^2}(2\overline{a}_1_3+\overline{w}_2\mu _1\overline{w}_1\mu _2);$$
(4.11)
from now on we will choose $`\{\mu _1,\mu _2,\xi ,_3\}`$ as a set of independent fermionic variables. Finally, the equation $`L𝒜^{}=\mathrm{\Lambda }_b\mathrm{\Lambda }_f`$ reduces to
$`H(𝒜_f^{})_{12}`$ $`=`$ $`(\mathrm{\Lambda }_f)_{12}\overline{\mu }_1\mu _2\overline{\mu }_2\mu _1+2(\overline{}_3_1\overline{}_1_3),`$ (4.12)
$`H(𝒜_b^{})_{12}`$ $`=`$ $`(\mathrm{\Lambda }_b)_{12}\overline{w}_1𝒜_{00}w_2\overline{w}_2𝒜_{00}w_1,`$ (4.13)
where $`H=|w_1|^2+|w_2|^2+4(|a_3|^2+|a_1|^2)`$. Let us now write the BRST transformations of the bosonic ADHM variables:
$`\mu _1`$ $`=`$ $`sw_1+𝒞_{12}w_2+𝒞_{00}w_1,`$
$`\mu _2`$ $`=`$ $`sw_2𝒞_{12}w_1+𝒞_{00}w_2,`$
$`\xi `$ $`=`$ $`sx_0,`$
$`_3`$ $`=`$ $`sa_3+2𝒞_{12}a_1,`$
$`_1`$ $`=`$ $`sa_12𝒞_{12}a_3.`$ (4.14)
The component $`𝒞_{12}`$ of the $`O(2)`$ connection
$$𝒞^{}=\left(\begin{array}{cc}0& 𝒞_{12}\\ 𝒞_{12}& 0\end{array}\right)$$
(4.15)
can be simply obtained by plugging the right hand sides of (4.1) into the fermionic constraint $`(\mathrm{\Delta }^{})_{12}=(\mathrm{\Delta }^{})_{21}`$ and solving for $`𝒞_{12}`$. Actually, the terms containing $`𝒞_{00}`$, which is given by (3.70), can be discarded, since they do not contribute upon integration on the instanton moduli space. This way we get
$$𝒞_{12}=\frac{1}{H}\left[\overline{w}_1sw_2\overline{w}_2sw_1+2(\overline{a}_3sa_1\overline{a}_1sa_3)\right].$$
(4.16)
Eliminating $`sa_1`$ via (4.10) (in the gauge $`\mathrm{\Sigma }=0`$), one can rewrite $`𝒞_{12}`$ in terms of differentials of independent bosonic moduli, thus obtaining
$`𝒞_{12}`$ $`=`$ $`{\displaystyle \frac{1}{2H}}[\overline{w}_1sw_2\overline{w}_2sw_14\overline{a}_1sa_3+`$ (4.17)
$`+`$ $`s\overline{w}_2w_1s\overline{w}_1w_24s\overline{a}_3a_1].`$
Two observations are in order. First, we remark that (4.17) clearly shows that $`𝒞_{12}`$ is real, as a connection of an orthogonal group should. Moreover, the r.h.s. of (4.17) does not depend on $`sx_0`$; for this reason and from (4.1) it immediately follows that in computing $`\mathrm{det}K`$ we can discard the variable $`\xi sx_0`$, which would contribute with the determinant of a unit matrix. We will then define a “reduced” fermionic matrix (of quaternions) $`\stackrel{~}{}`$ as
$$\stackrel{~}{}=\left(\begin{array}{c}\mu _1\\ \mu _2\\ _3\end{array}\right),$$
(4.18)
its bosonic counterpart being
$$\stackrel{~}{\mathrm{\Delta }}=\left(\begin{array}{c}w_1\\ w_2\\ a_3\end{array}\right).$$
(4.19)
The relation between $`\stackrel{~}{}`$ and $`s\stackrel{~}{\mathrm{\Delta }}`$ can be cast into the form
$$(\stackrel{~}{}_{\alpha \dot{\alpha }})_i=\sigma _{\alpha \dot{\alpha }}^\mu (K_{\mu \nu })_{ij}(s\stackrel{~}{\mathrm{\Delta }}_\nu )_j,$$
(4.20)
where $`i=1,2,3`$. Plugging (4.17) into (4.1), we get, after a little algebra, the following explicit expression for $`K`$,
$$(K_{\mu \nu })_{ij}=\left(\begin{array}{ccc}\delta _{\mu \nu }w_{2\mu }w_{2\nu }/H& w_{2\mu }w_{1\nu }/H& 4w_{2\mu }a_{1\nu }/H\\ \\ w_{1\mu }w_{2\nu }/H& \delta _{\mu \nu }w_{1\mu }w_{1\nu }/H& 4w_{1\mu }a_{1\nu }/H\\ \\ 2a_{1\mu }w_{2\nu }/H& 2a_{1\mu }w_{1\nu }/H& \delta _{\mu \nu }8a_{1\mu }a_{1\nu }/H\end{array}\right),$$
(4.21)
whose determinant we want now to compute.
To this end, let us write $`K`$ as
$$K=\text{1 l}zz^TP=(P^1zz^T)P,$$
(4.22)
where
$$P=\left(\begin{array}{ccc}\text{1 l}& 0& 0\\ 0& \text{1 l}& 0\\ 0& 0& 2\text{1 l}\end{array}\right),$$
(4.23)
and
$$z=\frac{1}{\sqrt{H}}\left(\begin{array}{c}w_2\\ w_1\\ 2a_1\end{array}\right).$$
(4.24)
It is easy to verify that the determinant of a matrix of the form
$$Q=\mathrm{diag}(\alpha _1,\mathrm{},\alpha _n)zz^T,z=\left(\begin{array}{c}z_1\\ \mathrm{}\\ z_n\end{array}\right),$$
(4.25)
is simply
$$detQ=\underset{i=1}{\overset{n}{}}\alpha _i\underset{i=1}{\overset{n}{}}\left(\underset{ji}{}\alpha _j\right)|z_i|^2,$$
(4.26)
from which it is straightforward to get
$$|detK|=\frac{4\left||a_3|^2|a_1|^2\right|}{H}.$$
(4.27)
Restoring the $`\mathrm{\Sigma }`$ dependence of $`a_1`$ and $`sa_1`$ in (4.17) we obtain
$`𝒞_{12}`$ $`=`$ $`{\displaystyle \frac{1}{2H}}[\overline{w}_1sw_2\overline{w}_2sw_14\overline{a}_1sa_3+`$ (4.28)
$`+`$ $`s\overline{w}_2w_1s\overline{w}_1w_24s\overline{a}_3a_1+s\mathrm{\Sigma }],`$
where $`\mathrm{\Sigma }=\mathrm{\Sigma }(w_1,w_2,a_3,x_0)`$. $`𝒞_{12}`$ contains also a term proportional to $`sx_0`$; however, this term turns out not to contribute to $`detK`$ and, in fact, we find
$$|detK|=\frac{4}{H}\left||a_3|^2|a_1|^2+\frac{1}{4}\frac{\mathrm{\Sigma }}{a_{3\mu }}a_{1\mu }+\frac{1}{8}\frac{\mathrm{\Sigma }}{w_{1\mu }}w_{2\mu }\frac{1}{8}\frac{\mathrm{\Sigma }}{w_{2\mu }}w_{1\mu }\right|.$$
(4.29)
It is now possible to write the terms containing $`\mathrm{\Sigma }`$ in a more compact way. To this end, recall that the action of the $`O(2)`$ reparametrization group on $`(w_1,w_2,a_3,a_1)`$ can be read from (3.12) and (3.16) with $`q=\text{1 l}`$ and
$$R=\left(\begin{array}{cc}\mathrm{cos}\theta & \mathrm{sin}\theta \\ \mathrm{sin}\theta & \mathrm{cos}\theta \end{array}\right).$$
(4.30)
It is straightforward to show that
$`a_{1\mu }`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{a_{3\mu }^\theta }{\theta }}|_{\theta =0},`$
$`w_{2\mu }`$ $`=`$ $`{\displaystyle \frac{w_{1\mu }^\theta }{\theta }}|_{\theta =0},`$
$`w_{1\mu }`$ $`=`$ $`{\displaystyle \frac{w_{2\mu }^\theta }{\theta }}|_{\theta =0},`$ (4.31)
so that we can finally rewrite (4.29) as
$$|detK|=\frac{4||a_3|^2|a_1|^2\frac{1}{8}\frac{\mathrm{\Sigma }^\theta }{\theta }|_{_{\theta =0}}|}{H}.$$
(4.32)
This result exactly gives the instanton measure for $`N=2`$ pure SYM, which was obtained in as a ratio of fermionic and bosonic Jacobians. In our formalism it is simply the determinant of a coordinate transformation, and it is possible to write it down with the precise knowledge of $`𝒞`$ alone. We will further clarify the rôle of $`𝒞`$ in the last section of this work, where we will construct the moduli space of self–dual gauge connections as a hyperkähler quotient space.
### 4.2 The Multi–Instanton Action from the TYM Action
When restricted to the zero–mode subspace, the TYM action vanishes up to the boundary term written in (2.49), which leads to the non–trivial multi–instanton action
$$S_{\mathrm{inst}}=sd^4x2^\mu \text{Tr}(\overline{\varphi }\psi _\mu )=4\pi ^2\underset{|x|\mathrm{}}{lim}|x|^3\frac{x^\mu }{|x|}s\text{Tr}(\overline{\varphi }\psi _\mu ).$$
(4.33)
In the limit $`|x|\mathrm{}`$, the only non–vanishing term is given by $`\text{Tr}[\overline{\varphi }(s\psi _\mu )]`$. Let us now calculate the asymptotic limits of $`\overline{\varphi }`$ and $`\psi _\mu `$ for winding number $`k`$. For the scalar field we trivially have $`lim_{|x|\mathrm{}}\overline{\varphi }=\overline{𝒜}_{00}`$, while for $`\psi _\mu `$
$$\underset{|x|\mathrm{}}{lim}\psi _\mu =\underset{|x|\mathrm{}}{lim}\left(U^{}fb\overline{\sigma }_\mu UU^{}b\sigma _\mu f^{}U\right).$$
(4.34)
Knowing that asymptotically
$`U_0`$ $``$ $`\sigma _0,`$
$`U_k`$ $``$ $`{\displaystyle \frac{1}{|x|^2}}x\overline{w}_kU_0,`$ (4.35)
$`f_{kl}`$ $``$ $`{\displaystyle \frac{1}{|x|^2}}\delta _{kl},`$
we get
$`\underset{|x|\mathrm{}}{lim}x^\mu \psi _\mu `$ $`=`$ $`\underset{|x|\mathrm{}}{lim}x^\mu {\displaystyle \underset{p,l,m=1}{\overset{k}{}}}[U_0^{}\mu _p\left({\displaystyle \frac{1}{|x|^2}}\delta _{pl}\right)b_{lm}\overline{\sigma }_\mu U_m+`$ (4.36)
$`U_p^{}b_{pl}\sigma _\mu \left({\displaystyle \frac{1}{|x|^2}}\delta _{lm}\right)\overline{\mu }_mU_0]=`$
$`=`$ $`{\displaystyle \frac{1}{|x|^2}}{\displaystyle \underset{l=1}{\overset{k}{}}}(\mu _l\overline{w}_lw_l\overline{\mu }_l);`$
from this we conclude that
$`[S_{\mathrm{inst}}]_k`$ $`=`$ $`4\pi ^2s\text{Tr}\left[\overline{𝒜}_{00}{\displaystyle \underset{l=1}{\overset{k}{}}}\left(\mu _l\overline{w}_lw_l\overline{\mu }_l\right)\right]`$ (4.37)
$`=`$ $`4\pi ^2s\text{Tr}\left[\overline{𝒜}_{00}(\mathrm{\Delta }^{}\mathrm{\Delta }^{})_{00}\right]=4\pi ^2s\text{Tr}\left[\overline{𝒜}_{00}\left(\mathrm{\Delta }\stackrel{}{𝒮}\mathrm{\Delta }^{}\right)_{00}\right],`$
which explicitly gives the instanton action as the total BRST variation of a function of the bosonic and fermionic collective coordinates. In the second equality, the subscript $`00`$ stands for the upper left entry of the matrix in parentheses, and $`𝒮`$ is the covariant derivative on instanton moduli space defined in (3.47). Note that only some moduli are involved in this expression, more precisely, the unconstrained ones.
It is easy to convince oneself that (4.37) reproduces the instanton action for $`N=2`$ SYM as written in . To this end, let us now act with the operator $`s`$ on the moduli. From (3.2) and (3.29) it follows that
$`sw_l`$ $`=`$ $`\mu _l(𝒞_{00}w_l{\displaystyle \underset{p=1}{\overset{k}{}}}w_p𝒞_{pl}^{}),`$ (4.38)
$`s\mu _l`$ $`=`$ $`𝒜_{00}w_l{\displaystyle \underset{p=1}{\overset{k}{}}}w_p𝒜_{pl}^{}𝒞_{00}\mu _l;`$
we get
$`s\text{Tr}\left[\overline{𝒜}_{00}{\displaystyle \underset{l=1}{\overset{k}{}}}\left(\mu _l\overline{w}_lw_l\overline{\mu }_l\right)\right]`$ $`=`$ $`\text{Tr}[2(\overline{𝒜}_{00}𝒜_{00}){\displaystyle \underset{l=1}{\overset{k}{}}}|w_l|^22\overline{𝒜}_{00}{\displaystyle \underset{l=1}{\overset{k}{}}}\mu _l\overline{\mu }_l+`$ (4.39)
$`2\overline{𝒜}_{00}{\displaystyle \underset{l,p=1}{\overset{k}{}}}w_l𝒜_{lp}^{}\overline{w}_p+`$
$`+{\displaystyle \underset{l=1}{\overset{k}{}}}\left(\overline{𝒜}_{00}𝒞_{00}+𝒞_{00}\overline{𝒜}_{00}\right)\mu _l\overline{w}_l+`$
$`+{\displaystyle \underset{l=1}{\overset{k}{}}}\left(\overline{𝒜}_{00}𝒞_{00}𝒞_{00}\overline{𝒜}_{00}\right)w_l\overline{\mu }_l,`$
where we also used the fact that $`s𝒜_{00}=0`$. The last two terms in (4.39) vanish by virtue of (3.70); therefore, we conclude that
$$\left[S_{\mathrm{inst}}\right]_k=4\pi ^2\text{Tr}\left[2(\overline{𝒜}_{00}𝒜_{00})\underset{l=1}{\overset{k}{}}|w_l|^22\overline{𝒜}_{00}\underset{l=1}{\overset{k}{}}\mu _l\overline{\mu }_l+\underset{l,p=1}{\overset{k}{}}(\overline{w}_l\overline{𝒜}_{00}w_p\overline{w}_p\overline{𝒜}_{00}w_l)𝒜_{lp}^{}\right]$$
(4.40)
which exactly reproduces the $`N=2`$ SYM action in moduli space<sup>15</sup><sup>15</sup>15Note that in our notations $`2\text{Tr}(\overline{𝒜}_{00}𝒜_{00})=|v|^2`$. . $`\left[S_{\mathrm{inst}}\right]_k`$ can be decomposed as $`\left[S_{\mathrm{inst}}\right]_k=\left[S_B\right]_k+\left[S_F\right]_k`$, where $`\left[S_B\right]_k`$ ($`\left[S_F\right]_k`$) is the Higgs action (the Yukawa action) for instanton number $`k`$. Explicitly
$`\left[S_B\right]_k`$ $`=`$ $`4\pi ^2\text{Tr}\left[2(\overline{𝒜}_{00}𝒜_{00}){\displaystyle \underset{l=1}{\overset{k}{}}}|w_l|^2+{\displaystyle \underset{l,p=1}{\overset{k}{}}}(\overline{w}_l\overline{𝒜}_{00}w_p\overline{w}_p\overline{𝒜}_{00}w_l)(𝒜_b^{})_{lp}\right],`$ (4.41)
$`\left[S_F\right]_k`$ $`=`$ $`4\pi ^2\text{Tr}\left[2\overline{𝒜}_{00}{\displaystyle \underset{l=1}{\overset{k}{}}}\mu _l\overline{\mu }_l+{\displaystyle \underset{l,p=1}{\overset{k}{}}}(\overline{w}_l\overline{𝒜}_{00}w_p\overline{w}_p\overline{𝒜}_{00}w_l)(𝒜_f^{})_{lp}\right],`$ (4.42)
$`𝒜_b^{}`$ and $`𝒜_f^{}`$ being defined in (3.34).
We are now ready to perform explicit instanton calculations in our framework.
## 5 Computation of Instanton–Dominated Correlators in the Seiberg–Witten Model
The strategy for computing instanton–dominated correlators in our set–up has been described at the beginning of sec. 4. Here we focus the attention on the Green’s function $`\text{Tr}\varphi ^2`$, which is relevant for the computation of the Seiberg–Witten low–energy effective action. To begin with, notice that the group of translations in $`\text{IR}^4`$ is a symmetry of the theory even in the case $`v0`$; as a consequence $`x_0`$ and its supersymmetric counterpart $`sx_0\xi `$ (which is naturally expressed as the BRST variation of the instanton configuration center $`x_0`$) do not appear in $`S_{\mathrm{inst}}`$, as a direct check of (4.40) also shows. They will then have to be soaked up by selecting the translational part in the correlator insertion $`\text{Tr}\varphi ^2`$; this amounts to performing the replacement $`\varphi F_{\mu \nu }(1/2)(sx_{0\mu }sx_{0\nu })`$, i.e.
$$\text{Tr}\varphi ^2\frac{1}{2}\text{Tr}(F_{\mu \nu }\stackrel{~}{F}_{\mu \nu })s^4x_0.$$
(5.1)
The integral over these collective coordinates can now be easily performed giving the winding number ,
$$_{\{x_0\}}\frac{\text{Tr}\varphi ^2}{8\pi ^2}k.$$
(5.2)
Therefore, we get
$$<\frac{\text{Tr}\varphi ^2}{8\pi ^2}>_k=k_{^+\backslash \{x_0\}}e^{[S_{\mathrm{inst}}]_k}.$$
(5.3)
The last step consists in integrating the exponential of the instanton action over the remaining collective coordinates, i.e. over the “reduced” moduli space $`^+\backslash \{x_0\}`$, whose dimension is $`4n`$ where $`n=2k1`$. Let us call $`\stackrel{~}{\mathrm{\Delta }}_i`$, $`i=1,\mathrm{},n`$ the ADHM data for $`^+\backslash \{x_0\}`$, and $`\stackrel{~}{}_i`$, $`i=1,\mathrm{},n`$ their fermionic counterpart (therefore $`\xi `$ is not included in the $`\stackrel{~}{}_i`$’s). The $`\stackrel{~}{\mathrm{\Delta }}_i`$’s and the $`\stackrel{~}{}_i`$’s are respectively the generalizations of (4.19) and (4.18) for instanton number $`k`$. After substituting the solutions to the fermionic constraint (3.19) in (4.42), $`[S_F]_k`$ can be written as
$$[S_F]_k=\overline{\stackrel{~}{}}_i^{\dot{A}\alpha }(h_{ij})_\alpha {}_{}{}^{\beta }(\stackrel{~}{}_j)_{\beta \dot{A}}^{},$$
(5.4)
where $`i,j=1,\mathrm{},n`$ and<sup>16</sup><sup>16</sup>16In the following equation we denote by $`h^{}`$ the hermitean conjugate matrix obtained without complex conjugating $`v`$, i.e. treating $`v`$ as real. $`h=h^{}`$; let us also define $`h_{ij}=8\pi ^2\widehat{h}_{ij}`$, for the sake of future convenience. In the $`k=1`$ case one simply has $`\widehat{h}=\overline{v}`$, whereas for $`k=2`$ the explicit expression for $`\widehat{h}_{ij}`$ will be written in (5.15).
The exponential of $`[S_F]_k`$ can now be expanded in powers. Under the integration over the reduced moduli space $`^+\backslash \{x_0\}`$ the only surviving term of the expansion will be the one that, after using (3.56), produces the top form on $`^+\backslash \{x_0\}`$. It is crucial to remark that all the terms containing $`𝒞_{00}`$ do not contribute to the amplitudes since the parameter $`\lambda `$ introduced in (2.33) does not belong to the moduli space. In order to better perform this expansion, let us define
$$\stackrel{~}{}_i=\left(\begin{array}{cc}(\stackrel{~}{}_i)_4+i(\stackrel{~}{}_i)_3& i(\stackrel{~}{}_i)_1+(\stackrel{~}{}_i)_2\\ i(\stackrel{~}{}_i)_1(\stackrel{~}{}_i)_2& (\stackrel{~}{}_i)_4i(\stackrel{~}{}_i)_3\end{array}\right)=\left(\begin{array}{cc}(\eta _i)_1& (\overline{\eta }_i)_2\\ (\eta _i)_2& (\overline{\eta }_i)_1\end{array}\right),$$
(5.5)
where $`(\stackrel{~}{}_i)_\mu `$, $`\mu =1,\mathrm{},4`$ are the Cartesian components of $`\stackrel{~}{}_i`$. (5.4) can then be written as
$`[S_F]_k`$ $`=`$ $`\overline{\eta }_i^\alpha \left[(h_{ij})_\alpha {}_{}{}^{\beta }+ϵ_{\alpha \delta }(h_{ji})_\sigma {}_{}{}^{\delta }ϵ_{}^{\beta \sigma }\right](\eta _j)_\beta =\overline{\eta }_i^\alpha [(hh^{})_{ij}]_\alpha {}_{}{}^{\beta }(\eta _j)_{\beta }^{}`$ (5.6)
$`=`$ $`2\overline{\eta }_i^\alpha (h_{ij})_\alpha {}_{}{}^{\beta }(\eta _j)_{\beta }^{},`$
since $`h`$ is anti–hermitean. In order to recognize the coefficient of the top form, we now explicitly expand $`\mathrm{exp}([S_F]_k)`$. After a little algebra, one finds
$$e^{[S_F]_k}|_{topform}=(32\pi ^2)^{2n}det\widehat{h}\underset{i=1}{\overset{n}{}}\left[(\stackrel{~}{}_i)_1(\stackrel{~}{}_i)_2(\stackrel{~}{}_i)_3(\stackrel{~}{}_i)_4\right],$$
(5.7)
where we used (5.6) and $`\eta _1\eta _2\overline{\eta }_1\overline{\eta }_2=4\stackrel{~}{}_1\stackrel{~}{}_2\stackrel{~}{}_3\stackrel{~}{}_4`$. The coefficient of the top form on the reduced moduli space is then proportional to the determinant of the matrix $`H`$. However, one more ingredient now emerges: the matrix $`K`$ of the change of coordinate basis between $`_i`$ and $`s\mathrm{\Delta }_i`$; recalling (4.4), we conclude in fact that
$$e^{[S_F]_k}|_{topform}=(32\pi ^2)^{2n}det\widehat{h}|detK|s^{4n}\stackrel{~}{\mathrm{\Delta }}.$$
(5.8)
Finally, inserting (5.8) in (5.3) we get
$$<\frac{\text{Tr}\varphi ^2}{8\pi ^2}>_k=k(32\pi ^2)^{2n}_{^+\backslash \{x_0\}}s^{4n}\stackrel{~}{\mathrm{\Delta }}det\widehat{h}|detK|e^{\left[S_B\right]_k}.$$
(5.9)
where $`[S_B]_k`$ is written in (4.41). Note that in (5.8), (5.9) the instanton integration measure $`|detK|`$ has naturally come out.
This is our starting point. Let us now perform the $`k=1`$ computation explicitly.
### 5.1 The $`k=1`$ Case in the Bulk
In the $`k=1`$ case, the structure of the instanton moduli space $`^+`$ has been thoroughly investigated, and it is explicitly known to be the manifold $`\text{IR}^4\times \text{IR}^+\times S^3/\text{}_2`$ ; the three factors correspond respectively to the instanton center ($`x_0`$), scale ($`|w|`$) and orientation in color space ($`w/|w|`$). The reduced moduli space $`^+\backslash \{x_0\}`$ is then the 4–dimensional manifold $`\text{IR}^+\times S^3/\text{}_2`$, obtained after first integrating out the collective coordinate $`x_0`$ .
The ADHM bosonic and fermionic matrices are written in (4.6), and the action (4.40) calculated on the one–instanton background is given by
$$\left[S_{\mathrm{inst}}\right]_{k=1}=4\pi ^2\left[|v|^2|w|^22\text{Tr}(\overline{v}\mu \overline{\mu })\right].$$
(5.10)
Taking into account that
$$\text{Tr}(\mu \overline{\mu }\overline{v})=sw_\mu sw_\nu \underset{a=1}{\overset{3}{}}\eta _{\mu \nu }^a(v^a)^{},$$
(5.11)
where $`\eta _{\mu \nu }^a`$ are the ’t Hooft symbols, we get, after a little algebra,
$`<{\displaystyle \frac{\text{Tr}\varphi ^2}{8\pi ^2}}>_{k=1}`$ $`=`$ $`{\displaystyle _{^+\backslash \{x_0\}}}e^{[S_{\mathrm{inst}}]_{k=1}}={\displaystyle _{^+\backslash \{x_0\}}}e^{4\pi ^2\left[|v|^2|w|^22\mathrm{T}\mathrm{r}(\mu \overline{\mu }\overline{v})\right]}`$ (5.12)
$`=`$ $`{\displaystyle \frac{(8\pi ^2)^2}{2!}}{\displaystyle _{^+\backslash \{x_0\}}}e^{4\pi ^2|v|^2|w|^2}\text{Tr}(\mu \overline{\mu }\overline{v})\text{Tr}(\mu \overline{\mu }\overline{v})`$
$`=`$ $`(8\pi ^2)^22^2{\displaystyle _{^+\backslash \{x_0\}}}(v^{})^2e^{4\pi ^2|v|^2|w|^2}s^4w={\displaystyle \frac{8\pi ^2}{v^2}},`$
which is the expected result . In (5.10) and hereafter we use the shorthand notation $`\overline{v}=\overline{𝒜}_{00}=v^{}\sigma ^3/2i`$.
### 5.2 The $`k=2`$ Case in the Bulk
In this subsection we will describe the $`k=2`$ computation in the bulk i.e. without using the property that the action is BRST exact. All the features of the topological approach will now become apparent.
The action (4.40) calculated on the two–instanton background is given by
$`\left[S_{\mathrm{inst}}\right]_{k=2}`$ $`=`$ $`[S_B+S_F]_{k=2}=4\pi ^2|v|^2(|w_1|^2+|w_2|^2)16\pi ^2{\displaystyle \frac{|\omega |^2}{H}}`$
$`+`$ $`8\pi ^2\text{Tr}\left\{\overline{\mu }_1\overline{v}\mu _1+\overline{\mu }_2\overline{v}\mu _2+{\displaystyle \frac{\overline{\omega }}{H}}\left[\overline{\mu }_1\mu _2\overline{\mu }_2\mu _1+2(\overline{}_3_1\overline{}_1_3)\right]\right\},`$
where we have defined
$$\omega =\mathrm{\Lambda }_{b12}(𝒜_{00})=\overline{w}_1𝒜_{00}w_2+\overline{w}_2𝒜_{00}w_1.$$
(5.14)
After substituting the fermionic constraint (4.11) in (5.2), $`[S_F]_{k=2}`$ can be written as in (5.4); the indices $`i,j`$ run from 1 to 3, the $`\stackrel{~}{}_i`$’s are defined in (4.18), and $`h_{ij}=8\pi ^2\widehat{h}_{ij}`$, where explicitly
$$\widehat{h}_{ij}=\left(\begin{array}{ccc}\overline{v}& \frac{\overline{\omega }}{H}& \frac{\overline{\omega }}{H|a_3|^2}w_2\overline{a}_3\\ \\ \frac{\overline{\omega }}{H}& \overline{v}& \frac{\overline{\omega }}{H|a_3|^2}w_1\overline{a}_3\\ \\ \frac{\overline{\omega }}{H|a_3|^2}a_3\overline{w}_2& \frac{\overline{\omega }}{H|a_3|^2}a_3\overline{w}_1& \frac{2\overline{\omega }}{H|a_3|^2}(a_3\overline{a}_1a_1\overline{a}_3)\end{array}\right).$$
(5.15)
Moreover, specializing (5.8) to the $`k=2`$ case we get
$$e^{[S_F]_{k=2}}|_{topform}=(32\pi ^2)^6det\widehat{h}|detK|s^4w_1s^4w_2s^4a_3.$$
(5.16)
The determinant of the matrix $`K`$ in (4.21) was explicitly computed in (4.27). We want now to calculate the determinant of $`\widehat{h}`$. To this end, note that this matrix has the form<sup>17</sup><sup>17</sup>17In the following equation we denote by a bar over a quaternion the hermitean conjugate quaternion obtained without complex conjugating $`v`$; in other words, if $`q\text{}`$ and $`v\text{}`$, then we define $`\overline{vq}=v\overline{q}`$.
$$\widehat{h}=\left(\begin{array}{ccc}F& B& C\\ B& F& D\\ \overline{C}& \overline{D}& E\end{array}\right),$$
(5.17)
where $`\overline{F}=F`$ and $`B=\overline{B}`$. By means of elementary operations on the rows and columns of the matrix $`\widehat{h}`$ (i.e. by multiplying rows by quaternions and then adding and subtracting rows) we can write
$$\widehat{h}=h_1h_2=\left(\begin{array}{ccc}\frac{F}{|F|^2|B|^2}& \frac{F}{|B|^2}(\frac{1}{|F|^2}\frac{\alpha \beta ^1}{|C|^2})& \frac{FC}{|C|^2}\\ 0& \frac{B\alpha \beta ^1}{|B|^2|C|^2}& \frac{BC}{|C|^2}\\ 0& 0& 1\end{array}\right)\left(\begin{array}{ccc}0& 0& \gamma \beta \alpha ^1\delta \\ 0& \beta & \beta \alpha ^1\delta \\ \overline{C}& \overline{D}& E\end{array}\right),$$
(5.18)
where
$`\alpha `$ $`=`$ $`|C|^2\overline{B}F+|B|^2C\overline{D},`$
$`\beta `$ $`=`$ $`|B|^2\overline{F}B+|F|^2\overline{B}F,`$
$`\gamma `$ $`=`$ $`|B|^2\overline{F}C+|F|^2\overline{B}D,`$
$`\delta `$ $`=`$ $`|C|^2\overline{B}D|B|^2CE.`$ (5.19)
Using (5.18) one finds, after some algebra,
$$det\widehat{h}=\left(\frac{\overline{\omega }v^{}}{2H}\right)^2\frac{1}{|a_3|^4}det\left\{\frac{v^{}}{2}(\overline{w}_1w_2\overline{w}_2w_1)\frac{i\overline{\omega }}{H}(\overline{w}_1\sigma ^3w_1+\overline{w}_2\sigma ^3w_2)\right\},$$
(5.20)
where $`\sigma ^3`$ is the third Pauli matrix. (5.20) is the determinant of a quaternion, i.e. the squared absolute value of the quaternion itself. The final result is
$$det\widehat{h}=\left(\frac{\overline{\omega }v^{}}{2H}\right)^2\frac{1}{|a_3|^4}\left\{2\overline{\omega }^2\tau _1+\left(\frac{v^{}}{2}\right)^2|\mathrm{\Omega }|^2+\overline{\omega }^2\left[\tau _2+\frac{1}{(v^{}/2)^2}\left(\frac{\overline{\omega }}{H}\right)^2\right]\right\},$$
(5.21)
where
$`\mathrm{\Omega }`$ $`=`$ $`w_1\overline{w}_2w_2\overline{w}_1,`$
$`L`$ $`=`$ $`|w_1|^2+|w_2|^2,`$
$`\tau _1`$ $`=`$ $`{\displaystyle \frac{L}{H}},`$
$`\tau _2`$ $`=`$ $`{\displaystyle \frac{L^2|\mathrm{\Omega }|^2}{H^2}}.`$ (5.22)
(5.21) reproduces the result known in literature .
With the aid of (5.16) we can now compute
$$<\frac{\text{Tr}\varphi ^2}{8\pi ^2}>_{k=2}=2(32\pi ^2)^6_{^+\backslash \{x_0\}}s^4w_1s^4w_2s^4a_3det\widehat{h}|detK|e^{[S_B]_{k=2}}.$$
(5.23)
Using (4.27) and (5.21), it is easy to see that the integral over the bosonic moduli which appears in (5.23) is the same which was found in after integrating out the fermionic zero–modes. As in , (5.23) thus leads to $`\text{Tr}\varphi ^2/(8\pi ^2)_{k=2}=5(8\pi ^2)^3/(4v^6)`$ , which agrees with the results found by Seiberg and Witten in .
### 5.3 On the Use of the Operator $`s`$
As we have observed in the previous sections, the operator $`s`$ is nilpotent. Moreover, it is possible to write the action $`S_{\mathrm{inst}}`$ as the operator $`s`$ acting on a certain function of the moduli as in (4.37). This enables one to write the correlator $`<\text{Tr}\varphi ^2>_k`$ as an integral over the boundary of the instanton moduli space. Since
$$\left[S_{\mathrm{inst}}\right]_k=\left[S_B+S_F\right]_k=4\pi ^2s\left\{\text{Tr}\left[\overline{v}(\underset{i=1}{\overset{k}{}}\mu _i\overline{w}_iw_i\overline{\mu }_i)\right]\right\},$$
(5.24)
and $`s[S_{\mathrm{inst}}]_k=0`$, we obtain
$`e^{[S_{\mathrm{inst}}]_k}`$ $`=`$ $`{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{()^l}{l!}}([S_B]_k+[S_F]_k)([S_B]_k+[S_F]_k)^{l1}`$
$`=`$ $`4\pi ^2s\left\{\text{Tr}\left[\overline{v}({\displaystyle \underset{i=1}{\overset{k}{}}}\mu _i\overline{w}_iw_i\overline{\mu }_i)\right]{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{()^l}{l!}}([S_B]_k+[S_F]_k)^{l1}\right\}`$
$`=`$ $`4\pi ^2s\left\{\text{Tr}\left[\overline{v}({\displaystyle \underset{i=1}{\overset{k}{}}}\mu _i\overline{w}_iw_i\overline{\mu }_i)\right]{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{()^l}{l!}}{\displaystyle \underset{p=0}{\overset{l1}{}}}\left(\begin{array}{c}l1\\ p\end{array}\right)([S_B]_k)^{l1p}([S_F]_k)^p\right\}.`$
Since $`[S_F]_k`$ is a fermion bilinear, in order to build a fermionic top form on the $`(8k4)`$–dimensional reduced moduli space we must have $`p=4k3`$, leading to
$`e^{[S_{\mathrm{inst}}]_k}|_{topform}`$ $`=`$ $`4\pi ^2s\left\{\text{Tr}\left[\overline{v}({\displaystyle \underset{i=1}{\overset{k}{}}}\mu _i\overline{w}_iw_i\overline{\mu }_i)\right]{\displaystyle \frac{([S_F]_k)^{4k3}}{(4k3)!}}{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{()^l}{(l+4k2)}}{\displaystyle \frac{([S_B]_k)^l}{l!}}\right\}`$ (5.26)
$`=`$ $`4\pi ^2s\{\text{Tr}\left[\overline{v}({\displaystyle \underset{i=1}{\overset{k}{}}}\mu _i\overline{w}_iw_i\overline{\mu }_i)\right]([S_F]_k)^{4k3}([S_B]_k)^{4k+2}`$
$``$ $`(1e^{[S_B]_k}{\displaystyle \underset{l=0}{\overset{4k3}{}}}{\displaystyle \frac{([S_B]_k)^l}{l!}})\}.`$
As we stated in the introduction, writing the correlator as a total derivative over the moduli space can lead to interesting results. The $`8k`$–dimensional moduli space $`_k`$, can be compactified according to . If we denote this compactification by $`\overline{}_k`$, it is well known that the boundary $`\overline{}_k`$ can be decomposed into a union of lower moduli spaces, so that we can write
$$\overline{}_k=_k\text{IR}^4\times _{k1}S^2\text{IR}^4\times _{k2}\mathrm{}S^k\text{IR}^4$$
(5.27)
where $`S^i\text{IR}^4`$ denotes the $`i^{th}`$ symmetric product of points of $`\text{IR}^4`$. The curvature density in $`S^l\text{IR}^4\times _{kl}`$ is
$$|F_k|^2=|F_{kl}|^2+\underset{i=1}{\overset{l}{}}8\pi ^2\delta (xy_i)$$
(5.28)
where $`y_iS^i\text{IR}^4`$ are the centers of the instanton. We will check (5.28) in the $`k=2`$ case using (3.9). Given
$$\mathrm{\Delta }^{}\mathrm{\Delta }=\left(\begin{array}{cc}|w_1|^2+(x_1x)^2+|a_1|^2& \overline{w}_1w_2+(\overline{x}_1\overline{x})a_1+\overline{a}_1(x_2x)\\ \\ \overline{w}_2w_1+(\overline{x}_2\overline{x})a_1+\overline{a}_1(x_1x)& |w_2|^2+(x_2x)^2+|a_1|^2\end{array}\right),$$
(5.29)
we observe that one part of the boundary is given by $`|w_1|0`$. Using (4.10) we have
$$\underset{|w_1|0}{lim}\mathrm{\Delta }\left(\begin{array}{cc}0_{2\times 1}& \mathrm{\Delta }_{k=1}\\ x_1x& 0_{1\times 1}\end{array}\right),$$
(5.30)
and
$$\underset{|w_1|0}{lim}det\mathrm{\Delta }^{}\mathrm{\Delta }=(x_1x)^2[|w_2|^2+(x_2x)^2].$$
(5.31)
Then
$`\underset{|w_1|0}{lim}\text{Tr}(FF)_{k=2}={\displaystyle \frac{1}{2}}\underset{|w_1|0}{lim}\mathrm{}\mathrm{}\mathrm{log}det(\mathrm{\Delta }^{}\mathrm{\Delta })_{k=2}d^4x`$
$`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{}\mathrm{}\mathrm{log}(x_1x)^2d^4x+\text{Tr}(FF)_{k=1}=\text{Tr}(FF)_{k=1}+8\pi ^2\delta ^4(xx_1).`$
Extending these computations to encompass all boundaries one can check (5.28). We leave the application of these considerations and of (5.26) to cases with $`k>1`$ to future work and here we limit ourselves to a simple check of (5.26) in the $`k=1`$ case.
From the analyses of , it is known that the boundary of the $`k=1`$ moduli space consists of instantons of zero “conformal” size; this means that if we projectively map the Euclidean flat space $`\text{IR}^4`$ onto a four sphere $`S^4`$, the boundary of the corresponding transformed $`k=1`$ instanton moduli space is given by instantons of zero conformal size $`\tau `$, where $`\tau `$ is obtained from $`|w|`$ through a projective transformation ($`|w|`$ itself does not represent a globally defined coordinate on the $`S^4`$ instanton moduli space). In terms of the size $`|w|`$ of the $`\text{IR}^4`$ instanton, the $`\tau 0`$ limit corresponds to $`|w|0,\mathrm{}`$. Specializing (5.26) to $`k=1`$ and inserting it in (5.3) we get
$`<{\displaystyle \frac{\text{Tr}\varphi ^2}{8\pi ^2}}>_{k=1}`$ $`=`$ $`{\displaystyle _{^+\backslash \{x_0\}}}e^{[S_{\mathrm{inst}}]_{k=1}}={\displaystyle _{^+\backslash \{x_0\}}}e^{4\pi ^2\left[|v|^2|w|^22\mathrm{T}\mathrm{r}(\mu \overline{\mu }\overline{v})\right]}`$ (5.33)
$`=`$ $`4\pi ^2{\displaystyle _{^+\backslash \{x_0\}}}s\{\text{Tr}\left[\overline{v}(\mu \overline{w}w\overline{\mu })\right][S_F]_{k=1}{\displaystyle \frac{1}{[S_B]_{k=1}^2}}`$
$`(1e^{[S_B]_{k=1}}[S_B]_{k=1}e^{[S_B]_{k=1}})\}`$
$`=`$ $`32\pi ^4{\displaystyle _{^+\backslash \{x_0\}}}s\{\text{Tr}\left[\overline{v}(\mu \overline{w}w\overline{\mu })\right]\text{Tr}(\mu \overline{\mu }\overline{v}){\displaystyle \frac{1}{(4\pi ^2|v|^2|w|^2)^2}}`$
$`(1e^{4\pi ^2|v|^2|w|^2}4\pi ^2|v|^2|w|^2e^{4\pi ^2|v|^2|w|^2})\},`$
where we have used (5.26) with
$`\left[S_B\right]_{k=1}`$ $`=`$ $`4\pi ^2|v|^2|w|^2,`$ (5.34)
$`\left[S_F\right]_{k=1}`$ $`=`$ $`8\pi ^2\text{Tr}(\mu \overline{\mu }\overline{v}).`$ (5.35)
Using Stokes’ theorem, we can compute $`\text{Tr}\varphi ^2/(8\pi ^2)_{k=1}`$ as an integral over the boundary $`\left(^+\backslash \{x_0\}\right)`$, which for $`k=1`$ is $`\text{IR}^+\times S^3/\text{}_2`$. Since
$`\text{Tr}(\mu \overline{\mu }\overline{v})`$ $`=`$ $`sw_\mu sw_\nu \eta _{\mu \nu }^a(v^a)^{}=(\sigma _w^as|w|^2|w|^2s\sigma _w^a)(v^a)^{},`$
$`\text{Tr}[\overline{v}(\mu \overline{w}w\overline{\mu })]`$ $`=`$ $`2w_\mu sw_\nu \eta _{\mu \nu }^a(v^a)^{}=2|w|^2\sigma _w^a(v^a)^{},`$ (5.36)
we get
$$\text{Tr}(\mu \overline{\mu }\overline{v})\text{Tr}[\overline{v}(\mu \overline{w}w\overline{\mu })]=4|w|^4(v^av^a)^{}\sigma _w^1\sigma _w^2\sigma _w^3.$$
(5.37)
Here $`\sigma _w^a`$ are the left–invariant 1–forms, defined as $`\sigma _w^a=|w|^2\eta _{\mu \nu }^aw_\mu sw_\nu `$, and satisfy the relation $`\sigma _w^a\sigma _w^b=ϵ^{abc}s\sigma _w^c`$. Plugging (5.37) into (5.33) and recalling that $`_{S^3/\text{}_2}\sigma _w^1\sigma _w^2\sigma _w^3=\pi ^2`$, we get
$`<{\displaystyle \frac{\text{Tr}\varphi ^2}{8\pi ^2}}>_{k=1}`$ $`=`$ $`{\displaystyle \frac{8\pi ^2}{v^2}}\left(1e^{4\pi ^2|v|^2|w|^2}4\pi ^2|v|^2|w|^2e^{4\pi ^2|v|^2|w|^2}\right)|_{|w|=0}^{|w|=\mathrm{}}`$ (5.38)
$`=`$ $`{\displaystyle \frac{8\pi ^2}{v^2}}`$
which is the result obtained in (5.12).
## 6 Topological Correlators in Witten’s Topological Field Theory
In this section we focus the attention on Witten’s twisted formulation of $`N=2`$ SYM theory. We will put to zero the v.e.v. $`v`$ of the complex scalar field. For winding number $`k=1`$, the top form on the (8–dimensional) instanton moduli space is $`\text{Tr}\varphi ^2(x_1)\text{Tr}\varphi ^2(x_2)`$, and one can compute the Green’s function $`<\text{Tr}\varphi ^2(x_1)\text{Tr}\varphi ^2(x_2)>`$. The prescription (4.1) for computing Green’s function gives in this case
$$\text{Tr}\varphi ^2(x_1)\text{Tr}\varphi ^2(x_2)=_^+\left[\text{Tr}\varphi ^2(x_1)\text{Tr}\varphi ^2(x_2)\right]_{zeromodesubspace},$$
(6.1)
where we have recalled that the boundary term (2.49) in $`S_{\mathrm{TYM}}`$ vanishes when $`v=0`$.
We could then proceed and compute explicitly the r.h.s. of (6.1). However, the observation that the BRST operator $`s`$ is on $`^+`$ the exterior derivative leads us to consider, as in subsec. 5.3, the possibility of computing correlators of $`s`$–exact operators as integrals of forms on the boundary of $`^+`$. Indeed, recall that we can write
$$\text{Tr}\varphi ^2=sK_c,K_c=\text{Tr}(csc+\frac{2}{3}ccc),$$
(6.2)
an expression which parallels the well–known relation
$$\text{Tr}F^2=dK_A,K_A=\text{Tr}(AdA+\frac{2}{3}AAA).$$
(6.3)
Using Stokes’ theorem, one is led to re–express the r.h.s. of (6.1) as
$$_^+\text{Tr}\varphi ^2\text{Tr}\varphi ^2=_^+K_c\text{Tr}\varphi ^2.$$
(6.4)
We are then faced with two different computational strategies:
* the bulk calculation, and
* the boundary calculation.
Let us explore in detail both possibilities.
### 6.1 The Calculation of $`<\text{Tr}\varphi ^2(x_1)\text{Tr}\varphi ^2(x_2)>`$ in the Bulk
From the last equation in (3.2) we know that the zero–mode configuration for $`\varphi `$ is
$$\varphi =U^{}f^{}U+U^{}𝒜U.$$
(6.5)
The parametrization for a $`k=1`$ instanton has been described in sec. 4.1; from this it turns out that $`f(x)=(\mathrm{\Delta }^{}\mathrm{\Delta })^1=[(xx_0)^2+w^2]^1`$. Plugging the expression (4.7) for $``$ into (6.5) and recalling that $`𝒜=0`$ when $`k=1`$, we get
$$\varphi =U^{}(s\mathrm{\Delta })f(s\mathrm{\Delta })^{}U.$$
(6.6)
It then follows that
$$\text{Tr}\varphi ^2=\text{Tr}\left[Ps\mathrm{\Delta }f(s\mathrm{\Delta })^{}Ps\mathrm{\Delta }f(s\mathrm{\Delta })^{}\right],$$
(6.7)
where $`P`$ has been introduced in (3.11). After a little algebra, (6.7) becomes
$$\text{Tr}\varphi ^2(x)=48w^4f^4(x)\underset{\mu =1}{\overset{4}{}}\mathrm{\Gamma }_\mu (x),$$
(6.8)
where the quaternionic 1–form $`\mathrm{\Gamma }(x)`$ is defined by
$$\mathrm{\Gamma }(x)=sx_0+\frac{(xx_0)\overline{w}}{|w|^2}sw.$$
(6.9)
It is easy to convince oneself that
$$\underset{\mu =1}{\overset{4}{}}\mathrm{\Gamma }_\mu (x_1)\underset{\nu =1}{\overset{4}{}}\mathrm{\Gamma }_\nu (x_2)=J(x_1x_2)s^4x_0s^4w,$$
(6.10)
with $`J(x_1x_2)=(x_1x_2)^4/w^4`$; we can then write
$$\text{Tr}\varphi ^2(x_1)\text{Tr}\varphi ^2(x_2)=(48)^2w^4(x_1x_2)^4f^4(x_1)f^4(x_2)s^4x_0s^4w.$$
(6.11)
Plugging this expression into the r.h.s. of (6.1), it follows that
$$<\text{Tr}\varphi ^2(x_1)\text{Tr}\varphi ^2(x_2)>=(48)^2(x_1x_2)^4_^+s^4x_0s^4ww^4f^4(x_1)f^4(x_2).$$
(6.12)
The structure of the $`k=1`$ moduli space has been discussed in subsec. 5.1, where we have learnt that $`_{k=1}^+=\text{IR}^4\times \text{IR}^+\times S^3/\text{}_2`$. (6.12) then becomes
$$_{\text{}^+\times S^3/\text{}_2}s^4ww^4_\text{}^4s^4x_0f^4(x_1)f^4(x_2)=\frac{\pi ^4}{72}\frac{1}{(x_1x_2)^4},$$
(6.13)
from which we finally get
$$<\frac{\text{Tr}\varphi ^2(x_1)}{8\pi ^2}\frac{\text{Tr}\varphi ^2(x_2)}{8\pi ^2}>=\frac{1}{2}.$$
(6.14)
We remark that a hasty analysis would lead to the conclusion that, in the limit $`|x_1x_2|0`$, the Green’s function (6.14) is singular due to the behavior of (6.13). This is contrary to the geometrical interpretation of this correlator as a component of the Chern class of the bundle with curvature (2.11) . With a little more thinking one gets convinced that this singularity is only apparent, as we will show in the next section. In our opinion, this interpretation of the above–computed Green’s function makes it unnatural the application to it of clustering arguments, as recently argued in .
We now turn to describe the same calculation performed on the boundary of instanton moduli space.<sup>18</sup><sup>18</sup>18We thank Gian Carlo Rossi for many fruitful discussions and suggestions on the calculations described in the next section.
### 6.2 The Calculation of $`<\text{Tr}\varphi ^2(x_1)\text{Tr}\varphi ^2(x_2)>`$ on the Boundary of $`^+`$
We start off by considering (6.4), which allows us to write
$$<\text{Tr}\varphi ^2(x_1)\text{Tr}\varphi ^2(x_2)>=_^+K_c(x_1)\text{Tr}\varphi ^2(x_2).$$
(6.15)
The expression of the current $`K_c`$ (which is a 3–form) is a trivial extension of (3.10), and reads, for instanton number $`k`$,
$$K_c=\text{Tr}\left[PsD(sD)^{}D(sD)^{}+\frac{1}{3}(D^{}sD)(D^{}sD)(D^{}sD)\right].$$
(6.16)
For $`k=1`$ one simply has
$$D(x)=f^{\frac{1}{2}}(x)\left(\begin{array}{c}w\\ x_0x\end{array}\right),$$
(6.17)
and after a lengthy algebra, one finds
$`K_c(x)`$ $`=`$ $`2f^3(x)[2w^4(w^2+3y^2)\overline{\sigma }_w^1\overline{\sigma }_w^2\overline{\sigma }_w^3+2y^4(y^2+3w^2)\overline{\sigma }_y^1\overline{\sigma }_y^2\overline{\sigma }_y^3`$
$`+`$ $`2w^2y^2(w^2+y^2)\left({\displaystyle \frac{sy^2}{2y^2}}{\displaystyle \frac{sw^2}{2w^2}}\right)(\overline{\sigma }_y^a\overline{\sigma }_w^a)`$
$`+`$ $`w^2y^2(y^2w^2)s(\overline{\sigma }_y^a\overline{\sigma }_w^a)],`$
where we set $`y=x_0x`$. The right–invariant 1–forms $`\overline{\sigma }_z^a`$ are defined as $`\overline{\sigma }_z^a=|z|^2\overline{\eta }_{\mu \nu }^az_\mu sz_\nu `$, and satisfy the relation $`\overline{\sigma }_z^a\overline{\sigma }_z^b=ϵ^{abc}s\overline{\sigma }_z^c`$. The next step consists in computing the product $`K_c(x_1)\text{Tr}\varphi ^2(x_2)`$. The calculation is greatly simplified if one sets $`x_1=x_2`$. Moreover, one has to take into account only the terms that yield a non–vanishing result when integrated on the boundary of instanton moduli space. If we do this, we get
$$[K_c\text{Tr}\varphi ^2](x)192y^4w^4f^4(x)(\overline{\sigma }_w^1\overline{\sigma }_w^2\overline{\sigma }_w^3)\left(\frac{sy^2}{2y^2}\frac{sw^2}{2w^2}\right)(\overline{\sigma }_y^1\overline{\sigma }_y^2\overline{\sigma }_y^3).$$
(6.19)
Note that $`y^4(sy^2/2y^2)(\overline{\sigma }_y^1\overline{\sigma }_y^2\overline{\sigma }_y^3)=s^4(x_0x)=s^4x_0`$. It then follows that
$$_^+[K_c\text{Tr}\varphi ^2](x)=192_{S^3/\text{}_2}\overline{\sigma }_w^1\overline{\sigma }_w^2\overline{\sigma }_w^3\underset{|w|0}{lim}_\text{}^4s^4x_0w^4f^4(x).$$
(6.20)
Since
$`\underset{|w|0}{lim}{\displaystyle _\text{}^4}s^4x_0w^4f^4(x)`$ $`=`$ $`\underset{|w|0}{lim}{\displaystyle _\text{}^4}s^4x_0{\displaystyle \frac{w^4}{[w^2+(xx_0)^2]^4}}=C{\displaystyle _\text{}^4}s^4x_0\delta ^{(4)}(xx_0)`$ (6.21)
$`=`$ $`C,`$ (6.23)
where
$$C=_\text{}^4s^4x\frac{1}{(1+x^2)^4}=\frac{\pi ^2}{6},$$
(6.24)
we conclude that
$$_^+[K_c\text{Tr}\varphi ^2](x)=192\pi ^2\frac{\pi ^2}{6}=(8\pi ^2)^2\frac{1}{2},$$
(6.25)
and the final result is
$$<\frac{\text{Tr}\varphi ^2}{8\pi ^2}\frac{\text{Tr}\varphi ^2}{8\pi ^2}>=_^+\frac{K_c}{8\pi ^2}\frac{\text{Tr}\varphi ^2}{8\pi ^2}=\frac{1}{2}.$$
(6.26)
(6.26) coincides with the result found in (6.14). The limit of coincident points is thus well–defined, as we observed at the end of the previous section.
## 7 The ADHM Construction and Hyperkähler Quotients
In this section we construct the moduli space of self–dual connections on flat space $`\text{IR}^4`$ in terms of hyperkähler quotients following . This will allow us to clarify the geometrical meaning of the algebraic construction of the BRST transformations presented in sec. 3.4. As explained in , the hyperkähler construction of a quotient space can be regarded from a physicist’s point of view as the gauging of a non–linear sigma model. The corresponding connection is obtained in a purely geometrical way directly from the isometries of the constraint equation (3.7) which imposes the self–duality of the gauge field strength (expressed in the ADHM formalism), and coincides with the connection $`𝒞`$ introduced in sec.3.4 and worked out explicitly in sec. 4.1 for the $`k=2`$ case. We will show that the square root of the determinant of the metric on instanton moduli space gives the bosonic Jacobian involved in the transformation of the functional integral into an integration over instanton moduli. For the construction of a gravitational instanton with this method see , while for an introduction to hyperkähler quotients in physicists’ language see .
The starting point is the ADHM matrix $`a`$, which for the case of $`SU(2)`$ instantons was defined in (3.2). Actually, for the present discussion it is more convenient to adopt a different parametrization for $`a`$; we rewrite it as
$$a=\left(\begin{array}{cc}t& s^{}\\ A& B^{}\\ B& A^{}\end{array}\right),$$
(7.1)
where $`A,B`$ are $`k\times k`$ complex matrices and $`s,t`$ are $`N\times k`$ and $`k\times N`$ dimensional matrices. Let us introduce the $`4k^2+4kN`$–dimensional hyperkähler manifold $`M=\{A,B,s,t\}`$. Given the three complex structures $`J_{ab}^i`$ where $`i=1,2,3`$ and $`a,b=1,\mathrm{},\mathrm{dim}M`$, we can build the 2–forms $`\omega ^i=J_{ab}^idx^adx^b`$, where $`x`$ is a choice of coordinates on $`M`$. The real forms $`\omega ^i`$ allow one to define a $`(2,0)`$ and a $`(1,1)`$ form
$`\omega _{\text{}}`$ $`=`$ $`\text{Tr}dAdB+\text{Tr}dsdt,`$
$`\omega _{\text{}}`$ $`=`$ $`\text{Tr}dAdA^{}+\text{Tr}dBdB^{}+\text{Tr}dsds^{}\text{Tr}dt^{}dt.`$ (7.2)
The transformations
$`A`$ $``$ $`QAQ^{},`$
$`B`$ $``$ $`QBQ^{},`$
$`s`$ $``$ $`QsR^{},`$
$`t`$ $``$ $`RtQ^{},`$ (7.3)
with $`QU(k),RU(N)`$ leave $`\omega _{\text{}},\omega _{\text{}}`$ invariant, and are the analogous of (3.12). Let $`\xi `$ be a generator of the algebra which leaves $`\omega ^i`$ invariant,
$$L_\xi \omega ^i=0,$$
(7.4)
where $`L_\xi `$ is the Lie derivative along $`\xi `$. As $`\omega ^i`$ is Kähler, (7.4) gives rise to conserved quantities, called momentum maps, defined as
$$i(\xi )\omega ^i=d\mu _\xi ^i,$$
(7.5)
where $`\mu _\xi ^i=\mu _a^i\xi ^a`$; in complex notation
$`\mu _{\text{}}`$ $`=`$ $`[A,B]+st,`$
$`\mu _{\text{}}`$ $`=`$ $`[A,A^{}]+[B,B^{}]+ss^{}t^{}t.`$ (7.6)
$`\mu _\xi ^i=0`$ defines a hypersurface
$$𝒩^+=\{\{A,B,s,t\}=xM:\mu _\xi ^i=0\}$$
(7.7)
of dimension $`\mathrm{dim}𝒩^+=k^2+4kN`$; using (7.1), one can immediately see that these equations are the equivalent of (3.15). The moduli space of self–dual gauge connections, $`^+`$, is obtained by modding $`𝒩^+`$ by the reparametrizations defined in (7). It has dimension $`\mathrm{dim}^+=4kN`$ and it is hyperkähler.<sup>19</sup><sup>19</sup>19The metric on $`^+`$ could also be obtained from the Kähler form $`\omega _^+`$, which in turn is expressed in terms of the Kähler potential $`𝒦`$, as $`\omega _^+=\overline{}𝒦=\frac{1}{2}\overline{}\text{Tr}\left[a^{}(1+P_{\mathrm{}})a\right]`$, where $`P_{\mathrm{}}=1bb^{}`$ is the asymptotic expression of the projector $`P`$ defined in (3.11).
In the following, we will focus on the $`k=2`$ case with gauge group $`SU(2)`$. For the explicit computations we go back to the parametrization of the ADHM moduli space introduced in sec. 3.1 and exploited for the $`k=2`$ case in sec. 4.1; the matrix $`a`$ is written in (4.8). We introduce a 20–dimensional hyperkähler manifold $`M=(w_1,w_2,a_3,a_1,x_0)`$.<sup>20</sup><sup>20</sup>20Notice that, since we are using a different parametrization of the ADHM space with respect to (7.1), the dimension of the manifolds $`M`$ and $`𝒩^+`$ is not that of the previous discussion. However, also the reparametrization groups are different, in such a way that the final dimension of the moduli space of self–dual gauge connections is the same, as it must be. Actually, since the theory is invariant under the group of translations in $`\text{IR}^4`$, one can fix $`x_0`$ and restrict the analysis to the 16–dimensional hyperkähler manifold $`M\backslash \{x_0\}`$ parametrized by the quaternionic coordinates $`m^I=(w_1,w_2,a_3,a_1)`$, endowed with a flat metric
$$ds^2=\eta _{I\overline{J}}dm^Id\overline{m}^{\overline{J}}=|dw_1|^2+|dw_2|^2+|da_3|^2+|da_1|^2.$$
(7.8)
To keep the notation as simple as possible, we rename $`^+\backslash \{x_0\}`$ and $`𝒩^+\backslash \{x_0\}`$ as $`^+`$, $`𝒩^+`$ respectively.
For $`k=2`$, the ADHM bosonic constraint (3.15) reads
$$\overline{w}_2w_1\overline{w}_1w_2=2(\overline{a}_3a_1\overline{a}_1a_3),$$
(7.9)
and, as discussed in sec. 3.1, it is invariant under the reparametrization group $`O(2)`$, whose action on the $`k=2`$ quaternionic coordinates is
$`(w_1^\theta ,w_2^\theta )=(w_1,w_2)R_\theta ,`$
$`(a_3^\theta ,a_1^\theta )=(a_3,a_1)R_{2\theta },`$ (7.10)
with
$$R_\theta =\left(\begin{array}{cc}\mathrm{cos}\theta & \mathrm{sin}\theta \\ \mathrm{sin}\theta & \mathrm{cos}\theta \end{array}\right).$$
(7.11)
The construction of the reduced bosonic moduli space $`^+`$ proceeds now in two steps. First, we solve explicitly (7.9) in an $`O(2)`$ invariant way. Since the constraint (3.7) corresponds to $`3k(k1)/2`$ equations, $`𝒩^+`$ turns out to be a 13–dimensional manifold for $`k=2`$, described by the set of coordinates $`(w_1,w_2,a_3,\mathrm{\Sigma })`$, where $`\mathrm{\Sigma }`$ is the auxiliary real variable related to the $`O(2)`$ reparametrization symmetry. Second, we mod out this isometry group of $`𝒩^+`$ by means of the hyperkähler quotient procedure. The instanton moduli space is then $`^+=𝒩^+/O(2)`$, and it has dimension $`\mathrm{dim}^+=\mathrm{dim}𝒩^+k(k1)/2|_{k=2}=12`$. As anticipated, the construction of the quotient space $`^+`$ can be seen as the gauging of a non–linear sigma model. The corresponding connection is given by
$$𝒞=\frac{1}{|k|^2}\eta _{I\overline{J}}\left(\overline{k}^{\overline{J}}dm^I+d\overline{m}^{\overline{J}}k^I\right),$$
(7.12)
where $`k^I_I+\overline{k}^{\overline{I}}\overline{}_{\overline{I}}`$ is the $`O(k)`$ Killing vector with $`|k|^2=\eta _{I\overline{J}}k^I\overline{k}^{\overline{J}}`$. The components of the $`O(2)`$ Killing vector on $`M`$ are easily deduced from (7.10):
$$k^I=(w_2,w_1,2a_1,2a_3).$$
(7.13)
Substituting (7.13) into (7.12), we get
$`𝒞`$ $`=`$ $`{\displaystyle \frac{1}{H}}(\overline{w}_1dw_2\overline{w}_2dw_1+2\overline{a}_3da_12\overline{a}_1da_3+`$ (7.14)
$`+d\overline{w}_2w_1d\overline{w}_1w_2+2d\overline{a}_1a_32d\overline{a}_3a_1).`$
Notice that this is exactly the connection (4.16) obtained in sec. 4.1 by solving the fermionic constraint (3.19). Therefore, this procedure clarifies the geometrical meaning of the connection $`𝒞`$ introduced in sec. 3.1, providing a very simple method to compute it directly from the isometries of the ADHM moduli space, without referring to the constraint equation (3.19).
The metric $`g_{I\overline{J}}^{𝒩^+}`$ on the constrained hypersurface $`𝒩^+`$ is obtained plugging (4.10) into (7.8), and gets simplified if we introduce the variable
$$W=\overline{w}_2w_1.$$
(7.15)
The hypersurface $`𝒩^+`$ is now described by the new set of coordinates $`(w_1,U,V,a_3,\mathrm{\Sigma })`$, where
$`U={\displaystyle \frac{W+\overline{W}}{2}},`$
$`V={\displaystyle \frac{W\overline{W}}{2}},`$ (7.16)
are respectively the real and the imaginary part of $`W`$. The Jacobian factor associated to this change of variables is
$$d^4w_1dUd^3V=|w_1|^4d^4w_1d^4w_2.$$
(7.17)
In the new variables, (7.8) reads
$`ds^2`$ $`=`$ $`\left(1+{\displaystyle \frac{|w_2|^2}{|w_1|^2}}\right)|dw_1|^2+{\displaystyle \frac{dU^2}{|w_1|^2}}+{\displaystyle \frac{|dV|^2}{|w_1|^2}}+`$ (7.18)
$`{\displaystyle \frac{dU}{|w_1|^2}}(\overline{w}_2dw_1+d\overline{w}_1w_2)+{\displaystyle \frac{dV}{|w_1|^2}}(\overline{w}_2dw_1d\overline{w}_1w_2)+`$
$`+|da_3|^2+|da_1|^2,`$
which, inserting (4.10), becomes
$`ds^2`$ $`=`$ $`\left(1+{\displaystyle \frac{|w_2|^2}{|w_1|^2}}\right)|dw_1|^2+{\displaystyle \frac{dU^2}{|w_1|^2}}+{\displaystyle \frac{|dV|^2}{|w_1|^2}}+`$ (7.19)
$`{\displaystyle \frac{dU}{|w_1|^2}}(\overline{w}_2dw_1+d\overline{w}_1w_2)+{\displaystyle \frac{dV}{|w_1|^2}}(\overline{w}_2dw_1d\overline{w}_1w_2)+`$
$`+\left(1+{\displaystyle \frac{|a_1|^2}{|a_3|^2}}\right)|da_3|^2+{\displaystyle \frac{d\mathrm{\Sigma }^2}{16|a_3|^2}}+{\displaystyle \frac{|dV|^2}{4|a_3|^2}}+`$
$`{\displaystyle \frac{d\mathrm{\Sigma }}{4|a_3|^2}}(\overline{a}_1da_3+d\overline{a}_3a_1){\displaystyle \frac{dV}{2|a_3|^2}}(\overline{a}_1da_3d\overline{a}_3a_1).`$
The r.h.s. of (7.19) can be regarded as the Lagrangian density of a zero–dimensional non–linear sigma model with target space $`𝒩^+`$. In real coordinates $`m^A=(w_1^\mu ,U,V^i,a_3^\mu ,\mathrm{\Sigma })`$, the $`O(2)`$ Killing vector on this manifold has components
$$k^A=(w_2^\mu ,|w_1|^2|w_2|^2,0,2a_1^\mu ,8(|a_3|^2|a_1|^2)).$$
(7.20)
The global $`O(2)`$ symmetry can be promoted to a local one by introducing the connection (7.12), which on $`𝒩^+`$ is written as
$`𝒞`$ $`=`$ $`{\displaystyle \frac{g_{AB}^{𝒩^+}k^B}{H}}dm^A=`$ (7.21)
$`=`$ $`{\displaystyle \frac{1}{H}}\left(2w_2^\mu dw_1^\mu +dU4a_1^\mu da_3^\mu +{\displaystyle \frac{d\mathrm{\Sigma }}{2}}\right),`$
where the metric $`g_{AB}^{𝒩^+}`$ is obtained by rewriting (7.19) in the coordinates $`\{m_A\}`$. Writing $`U`$ in terms of $`w_1,w_2`$ by means of (7.15) and (7.16), the connection (7.21) becomes
$$𝒞=\frac{1}{H}\left(w_1^\mu dw_2^\mu w_2^\mu dw_1^\mu 4a_1^\mu da_3^\mu +\frac{d\mathrm{\Sigma }}{2}\right).$$
(7.22)
From the gauged version of the Lagrangian (7.19) we can read off the metric on $`^+=𝒩^+/O(2)`$ written in the $`\{m^A\}`$ coordinates, namely
$$g_{AB}^^+=g_{AB}^{𝒩^+}\frac{g_{AC}^{𝒩^+}g_{BD}^{𝒩^+}k^Ck^D}{g_{EF}^{𝒩^+}k^Ek^F}.$$
(7.23)
The local $`O(2)`$ isometry allows one to put $`\mathrm{\Sigma }`$ to zero; notice that in this gauge (7.22) leads to the connection (4.17). Finally, by using translational invariance to restore the dependence on $`x_0`$, and taking into account the Jacobian factor (7.17), we write the volume form on the moduli space of self–dual gauge connections with winding number $`k=2`$ as
$$|w_1|^4\sqrt{g_{\mathrm{\Sigma }=0}^^+}d^4w_1d^4w_2d^4a_3d^4x_0=\frac{H}{|a_3|^4}\left||a_3|^2|a_1|^2\right|d^4w_1d^4w_2d^4a_3d^4x_0,$$
(7.24)
which reproduces the well–known result of Osborn .
## Acknowledgements
We are particularly indebted to G.C. Rossi for many valuable discussions over a long time, and for enlightening comments and suggestions on a preliminary version of this paper. We are also grateful to D. Anselmi, C.M. Becchi, P. Di Vecchia, S. Giusto, C. Imbimbo, V.V. Khoze, M. Matone, S.P. Sorella and R. Stora for many stimulating conversations. D.B. was partly supported by the Angelo Della Riccia Foundation.
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# Photometry of Nova V 1493 Aql Based on data collected at the Osservatorio Astrofisico di Catania, stazione M. G. Fracastoro, Serra la Nave (Etna), Italia
## 1 Introduction
V 1493 Aql (= Nova Aql 1999 n. 1) was discovered by Tago (1999) as an 8.8 mag object on two films taken on July 13.56 UT with a 55-mm f/3 camera lens. Nothing was visible on a film taken 4 days before. The precursor of the nova was too faint in quiescence to be recorded by the 1.2 m Palomar Schmidt, which sets the amplitude of the outburst in the $`\delta \mathrm{m}12`$ mag range (Moro et al, 1999). Low-resolution spectra taken on July 14.6 UT by Ayani and Kawabata (1999) showed strong, broad Balmer lines with FWHM about 3400 $`\mathrm{kms}^1`$. Similar values are reported by Tomov et al (1999) on July 15.9 UT who, in addition, remark the presence of FeII emission lines and the lack of absorption components in the Balmer emission lines. The decline of the nova was very fast and already on July 17.02 UT was fainter than 11.0, thus indicating a $`t2`$ of the order of 3 days. On August 3.89 the nova was already at $`V=13.0`$ (Lehky, 1999). We took advantage of an observing run at the Catania Observatory to monitor the following stages in the decline of the nova. The object underwent a secondary outburst which began shortly before our observations. In this letter we report our and other photometric observations of V 1493 Aql and discuss the implication of these data in the context of the classical nova phenomenon.
## 2 Observations and data reduction
From August 5th 1999 to August 15th we observed V 1493 Aql with the 91 cm telescope of the Osservatorio Astrofisico di Catania at the M. G. Fracastoro mountain station on Mt. Etna, equipped with a single channel photometer with $`UBV`$ filters. Standard stars from the lists of Landolt (1983,1992) were observed each night for calibration purposes. The full data, as well as further details on the observations, are available in electronic form.
Each night we observed at least two among four stars which are angularly near to V 1493 Aql (GSC 01048-0098, GSC 01048-01359, HD 230704, HD 178263) for comparison. Unfortunately all the four stars showed variations in our photometry which were above the expected measurement errors. The Tycho Catalogue (Perryman et al 1997) reports a photometric variability for all of them. We made no use of the differential photometry, although all the data is available electronically. Figure 1 shows photometric data from the present work and from the IAU circulars, while Figure 2 shows the Serra La Nave data in more detail. In the time span by our observations V 1493 Aql showed a brightening of about 0.4 mag. Fluctuations from night to night and even within the night were larger than the estimated errors. We searched for a periodicity in the photometric data and may safely exclude the presence of periodicities of 5 days, or shorter. The periodograms relative to $`U`$, $`B`$ and $`V`$ are very similar.
## 3 The photometric behavior during decline
It is well known that each nova has its own individuality and that no two novae are completely alike in their photometric behavior (Payne–Gaposhkin 1957; Mc Laughlin 1960; Duerbeck 1981; Van den Bergh and Younger 1987). Light curves tend to differ from each other immediately after the first phase of near exponential decay from maximum light showing either regular or irregular variations with oscillations, rapid decline (absorption) followed by recovery, re-brightening, and in many cases even a continuous nearly exponential smooth decline. See also Bianchini and Friedjung (1992) for an interpretation of the oscillatory behaviour.
The light curve of V 1493 Aql is characterized, after the first rapid decline to $`V12`$, by some evidence of the onset of small amplitude oscillations with the presence of two local minima at $`V=12.4`$ (JD 381.904, Reszelski 1999) and at $`V=13.0`$ (JD 394.89 Lehky, 1999). Then a conspicuous re-brightening, with some indication of a secondary oscillation takes place, almost in coincidence with our observations. The secondary maximum reaches $`V=11.5`$ on JD $`420`$ (Hanzl, 1999). AAVSO data confirm the general trend near the secondary maximum, but do not show any clear evidence of oscillations in the previous phase. It is possible that the rather large scatter in these data has masked oscillations of small amplitude. If we consider the whole light curve, the brightening looks like a secondary outburst with an increase in $`V1.5`$ mag with respect to the local minimum of $`V=13.0`$. From another perspective the secondary maximum at $`V=11.5`$ falls about 2.5 mag above the curve which corresponds to a smooth decline. The total duration of the secondary outburst is of the order of 40 days.
We have examined several light curves of novae in outburst, looking for similarities with V 1493 Aql and found only a marginal evidence for a similar behaviour. While the presence of more or less coherent oscillations is a characteristic shared by about 15 % of novae (with V 603 Aql and GK Per as prototypes), the long lasting re-brightening “profile” of V 1493 Aql is quite rare, if not unique. To the best of our knowledge, only a few objects bear some resemblance to V 1493 Aql: 1) N Dor 1971a (Van den Bergh & Younger 1987), but the brightening was of much shorter duration and smaller range in magnitude, and $`t2`$ much longer; 2) N Cep 1971 (ibidem) , similar to N Dor 1971 but with some evidence of an another small amplitude brightening; 3) N Her 1963 = V533 Her, the same, with additional small amplitude brightenings; 4) DK Lac (very long $`t2`$) presented many brightening peaks of very short duration at semiregular time intervals; 5) the recurrent nova T CrB showed a secondary maximum that took place about 100 days after the principal eruption.
Also the $`(BV)`$ observations, despite their coverage of the time interval of the nova decline being only partial, indicate a peculiar behaviour. Other novae generally become bluer as they evolve :$`(BV)`$ smoothly decreases indicating that the shrinking in radius of the emitting “photosphere” is accompanied by a gradual increase in temperature. This is generally interpreted as a gradual decrease in the opacity, due to a decrease in density, and the gradual uncovering of the hotter internal layers in the burning shell. Instead, V 1493 Aql is hotter at maximum light ( $`(BV)<0.4`$) than in the following phases: $`(BV)1`$ near the secondary maximum and $`0.7`$ in the following decline. This clearly indicates that the secondary maximum corresponds to a relative increase of the size of the “photosphere” and a decrease in temperature. Only five points, quite well spaced from the beginning of the outburst up to late stages, are available for the R band (Fig. 2), although, none is near the secondary maximum. Remarkably, they define a curve which has a similar slope as the $`V`$ curve. Actually, at a more accurate inspection (V-R) shows a decrease with time, a behavior that contrasts with that generally observed in other novae, e.g. Nova Vul 1984 ( Robb and Scarfe , 1995), where $`(VR)`$ steadily increases with the progress of the outburst.
## 4 Absolute magnitude, reddening and distance
Nova distances can be derived from a comparison between their intrinsic and observed luminosities and therefore require that the latter are corrected for the effect of interstellar extinction. Van den Bergh and Younger (1987) from an exhaustive $`UBV`$ study of novae at maximum, found that after correction for reddening, the intrinsic color $`(BV)_0`$ of novae two magnitudes below maximum, i.e. at time $`t2`$, is $`(BV)_0=0.02\pm 0.04`$, with an rms of 0.12. Comparison with the measured $`(BV)`$ allows to estimate $`E(BV)`$. For V 1493 Aql we have two points near $`t2`$ with $`(BV)=0.24`$ and $`(BV)=0.39`$ respectively. If we assume $`(BV)0.32`$ at $`t2`$ we obtain $`E(BV)0.33\pm 0.1`$ i.e. $`A_V1.04\pm 0.3`$, assuming $`R_V=3.15`$. Tomov et al. (1999) noted the presence of strong NaI D interstellar lines indicative of a large reddening and consistent with the above value.
At such low galactic latitudes ($`b=+2^{}.16`$) the use of the dust maps of Schlegel et al (1998), to estimate reddening is limited because for $`|b|<5^{}`$ the contaminating sources have not been removed from the maps, so that they provide at best an upper limit to the the reddening. The value derived for V 1493 Aql is $`E(BV)=1.67`$, which appears completely inconsistent with the observed colors.. We shall therefore adopt $`E(BV)=0.33`$ and $`A_V=1.04`$.
We use three methods to estimate the absolute magnitude of the nova in outburst:
1. the empirical relation between absolute magnitude at maximum and rate of decline (MMRD, Zwicky 1936; Della Valle & Livio 1995). Taking $`\mathrm{log}t20.5`$, ( which is likely to be an upper limit) from the MMRD relation of Della Valle & Livio (1995), we obtain $`M_V8.97`$. We point out that at very short $`t2`$ , as in our case, the relation is rather flat and produces changes on $`M_V`$ less than 0.1 mag. The observed maximum was at $`V=8.8`$ but we allow for a slightly brighter $`V`$ at maximum ( $`V8.3`$), since it seems reasonable that the real maximum has been missed. Adopting $`A_V=1.04`$ it is straightforward to obtain $`d=17.6\pm 3.2`$Kpc.
2. The absolute magnitude 15 days after maximum light appears to be independent of the speed class (Buscombe & de Vaucoleurs 1955) and is a good standard candle. As absolute magnitude at day 15 we take $`M_V(15)5.47\pm 0.2`$ as the average value of various estimates reported in Warner (1995, p.266). At day 15 we have $`V=12.6`$ implying $`d=25.5\pm 3.5`$ Kpc. We point out that day 15 lies inside the secondary maximum. Had we used the $`V13.6`$ value corresponding with the smooth (unperturbed) decline then $`d=40.4`$ Kpc.
3. Another estimate of the nova distance can be obtained on the basis of the “theoretical” assumption that the nova luminosity at maximum is close to or exceeds the Eddington luminosity ($`L_{edd}`$) for a 1 $`\mathrm{M}_{}`$ object. In this framework fast novae should reach the highest peaks, up to $`4.7\times 10^5L_{}`$ (Warner 1995). If we take $`L_{max}6L_{edd}2\times 10^5L_{}`$ as a representative value for V 1493 Aql, we obtain M$`{}_{bol}{}^{}=8.55`$. Near maximum light novae radiate mostly in the optical and the bolometric correction BC is quite small and close to $`0.2`$. If we take $`M_V=8.35`$ together with $`A_V=1.04`$ and $`V=8.3`$ we obtain $`d=13.2`$Kpc.
It is disturbing that the absolute visual magnitude obtained with the MMRD method ($`M_V=8.95`$) is smaller than the absolute bolometric magnitude obtained in the assumption $`L6L_{edd}(M_{bol}=8.5)`$. We recall however that the intrinsic dispersion in the MMRD curve of Della Valle and Livio (1995) is of the order of $`0.4`$ magnitudes. The two estimates can be reconciled if a bolometric luminosity close to 10 $`L_{edd}`$ is assumed (this gives $`M_V=8.95`$, if $`BC0.2`$ ). We recall that Duerbeck (1981) indicated that very fast novae should radiate at this high rate . Also Livio (1992), on the basis of numerical calculations, suggested $`L_{max}/L_{edd}4.63\times M_{wd}^3`$. Since a rather massive white dwarf ($`1.2`$ $`\mathrm{M}_{}`$) is theoretically required to produce a very fast nova such as V 1493 Aql, then also Livio’s relation points to $`L_{max}/L_{edd}10`$. By these arguments the distance $`d=13.2`$Kpc obtained with method 3. should be considered as a lower limit and the unweighted average, $`d=18.8\pm 3.6`$ Kpc of the three different estimates is rather conservative.
We have tried to estimate the total amount of the “excess” energy associated with the secondary outburst by measuring the area of the region delimited by the observed luminosity curve and the curve corresponding to a smooth, near exponential decline. The total time interval of the secondary brightening is 40 days and the extra energy contained in the bump associated with the secondary outburst is of the order of $`2\times 10^{43}`$ erg . We point out, however, that this value may be affected by a non negligible error because of the uncertainty in the estimate of the curve corresponding to the “smooth” decline.
## 5 Discussion
The very fast character of the light curve of V 1493 Aql associated with the conspicuous re-brightening (peak increase in $`V2.5`$ mag with respect to the smoothed V curve and total duration of about 40 days) is not common in novae. This, together with the fact that V 1493 Aql is hotter at maximum light ($`(BV)<0.4`$) than in the following phases: $`(BV)1`$ near the secondary maximum and $`0.7`$ in the following decline) makes V 1493 Aql quite peculiar among novae. We considered the possibility that the rapid reddening of V 1493 Aql in the early phases might be due to the formation of dust, however we regard this unlikely because dust formation is associated with the presence of a “deep minimum” in the light curve, which is not observed. Moreover dust is generally found in slow novae.
The large distance we have derived requires a comment. Clearly $`d`$ would be significantly reduced if the apparent magnitude at maximum were much brighter than 8.8. If $`V_{max}6.0`$ the distance would be reduced to a more comfortable 6.2 Kpc. However, we consider it very unlikely that a $`V=6`$ object might have escaped the attention of vigilant sky-watchers during a few days. The lack of detection on films taken on 9.9 July UT (Tago 1999), implies that even if the nova managed to escape the attention of the astronomers, it had to decline by $`3`$ mag in an extremely short time. Considering that the observed decline already suggests a $`t2`$ less than 3 days, this would imply a decay by about 5 magnitudes in less than 6 days, thus making of V 1493 Aql the fastest known nova. The photometric and spectroscopic behaviour of V 1493 Aql is that of a nova; the presence of an emission line spectrum near the maximum of the outburst (Ayani & Kawabata 1999; Tomov et al 1999; Lynch et al 1999) is generally associated with the very fast character in the light curve. We are therefore reluctant to accept that this is a unique object which does not obey the empirical MMRD-like relations found for other members of its class.
The three distance estimates are compatible, within 1.7 $`\sigma `$. Although the distances estimated from the magnitude at maximum are uncertain, to the extent that the “true” maximum may have been missed, the fact that assuming $`V_{max}=8.3`$ we obtain distances which are compatible with the estimate from $`V`$ at day 15, suggests that the maximum may not have been missed by more than a few tenths of magnitude.
A distance of $`18.8`$ Kpc places V 1493 Aql at over 14 Kpc from the Galactic centre, thus the star is at the outskirts of our own galaxy. If however the distance is 25.5 Kpc, as derived from the magnitude at day 15, then the distance from the Galactic centre is about 20 Kpc and thus the nova is outside the Galaxy. In this case it would fall in a Local Group galaxy, which could be called the Aquila galaxy. That a Local Group galaxy at such low Galactic latitudes could have gone undetected is not implausible, so the fact that no Local Group galaxy is known in that direction does not allow to rule out this possibility. The existence of the Aquila galaxy may be disproved by number counts and/or radial velocity surveys.
###### Acknowledgements.
We are grateful to G. Carbonaro, A. Di Stefano and M. Puleo for the assistance during the observations. We wish to thank G. Masi for kindly providing an unpublished measure. We acknowledge with thanks use of the data from the AAVSO International Database .
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# The head-tail structure of high-velocity clouds
## 1 Introduction
HVCs – first discovered by Muller et al. (muller (1963)) – are defined as neutral atomic hydrogen clouds with unusual high radial velocities (relative to the local–standard–of–rest frame, LSR) which deviate significantly from a simple galactic rotation model.
Despite 36 years of eager investigations there is no general consensus on the origin and the basic physical parameters of HVCs. This is mainly because HVCs are difficult to detect in emission other than H i 21-cm line radiation. HVCs mostly appear as “pure” neutral atomic hydrogen clouds. Absorption line studies provide information on the ionization state and the metalicity of HVCs. Their results indicate that the bright and very extended HVC complexes consist (at least partly) of processed material, having $``$ 1/3 of the solar abundances (see Wakker & van Woerden wakker-rev (1997) for a recent review). Recent observational results present evidence for emission of ionized atoms associated with HVCs. Tufte et al. (tufte (1998)) for instance presented H$`\alpha `$ emission associated with the HVC complexes M, A and C. Also the Magellanic Stream was found to be bright in the H$`\alpha `$ line emission (Weiner & Williams weiner (1996)). In the soft X-ray regime evidence is presented that the HVC complexes M, C, D and GCN are associated with excess soft X-ray radiation, produced by a plasma of a temperature $`T_{\mathrm{plasma}}=10^{6.2}`$ K (Herbstmeier et al. herbstmeier (1995), Kerp et al. kerp (1999)). Also in the $`\gamma `$-ray regime the detection of excess $`\gamma `$-ray emission is claimed towards the HVC complex M (Blom et al. blom (1997)).
However, the most critical issue of HVC research is the distance uncertainty to the HVCs. Danly et al. (danly (1993)), Keenan et al. (keenan (1995)) and Ryans et al. (ryans (1997)) consistently determined an upper distance limit of z $``$ 3.5 kpc to HVC complex M. The most important step forward is the very recently determined distance bracket of 2.5 $``$ z $``$ 7 kpc towards HVC complex A by van Woerden et al. (van Woerden (1999)). These results clearly indicate that the HVC complexes M and A belong to the Milky Way and its gaseous halo.
Parallel to the growing evidence that a significant fraction of the HVC complexes belong to the Milky Way and its halo, Blitz et al. (blitz (1999)) supported the hypothesis that some HVCs are of extragalactic origin. They argued, that it is reasonable to assume that primordial gas – left over from the formation of the Local Group galaxies – may appear as HVCs. Observational evidence for such a kind of HVC may be found by the detection of the highly ionized high-velocity gas clouds by Sembach et al. (sembach (1999)), because of its very low pressure of about 5 K $`\mathrm{km}\mathrm{s}^1`$.
The Magellanic System is a special case outside both major approaches. The Magellanic Stream (MS) and the Leading Arm (LA) (Putman et al. putman (1998)) both form coherent structures over several tens of degrees having radial velocities in the HVC regime. Their gas represents most likely the debris caused by the tidal interaction of the Magellanic Clouds with the Galaxy at distances of tens of kpc.
The physical conditions of the HVCs located in the gaseous Galactic halo or in the intergalactic space as well as their chemical compositions should be significantly different. HVCs located in the intergalactic space are only exposed to the extragalactic radiation field. It is reasonable to assume that their column density distribution is dominantly modified by gravitational forces of the Local Group galaxies. In contrast, the gaseous distribution of the HVCs located within the neighborhood of the Milky Way is not only modified by gravitational forces but also by the ambient medium in the gaseous halo and by the strong radiation field consisting of stellar UV-photons, soft X-rays and cosmic-rays.
Meyerdierks (meyerdierks (1991)) detected a HVC that appears like a cometary shaped cloud with a central core and an asymmetric envelope of warm neutral atomic hydrogen (the particular HVC is denoted in literature as HVC A2). He interpreted this head-tail structure as the result of an interaction between the HVC and normal galactic gas at lower velocities. Towards HVC complex C Pietz et al. (pietz96 (1996)) discovered the so-called H i “velocity bridges” which seem to connect the HVCs with the normal rotating interstellar medium. The most straight forward interpretation for the existence of such structures is to assume that a fraction of the HVC gas was stripped-off from the main condensation.
In the present paper, we extend the investigations of Meyerdierks (meyerdierks (1991)) and Pietz et al. (pietz96 (1996)) over the entire sky covered by the new Leiden/Dwingeloo H i 21-cm line survey (henceforward abbreviated as LDS, Hartmann & Burton hartmann97 (1997)). For this aim, we investigated the shape and the column density distribution of a complete sample of HVCs to search for distortions in the HVC velocity fields accompanied by column density gradients.
In Sect. 2 we give a brief summary of observational parameters concerning the LDS. In Sect. 3 we present our HVC-sample selection and the characteristic parameters. In Sect. 4 we describe the detection process of head-tail structures in our sample and show the distribution of head-tail structures within the individual HVC complexes. In Sect. 5 we discuss possible implications for the existence of the head-tail structures. In Sect. 6 we summarize our results.
## 2 The Leiden/Dwingeloo survey
The Leiden/Dwingeloo survey of Galactic neutral atomic hydrogen (Hartmann & Burton hartmann97 (1997)) comprises observations of the entire sky north of $`\delta `$=$`30\mathrm{°}`$. It represents an improvement over earlier large-scale H i surveys by an order of magnitude or more in at least one of the principal parameters of sensitivity, spatial coverage or spectral resolution. Most important for our scientific aim is the correction of the LDS for the influence on stray radiation to the H i spectra (Kalberla et al. kalb (1980), Hartmann et al. hartmann96 (1996)). The survey parameters were compiled by Hartmann (hartmann (1994)) and Hartmann & Burton (hartmann97 (1997)). Here we summarize only the major properties important for this work.
The angular resolution of the survey is determined by the beam size of the 25-m Dwingeloo telescope to 36’. The observations were performed on a regular grid with a true-angle lattice spacing of $`0\stackrel{}{.}5`$ in both, $`l`$ and $`b`$. The velocity resolution was set to 1.03 $`\mathrm{km}\mathrm{s}^1`$per channel of the auto-correlator. The effective velocity coverage (measured with respect to the Local Standard of Rest, LSR) covers the range -450 $`\mathrm{km}\mathrm{s}^1`$$`v_{\mathrm{LSR}}`$ 400 $`\mathrm{km}\mathrm{s}^1`$. The characteristic RMS limit of the evaluated brightness-temperature intensities is about 0.07 K. The residual uncertainties, introduced for instance by baseline fitting, are about 0.04 K.
## 3 The HVC-sample
### 3.1 Selection criteria
The aim of the present work is to perform a systematic search for HVCs with a head-tail structure in the entire data base of the LDS (Brüns bruens (1998)). For a statistical analysis we need a well defined and representative sample. Therefore we define three conditions which all have to be fulfilled by the HVCs under consideration.
* We identify H i emission lines as high-velocity H i profiles if their radial velocities are at least $`|v_{\mathrm{LSR}}|`$ 90 $`\mathrm{km}\mathrm{s}^1`$and deviate at least 50 $`\mathrm{km}\mathrm{s}^1`$ from a simple galactic rotation model (the second condition is important for areas near the Galactic Plane).
* The HVC must be traceable within three individual H i spectra, to overcome residual baseline uncertainties and broad but faint residual interference signals. We demand further that the signal is not correlated with the observational grid. The minimum three H i spectra are discarded if they are aligned only in galactic longitude or latitude. Accordingly, the three spectra build up the smallest allowed map of a HVC of interest. This corresponds to cloud-sizes larger 1°.
* The minimum H i column density of the studied HVCs is $`N_\text{}\text{i}=110^{19}`$$`\mathrm{cm}^2`$ . This constraint is introduced to analyze only high-velocity H i profiles with a sufficient signal-to-noise ratio, allowing to study the shape of the H i emission lines in detail.
Applying the three conditions compiled above to the LDS data, we build up a new “bright source catalogue of extended HVCs” for the northern sky offering an angular resolution below $`1\mathrm{°}`$.
### 3.2 General properties of the HVC sample
Fig. 1 shows the distribution of the selected HVCs across the galactic sky. The dashed-dotted line represents the southern declination boundary of the LDS at $`\delta =30\mathrm{°}`$. Each individual marker represents a single HVC selected according to the three criteria mentioned above (Sect. 3.1). In total we identified 252 HVCs. Because of the applied selection criteria, each marker represents a unique line of sight to a HVC. The selected HVCs follow in detail the positional distribution of the well known HVC complexes (Wakker & van Woerden wakker-rev (1997)). The open circles mark the location of HVCs with a head-tail structure (see Sect. 4).
In Fig. 2a the radial velocity (in the local standard-of-rest frame) is plotted versus the galactic longitude. Most of the HVCs have velocities in the range of -200 $`\mathrm{km}\mathrm{s}^1`$$`v_{\mathrm{LSR}}`$ -100 $`\mathrm{km}\mathrm{s}^1`$or +100 $`\mathrm{km}\mathrm{s}^1`$$`v_{\mathrm{LSR}}`$ +200 $`\mathrm{km}\mathrm{s}^1`$. Only HVCs belonging to the Magellanic Stream, the anti-center complex and some clouds near the Galactic Center show radial velocities of v$`{}_{\mathrm{LSR}}{}^{}`$ -200 $`\mathrm{km}\mathrm{s}^1`$. All these more extreme velocity HVCs are located on the southern galactic sky. The absence of a significant number of HVCs with positive radial velocities is because of the limitation to $`\delta 30\mathrm{°}`$ of the LDS data. In particular the Magellanic Cloud System is traceable in H i 21-cm line radiation with positive $`v_{\mathrm{LSR}}`$ velocities (see Putman et al. putman (1998) or Putman putman99 (1999) for a recent overview).
Fig. 2b shows the radial velocity (in the local standard of rest) as a function of the galactic latitude. In this figure the north/south asymmetry is visible, i.e. the HVCs in the northern galactic sky have relatively low radial velocities while the HVCs on the southern galactic sky show in general much higher radial velocities.
Figs. 2c and d show our HVC sample in a different representation: radial velocities are transformed into the galactic standard-of-rest frame.
$`v_{\mathrm{GSR}}=v_{\mathrm{LSR}}+220\mathrm{sin}(l)\mathrm{cos}(b)`$ (1)
In Fig. 2c the radial velocity (v<sub>GSR</sub>) is plotted versus the galactic longitude. The HVC complexes are grouped now to coherent structures, which are even larger than the extent of an individual HVC complex. For example, the HVC complexes M, A and C build up the largest coherent structure in this representation. Most of the HT HVCs belong to this structure. The Magellanic Stream forms a parallel shifted feature. There is one cloud complex that is clearly outside the main distribution: the Smith cloud (l = 38°, b = –13°, Smith smith (1963)) shows radial velocities v$`{}_{\mathrm{GSR}}{}^{}`$ +250 $`\mathrm{km}\mathrm{s}^1`$while all other clouds have v$`{}_{\mathrm{GSR}}{}^{}`$ +75 $`\mathrm{km}\mathrm{s}^1`$. This may be a hint for a different origin of the Smith cloud. Bland-Hawthorn et al. (bland-hawthorn (1998)) claimed an association of the Smith cloud with the Sgr dwarf.
Fig. 2d shows the radial velocity (v<sub>GSR</sub>) as a function of the galactic latitude. The north/south asymmetry is still visible.
Fig. 3 shows histograms of the HVC distribution for the parameters peak column density, $`v_{\mathrm{LSR}}`$ and $`v_{\mathrm{GSR}}`$. The histogram of peak column density versus number distribution of the corresponding HVCs contain information on the studied ensemble. We evaluated the log($`N`$)-log($`S`$) correlation (where $`N`$ is the number of HVCs per flux interval and $`S`$ is the flux) and found a correlation coefficient of -0.91, clearly indicating a linear relation between both quantities. The linear equation is log($`N`$) = $`(0.53\pm \mathrm{\hspace{0.17em}0.07})\times `$ log($`S`$) + $`(3.00\pm \mathrm{\hspace{0.17em}0.06})`$. Within the uncertainties these numbers are equal to the values of Wakker & van Woerden (wakker91 (1991)) on the log($`N`$)-log($`S`$) of the population of positive and negative HVCs distributed across the northern sky. This result demonstrates, that our ensemble of HVCs is a representative flux-limited sample of HVCs.
Fig. 3b shows the number distribution with respect to the radial velocity ($`v_{\mathrm{LSR}}`$). Most of the HVCs have velocities in the intervals centered on $`v_{\mathrm{LSR}}`$ = –150 $`\mathrm{km}\mathrm{s}^1`$and –100 $`\mathrm{km}\mathrm{s}^1`$. Fig. 3c shows the number distribution with respect to the velocity ($`v_{\mathrm{GSR}}`$). The bulk of the studied HVCs reveal low radial velocities with respect to the galactic standard of rest frame. The mean $`v_{\mathrm{GSR}}`$ velocity is negative ($`\overline{v_{\mathrm{HVC}}}`$ = –62.3 $`\mathrm{km}\mathrm{s}^1`$$`\pm `$ 7.5 $`\mathrm{km}\mathrm{s}^1`$), indicating that the majority of HVCs in our sample are moving towards the Galactic Disk. Wakker & van Woerden (wakker91 (1991)) showed that the inclusion of the southern hemisphere H i data on HVCs in their sample does not change this general $`v_{\mathrm{GSR}}`$ velocity behavior.
## 4 Head-Tail structures
Our HVC sample, selected according to the conditions compiled in Sect. 3, provide information on the distribution of the HVCs relative to the observational parameters galactic longitude and latitude, radial velocity ($`v_{\mathrm{LSR}},v_{\mathrm{GSR}}`$) and peak column density.
In the following step we analyze each individual HVC of the sample with respect to the shape of its H i line profiles, for instance for the variation of the mean velocity and the variation of the column density across the HVC extent.
### 4.1 Velocity gradients
We find that about 40% of the HVCs show up with a significant velocity gradient. These velocity gradients are very frequently associated with a column density gradient (20%). The detected velocity gradients are by them-self not a priori indicators for an intrinsic distortion of the HVC velocity field. There are several effects that can produce a velocity gradient.
* The HVCs of our sample are, because of the applied selection criteria, extended objects. Accordingly, the angle between the line of sight and the solar velocity vector varies across the extent of the HVC. This effect can produce a velocity gradient of about 10 $`\mathrm{km}\mathrm{s}^1`$/$`[\mathrm{°}]`$.
* The same kind of velocity gradient is also expected from the HVC velocity vector, because the angle between the HVC velocity vector and the line of sight changes with position, too. This velocity gradient depends on the unknown 3D velocity of a HVC.
* Two or more HVCs may be superposed on the same line of sight with comparable group velocities. The probability for an accidental superposition of two independent HVCs is very low. Nevertheless it is known that HVCs consist of several clumps. The superposition of two clumps that belong to the same HVC has a much higher probability. High angular resolution H i studies may disclose this kind of arrangement.
* Even if there is no evidence for a rotating HVC so far, a rotating HVC would show up with a velocity gradient similar to rotation curves of galaxies. Revealing a red- and blue-shifted extension in the position-velocity diagram.
The effect related to the solar velocity vector is well defined and therefore easy to calculate according to Eq. 1. All other velocity gradients are related to the HVC phenomenon.
### 4.2 Definition of head-tail structures
Our aim was to search for real distortions in the velocity field of individual HVCs, which may be an indicator for an interaction of the HVC with the surrounding interstellar medium. Accordingly, we studied those HVCs which reveal a velocity and a column density gradient simultaneously. This kind of HVC appears like a comet in the position-velocity domain, and justifies the name head-tail HVC (HT HVC). The head is the region with the highest column density of the HVC, while the column density of the tail is in general much lower (by a factor of 2 - 4). Fig. 4 shows six examples of the studied HT HVCs: the HVC shown in Fig. 4a and b belong to the anti-center complex (HVC189-30-205 and HVC192-24-130). In Fig. 4c two HVCs of the Magellanic Stream are displayed. Both show a head-tail structure. The HVC displayed in Fig. 4d (HVC166+56-130) is a special HT HVC, because the HVC appears to be connected with the H i gas in the Galactic Disk; it forms a so called velocity bridge. Up to now it is unknown, whether the HVC gas is really physically connected to the Galactic Disk gas or just an accidental superposition on the same line of sight. Fig. 4e shows a HT HVC where the head has a relatively low intensity. Fig. 4f shows the HT HVC with the highest observed radial velocity in our sample (v<sub>LSR</sub> = –425 $`\mathrm{km}\mathrm{s}^1`$). The velocity resolution was reduced to 3 $`\mathrm{km}\mathrm{s}^1`$in this picture to give an idea of the extent of the tail. This HVC is a very-high-velocity cloud (VHVC) several degrees away from the galaxy M33 (–300 $`\mathrm{km}\mathrm{s}^1`$$``$ v$`{}_{\mathrm{LSR}}{}^{}(\mathrm{M33})`$ –75 $`\mathrm{km}\mathrm{s}^1`$).
### 4.3 Results
In total, we found 45 HVCs associated with a HT structure in our HVC-sample containing 252 individual HVCs. The column density contrast between the head and the tail varies in general between factors 2 – 4. It is a general behavior that the column density maximum, the head, is located at higher radial velocities than the tail. The HT HVCs are extended objects and cover in general several square degrees.
Fig. 1 shows the distribution of the head-tail HVCs across the galactic sky. The HT HVCs are marked by open circles. Obviously all HVC complexes reveal HT HVCs, even in the northern part of the Magellanic Stream HT HVCs are present. Only in HVC complex L no HT HVCs were found. Fig. 2 shows the distribution of the HT HVCs in comparison to the HVC sample in respect to the parameters galactic longitude, latitude and radial velocity ($`v_{\mathrm{LSR}}`$ and $`v_{\mathrm{GSR}}`$).
In Fig. 3 histograms are plotted for the whole HVC sample and the HT HVCs with respect to the parameters peak column density, $`v_{\mathrm{LSR}}`$ and $`v_{\mathrm{GSR}}`$. We discuss these histograms further in the following subsections.
#### 4.3.1 Peak column density
Fig. 3a shows histograms for the number of HVCs per column density interval. The histogram for the entire HVC sample is plotted in light gray, the histogram for the HT HVCs is plotted in dark gray. It is obvious that there are only relatively few HT HVCs at lower column densities (relative to the HVC sample).
A detailed analysis of the histogram for the peak column density revealed a linear relation between the fraction of head-tail HVCs and the peak column density. Fig. 5 shows this relation. The vertical error-bars are calculated with respect to the low statistics (there are very few HT HVCs with a high column density), the horizontal error-bars indicate the size of the individual $`N_\text{}\text{i}`$ intervals. The solid line represents the linear regression fit:
$`{\displaystyle \frac{\mathrm{\#}HTs}{\mathrm{\#}HVCs}}=(0.0317\pm 0.0028){\displaystyle \frac{N(\text{}\text{i})}{[10^{19}\mathrm{cm}^2]}}+`$
$`+(0.017\pm 0.012)`$ (2)
We like to emphasize that the lower limit of $`N_\text{}\text{i}=\mathrm{1\; 10}^{19}`$$`\mathrm{cm}^2`$ was chosen to derive statistical significant information on even the faintest HVCs in our sample (this limit is about a factor of 3 above the detection limit using the highest velocity resolution $`\mathrm{\Delta }`$v = 1.03 $`\mathrm{km}\mathrm{s}^1`$of the LDS data). The correlation displayed in Fig. 5 is not biased by a detection limit.
This empirical relation can be used to determine an expected number of HT HVCs for each individual HVC-complex in our sample. The comparison of the expected numbers of HT HVCs with the detected ones shows a good consistency. Only towards the Galactic Anti-Center a significantly larger number of HT HVCs is detected than expected. The correlation explains why we were not able to find HT HVCs in Complex L: this complex has only very few, faint HVCs (the expected number of HT HVCs for this complex is only 0.3).
#### 4.3.2 Radial velocity ($`v_{\mathrm{LSR}}`$)
Fig. 3b shows histograms for the number of HVCs per velocity ($`v_{\mathrm{LSR}}`$) interval. The histogram for the entire HVC sample is plotted in light gray, the histogram for the HT HVCs is plotted in dark gray.
A large number of HT HVCs have radial velocities in the intervals centered on –100 and –150 $`\mathrm{km}\mathrm{s}^1`$(62.2%). This is consistent with the fact that most HVCs have radial velocities in this regime (62.9%). The interval centered on –200 $`\mathrm{km}\mathrm{s}^1`$contains a relatively large fraction of HT HVCs. The large number is expected because of the fact that there is a high percentage of high column density HVCs in this velocity interval, increasing the probability to detect a HT HVC according to Eq. 2. It remains unclear whether there are too many high column density HVCs or too few low column density HVCs in this velocity interval.
#### 4.3.3 Radial velocity ($`v_{\mathrm{GSR}}`$)
Fig. 3c shows histograms for the number of HVCs per velocity ($`v_{\mathrm{GSR}}`$) interval. The histogram for the entire HVC sample is plotted in light gray, the histogram for the HT HVCs is plotted in dark gray.
The mean radial velocity $`v_{\mathrm{GSR}}`$ of the HT HVCs is more negative than the mean radial velocity of the HVC sample. A Gaussian fit to the histograms revealed mean velocities of $`\overline{v_{\mathrm{HTHVC}}}`$ = –86.1 $`\mathrm{km}\mathrm{s}^1`$$`\pm `$ 3.8 $`\mathrm{km}\mathrm{s}^1`$and $`\overline{v_{\mathrm{allHVC}}}`$ = –62.3 $`\mathrm{km}\mathrm{s}^1`$$`\pm `$ 7.5 $`\mathrm{km}\mathrm{s}^1`$. The difference between the mean velocities ($`\mathrm{\Delta }`$v=23.8$`\mathrm{km}\mathrm{s}^1`$) is significant. The probability to find a HT HVC increases with negative GSR velocity. We checked that this result is not caused by the relation shown in Fig. 5.
The velocity dispersion in Fig. 3c is nearly identical for the HT HVCs (FWHM = 174.0 $`\mathrm{km}\mathrm{s}^1`$$`\pm `$ 7.5 $`\mathrm{km}\mathrm{s}^1`$) and the complete HVC sample (FWHM = 179.9 $`\mathrm{km}\mathrm{s}^1`$$`\pm `$ 15.1 $`\mathrm{km}\mathrm{s}^1`$).
## 5 Discussion
In the following subsections we discuss possible interpretations for HVCs with a head-tail structure at different distances from the Galactic Disk.
### 5.1 HVCs in the gaseous Galactic halo
The gaseous halo of the Milky Way has a vertical scale height of 4.4 kpc (Kalberla & Kerp kalberla (1998)). At least the HVC complexes M (z $``$ 3.5 kpc, Danly et al. danly (1993), Keenan et al. (keenan (1995)) and Ryans et al. (ryans (1997))) and A (2.5 $``$ z $``$ 7 kpc, van Woerden et al. van Woerden (1999)) are located within the gaseous halo. As a natural consequence some kind of interaction is expected when a cloud has a high velocity relative to the ambient medium. In a ram-pressure model the head is the interacting HVC and the tail is gas which was recently stripped off from the HVC. In an absolutely homogeneous halo medium there should be a constant interaction rate for all HVCs with similar distances and velocities. The HVC complexes form coherent structures in position and in velocity, i.e. they may form also coherent structures in space. Therefore one would expect naively that all HVCs of a complex should show up with a head-tail structure. This is obviously not observed.
Kalberla & Kerp kalberla (1998) estimate that $``$ 10 % of the halo gas is neutral, having higher density than the surrounding plasma. This implies that only those HVCs are expected to show up with a head-tail structure that are passing through an area with locally higher halo density. Accordingly there is a 10 % probability for the creation of a HT HVC. Thus at least 10 % of the HVCs should have a head-tail structure.
In the following we estimate the life-time of a head-tail structure. The HVC tails have on the average a column density of a few times $`10^{19}`$ $`\mathrm{cm}^2`$ . Because of the low dust abundance in galactic HVCs (in particular in case of the galactic HVC complex M, Wakker & Boulanger wakker86 (1986)) photoelectric heating can be neglected as an excitation process. Only the diffuse X-ray radiation of the galactic halo and the cosmic-rays heat the HVCs. According to Wolfire et al. (wolfirea (1995)a and b) column densities of a few $`N_\text{}\text{i}10^{19}`$ $`\mathrm{cm}^2`$ have a high ionization probability. This corresponds to a life-time in the neutral state of about $`10^5`$$`10^6`$ years. Compared to the free-fall time of a galactic HVC, of about $`10^7`$ years, the life-time of the stripped-off matter is very short. The tail of a high peak column density HVC contains more material and will survive significantly longer against the ionizing radiation of the diffuse X-rays and cosmic rays. This is consistent with the observational fact that the high peak column density HVCs reveal much more frequently a HT HVC (Fig. 5).
In addition, it is possible that the production or the life-time of a tail depends on the small scale structure of a HVC. HVCs consist of two components: cold clumps are surrounded by an envelope of warm gas. HT HVCs with clumps in their tails will survive much longer than tails that contain only diffuse warm gas. The angular resolution of the LDS is too low to resolve the inner parts of the HVCs. H i observations with high angular resolution are mandatory to study the effects of the small-scale structure.
### 5.2 The Magellanic Stream
The Magellanic Stream (MS) is the only HVC complex with a generally accepted origin: it is build up by gaseous debris caused by the tidal interaction between the Magellanic Clouds and the Milky Way. The distance of the southern end of the Magellanic Stream is assumed to be very similar to the distance to the Magellanic Clouds (D $``$ 50 kpc). In the opposite, the distance of the northern end is unknown. If it is located at low z-heights, i.e. in the gaseous halo, the observed head-tail structures may be produced by the same process as the HT HVCs in the complexes M and A. In this scenario the MS covers a large distance bracket of several tens of kpc. There should be no head-tail structures in the southern part of the stream. On the other hand, the northern part of the MS could have distances similar to the distance of the Magellanic Clouds. The existence of HT HVCs cannot be explained by an interaction with the gaseous halo, because it would be far outside the gaseous halo. A possible interaction partner could be the debris from previous revolutions of the Magellanic System. The ionizing radiation field is much weaker at these distances. Accordingly, the life-time of a tail in the MS is orders of magnitudes longer compared to HVCs in the lower Galactic halo.
### 5.3 HVCs at extragalactic distances
There is evidence that some HVCs are at intergalactic distances (Braun & Burton braun (1999), Blitz et al. blitz (1999), Sembach et al. sembach (1999)).
Sembach et al. (sembach (1999)) found highly ionized high-velocity gas clouds. Several lines of evidence, including very low thermal pressures (P/k $``$ 2 cm<sup>-3</sup> K), favor a location for the highly ionised high-velocity gas clouds in the Local Group or very distant Galactic halo.
Braun & Burton (braun (1999)) searched for compact, isolated HVCs. They classified HVCs as compact if they have angular sizes less than 2 degrees FWHM. They are isolated in that sense, that they are are separated from neighboring H i emission by expanses where no emission is seen to the detection limit of the data. They found 66 of these compact, isolated HVCs. A comparison between their and our sample reveal that the sample of Braun & Burton is much more uniformly distributed in the parameter space (position vs. velocity) than the HVCs of our sample. The different distributions may be an indication for different objects in origin and evolution. Braun & Burton claimed, that their HVCs are most likely located at intergalactic distances. The HVCs of our sample probably contain a significant number of HVCs located in the gaseous halo of the Milky Way. 11 HVCs of their sample are also included in our sample. Moreover, two of them show up with a head-tail structure (HVC271+29+181 and HVC30-51-119).
If these HVCs are located in the intergalactic space, the life-time of the head-tail structures would be much longer because of the much weaker ionizing radiation field at these distances. On the other hand, it is not straight forward to explain what kind of interaction process may produce the observed features at distances of several hundreds of kpc. Especially these (probably very distant) head-tail HVCs should be observed with a higher angular resolution in the near future.
## 6 Summary and conclusion
We performed a systematic analysis of the HT HVCs across the whole sky that is covered by the LDS (about 75% of the entire sky). We selected a representative flux limited HVC sample with column densities $`N_\text{}\text{i}\mathrm{1\; 10}^{19}`$$`\mathrm{cm}^2`$ and minimum angular diameter $``$ 1°. In total we found 252 HVCs.
Each individual HVC was analyzed with respect to velocity gradients and asymmetries in the H i line profiles. The so called head tail HVCs have a cometary shape in the position-velocity domain. We found that 45 out of 252 HVCs show up with a head-tail structure. These HT HVCs are randomly distributed over the whole sky covered by the LDS.
A statistical evaluation of the HVC ensemble revealed that the probability to detect a HT structure increases linear with the peak column density and with increasing negative radial velocity in the GSR frame. There is no significant correlation between the other parameters like galactic longitude, latitude and $`v_{\mathrm{LSR}}`$.
The detection of HT HVCs in nearly all prominent HVC complexes implies qualitatively comparable physical processes in all of the HVC complexes. Individual HVC cores seem to interact with their ambient medium. In case of the HVC complexes located within the Galactic halo and the Magellanic Stream the interaction with the gaseous halo medium or gaseous debris distributed along the orbit of the Magellanic Clouds may serve as a straight forward explanation for the existence of the HT HVC. At intergalactic distances the detection of HT HVCs is difficult to interpret, because the characteristic time-scale for heating and cooling of the HVC matter is orders of magnitude shorter than the assumed age of these HVC complexes.
High angular resolution H i observations of the detected 45 HT HVCs in future will improve our knowledge on the temperature structure and small-scale column density distribution of this special kind of HVCs. In addition the comparison with other wavelengths, e.g. with H$`\alpha `$ emission, will help to understand the existence of HT HVCs.
###### Acknowledgements.
Part of this work was supported by the German *Deutsche Forschungsgemeinschaft, DFG* project number ME 745/19.
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# Phase-coherent effects in multiterminal superconductor/normal metal mesoscopic structures
## Abstract
In this report we analyse three effects which may arise in a mesoscopic multiterminal S/N structure (two normal and two superconducting reservoirs). We show that the Josephson critical current is a non-monotonic function of a voltage between the normal reservoirs. The influence of the spin-polarized electrons on the Josephson critical current was also studied. We show that if there is a temperature difference between the normal reservoirs, a voltage between the superconductors and normal conductors arises which oscillates as the phase difference varies. Its magnitude is much greater than the thermoemf in the case of the ordinary thermoelectric effect.
1. Introduction
In this report we present results of the theoretical study of transport in diffusive 4-terminal superconductor-normal metal (S/N) mesoscopic structures. A great deal of attention was paid recently to the problem of the influence a current between the N reservoirs has on the Josephson effect in the superconducting curcuit. This problem was analysed first in Ref. for a ballistic S/N system and in Ref. for a diffusive system. It was shown that the Josephson critical current $`I_c`$ changes sign at a certain voltage $`V_\text{ }`$between the N reservoirs, and the Josephson junction switches into a $`\pi `$-state ($`\pi `$-junction). The reason for the sign reversal effect is that the distribution function in the N wire, which can have an equilibrium form, is shifted with respect to the distribution function in the superconductors by the value $`V`$. Later the dependence $`I_c(V)`$ was calculated in various models and approximations in Refs. (diffusive case) and (ballistic case). The sign reversal effect in a diffusive Au/Nb mesoscopic structure was observed by Baselmans et al. .
In multi-terminal S/N structures one can observe not only the sign reversal effect, but also a number of other interesting phenomena. It was shown in Refs. that the measured critical current $`I_m`$ in a structure similar to that shown in Fig.1a may not only decrease, but also may increase with increasing voltage $`V`$. In particular one can observe the Josephson-like effects (plateau on the $`I_3(V_S)`$ curve, oscillations of the measured critical current $`I_m`$ in a magnetic field etc) even if the Josephson coupling between superconductors under equilibrium conditions is vanishing. The stimulation of the Josephson effect in this structure is related in particular to a phase dependence of the S/N interface or N wire conductance: $`\delta G\mathrm{cos}\phi .`$ It turns out that apart from the Josephson term $`I_c\mathrm{sin}\phi ,`$ an additional term $`I_{sg}\mathrm{cos}\phi `$ appears in a dynamic equation for $`\phi `$. This term leads to a change of the critical current from the value $`I_c`$ to a new (measured) value $`I_m=\sqrt{I_c^2+I_{sg}^2}`$,where $`I_{sg}`$ is the amplitude of the subgap current. The stimulation of the Josephson effect by a current between the N reservoirs was analysed in Refs. on a simplified model of gapless superconductors. Here we present results of the analysis for the real case where the superconductors have a gap and the density-of-states (DOS) has a singularity . We show that the measured critical current $`I_m`$ reaches a maximum at $`eV\mathrm{\Delta }`$ and depends weakly on the temperature $`T.`$ These results agree with recent observations of the stimulation of the Josephson critical current .
We also study the effect of the spin injection into the N wire from the ferromagnetic (F) reservoirs on the critical current $`I_c`$ for the structure shown in Fig.1b. If the magnetizations in the F reservoirs have the same direction, then there is no spin injection. In the case of antiparallel orientation, the total magnetization of the injected electrons in the N wire is finite. We show that in the case of a N wire short length $`2L`$ ($`L<\{L_ϵ,L_{sp}\}`$; here $`L_ϵ`$ and $`L_{sp}`$ are the energy and spin relaxation length) the spin injection does not affect the critical current and the dependence $`I_c(V)`$ is the same as in the case of normal (nonmagnetic) reservoirs. In a semi-mesoscopic structure $`(L_ϵ<L<L_{sp})`$ the $`\pi `$-junction is realised only in the case of the antiparallel orientation of the magnetizations in the F reservoirs when spins are accumulated in the N wire.
Finally we analyse the thermoelectric effect in the structure shown in Fig.1a. We assume that there is no current between the N reservoirs, but the reservoirs are maintained at different temperatures $`T=T_o\pm \delta T`$. It will be shown that the temperature gradient leads to a voltage difference $`V_{T\text{ }}`$between the normal and superconducting circuits. This voltage does not contain the small parameter $`T/ϵ_F`$ (as is the case for the ordinary thermoeffect in normal metals) and reaches a value of order $`\delta T/e`$. In addition, the voltage $`V_{T\text{ }}`$ depends on the phase difference $`\phi `$; it is zero at $`\phi =0`$ and oscillates with increasing $`\phi .`$
2. Basic equations
The Green’s function technique is perhaps the most convenient and powerful method for studing transport and nonequlibrium phenomena in S/N mesoscopic structures. If we are not interested in purely quantum interference effects arising as the result of multiple (non-Andreev) reflections, then we can use a simpler, well developed method of quasiclassical Green’s functions (see, for example, ). In this method the matrix Keldysh function $`\widehat{G}`$ is introduced along with the retarded (advanced) matrix Green’s functions $`\widehat{G}^{R(A)}=G^{R(A)}\widehat{\tau }_z+`$ $`\widehat{F}^{R(A)}`$, here $`\widehat{F}^{R(A)}`$ is the matrix condensate Green’s function. The functions $`\widehat{G}^{R(A)}`$determine the excitation spectrum and the DOS. The function $`\widehat{G}`$ describes the nonequilibrium properties and is related to the matrix distribution functions $`\widehat{f\text{ }}`$: $`\widehat{G}_\alpha =\widehat{G}_\alpha ^R\widehat{f\text{ }}_\alpha \widehat{f\text{ }}_\alpha \widehat{G}_\alpha ^A`$, where $`\widehat{f\text{ }}_\alpha =f_{\alpha +}\widehat{1}+f_\alpha \widehat{\tau }_z.`$ The components of $`\widehat{f\text{ }}_\alpha `$ can be expressed through the distribution functions of electrons $`n_\alpha `$ and holes $`p_\alpha `$ ($`\alpha `$ is the spin index). The function $`f_{\alpha +}`$ describes an electron and hole distribution symmetrical in branch populations: $`f_{\alpha +}=1(n_\alpha +p_{\overline{\alpha }})`$, here $`\overline{\alpha }=`$ if $`\alpha =`$. The function $`f_\alpha =(n_\alpha p_{\overline{\alpha }})`$ describes the branch imbalance and determines the electric potential.
The equations for $`\widehat{G}_\alpha `$ and $`\widehat{G}_\alpha ^{R(A)}`$ can be solved analitically in the case of a short structure ($`ϵ_{Th}D/L^2>>ϵ`$, where $`ϵ`$ is a characteristic energy; $`ϵ\mathrm{min}\{T,\mathrm{\Delta }\}`$). In the general case one needs to solve equations for the distribution functions $`f_{\alpha \pm }`$ and the Usadel equation . In the present paper we consider the case when the temperature may be both greater or less than the Thouless energy $`ϵ_{Th}`$ and the approximation of a short length is not valid. The distrubution functions in the structure shown in Fig.1a obey the kinetic equation (we consider only the dirty limit)
$$L_x[M_\pm _xf_\pm (x)+J_Sf_{}(x)\pm J_{an}_xf_{}(x)]=r[\stackrel{}{A_\pm }\delta (xL_1)+\stackrel{\mathrm{\_}}{A_\pm }\delta (x+L_1)].$$
(1)
where $`M_\pm =(1G^RG^A(\widehat{F}^R\widehat{F}^A)_1)/2;`$ $`J_{an}=(\widehat{F}^R\widehat{F}^A)_z/2,J_s=(1/2)(\widehat{F}^R_x\widehat{F}^R\widehat{F}^A_x\widehat{F}^A)_z,`$ $`A_\pm =(\nu \nu _S+g_1)(f_\pm f_{S\pm })(g_{z\pm }f_S+g_zf_{});g_{1\pm }=(1/4)[(\widehat{F}^R\pm \widehat{F}^A)(\widehat{F}_S^R\pm \widehat{F}_S^A)]_1;g_{z\pm }=(1/4)[(\widehat{F}^R\widehat{F}^A)(\widehat{F}_S^R\pm \widehat{F}_S^A)]_z;`$
The parameter $`r=R/R_b`$ is the ratio of the resistance of the N wire $`R`$ and S/N interface resistance $`R_b`$; the functions $`\stackrel{\mathrm{\_}}{A_{}}`$ and $`\underset{+}{\overset{\mathrm{\_}}{A}}`$ coincide with $`\stackrel{\mathrm{\_}}{A_{}},\underset{+}{\overset{\mathrm{\_}}{A}}`$ if we make a substitution $`\phi \phi `$. We introduced above the following notations $`(\widehat{F}^R\widehat{F}^A)_1=Tr(\widehat{F}^R\widehat{F}^A),`$ $`(\widehat{F}^R\widehat{F}^A)_z=Tr(\widehat{\tau }_z\widehat{F}^R\widehat{F}^A)`$ etc.; $`\nu ,`$ $`\nu _S`$ are the DOS in the N film at $`x=L_1`$ and in the superconductors. The functions $`f_{S\pm }`$ are the distribution functions in the superconductors which are assumed to have equilibrium forms. This means that $`f_{S+}f_{eq}=\mathrm{tanh}(ϵ\beta )`$ (here $`\beta =1/(2T)`$) and $`f_S=0`$, because we set the potential of the superconductors equal to zero (no branch imbalance in the superconductors). We assumed that the width of the S/N interface $`w`$ is much less than $`L_{1,2}`$ and introduced the $`\delta `$functions. Note that in some papers the last term in the left-hand side of Eq.(1) is missing. One can solve Eq.(1) and express the distribution functions $`f_{\alpha \pm }`$ in terms of the retarded (advanced) Green’s functions which should be found from the Usadel equation. Integrating Eq.(1) once, we obtain
$$M_\pm _xf_\pm (x)+J_Sf_{}(x)\pm J_{an}_xf_{}(x)=J_{1,2\pm }$$
(2)
where the indeces $`1,2`$ relates to the intervals $`(0,L_1)`$,$`(L_1,L)`$ and the constants $`J_{1,2\pm }`$ are partial flows. For example, the partial flow $`J_{1,2}`$ determines the electrical current in the intervals $`(0,L_1)`$,$`(L_1,L)`$
$$I_{1,2}=\sigma _0^{\mathrm{}}𝑑ϵJ_{1,2}$$
(3)
Here $`\sigma `$ is the conductivity of the N wire. As it follows from Eq.(1), the flows $`J_{1,2\pm }`$ are connected with each other
$$J_{2\pm }J_{1\pm }=rA_\pm $$
(4)
The functions $`rA_\pm `$ are the flows through the S/N interface. In particular, $`rA_{}`$ is the partial electrical current through the S/N interface. The first term $`(\nu \nu _S+g_{1+})f_{}`$ in the expression for $`A_{}`$ is the quasiparticle current (the term $`g_{1+}f_{}`$ contributes to the subgap current), and the second term $`(g_{z+}f_{S+}+g_zf_+)`$ is related to the supercurrent. One can show that the partial supercurrent $`J_S`$ is a constant in the interval $`(0,L_1)`$ and equals zero outside this interval.
We can solve Eqs.(2) and express the distribution functions $`f_\pm `$ through the constants $`J_{1,2\pm }.`$ At $`x=\pm L`$ the functions $`f_\pm `$ are equal to equilibrium distribution functions in the N reservoirs. Therefore we can find the functions $`f_\pm (x)`$ and substituting them into Eq.(4) determine $`f_\pm (L_1)`$ and $`J_{1,2\pm }`$. Thus the problem is reduced to solving the Usadel equation. The solution of the Usadel equation can be found either numerically or analitically in limiting cases. In the next sections we will employ this method to study the dependence of the critical current on the voltage between the N reservoirs and the thermoeffect.
3. Suppression and enhancement of the critical current
In this section we consider a symmetrical structure shown in Fig.1a. We assume that the electric potential at the N reservoirs is $`V`$ and at the superconductors is zero (the quasiparticle currents flow from the N reservoirs to the superconductors and the supercurrent flows between the superconductors). The distribution functions in the reservoirs have the equilibrium form: $`F_{V\pm }=[\mathrm{tanh}((ϵ+eV)\beta )\pm \mathrm{tanh}((ϵeV)\beta )]/2.`$
First we consider the case of a small $`r,`$ that is, the interface resistance is large and the condition
$$(ϵ\tau _ϵ)^1<<r<<1$$
(5)
should be satisfied. The first condition ensures that inelastic collision term can be neglected (here $`\tau _ϵ`$ is the energy relaxation time, $`ϵ\mathrm{min}\{T,ϵ_{Th}\}`$). Then the distribution function $`f_+`$ is constant and in the main approximation is $`f_+(F_{V+}+f_{eq}(r_2\nu \nu _s))/(1+r_2\nu \nu _s),`$ here $`r_2=rL_2/L`$. The function $`f_{}`$ is constant in the interval $`(0,L_1)`$ and in the main approximation equals $`f_{}F_V/(1+r_2\nu \nu _s)`$. Outside this interval the function $`f_{}`$ increases linearly to the value $`F_V.`$ Following the method presented above, we find the current through the S/N interface
$$I_3(V)=I_2I_1=I_o(V)+I_{sg}(V)\mathrm{cos}\phi I_c(V)\mathrm{sin}\phi $$
(6)
The first two terms are the quasiparticle current and the last term is the Josephson supercurrent. All the components of the current depend on $`V`$. This expression shows that at a given control voltage $`V`$ and zero voltage difference between superconductors ($`\phi `$ is constant in time) the current $`I_3`$ may vary with changing $`\phi `$ in some limits:$`I_3(V)I_o(V)I_m(V)`$. This means a plateau on the $`V_S(I_3)`$ characteristics (see ); here $`V_S=(\mathrm{}/2e)_t\phi `$ is the voltage difference between superconductors. We can write the phase-dependent part of $`I_3`$ in the form $`I_{3\phi }=I_m\mathrm{sin}(\phi +\alpha )`$, where $`I_m=\sqrt{I_c^2+I_{sg}^2}`$ is the measured critical current, $`\mathrm{cos}\alpha =I_c/I_m`$. In Fig.2 we plot the dependence of the measured critical current $`I_m`$ on the control voltage $`V`$ for different temperatures. It is seen that the temperature dependence of the maximum of $`I_m`$ which is achieved at $`eV\mathrm{\Delta }`$ is much weaker than the $`I_c(T)`$ dependence. These results qualitatively agree with the experimental data .
4. Spin injection and the critical Josephson current.
In this section we present results of the study of spin injection on the critical Josephson current $`I_c`$ for the structure shown in Fig.1b . We calculate the distribution function $`f_+`$ which determines $`I_c`$ for parallel $`()`$ and antiparallel $`()`$ orientations of the magnetization in the ferromagnetic (F) reservoirs. It is assumed that the voltage $`2V`$ is applied between the F reservoirs and a current flows between the N reservoirs. We consider different limits of the length of the N wire.
4.1. Mesoscopic limit: $`L<\{L_ϵ,L_{sp}\}.`$
In this limit the function $`f_+(f_++f_+)/2`$ does not depend on mutual orientation of ferromagnets and has the same form as in the nonmagnetic case: $`f_+=F_{V+}.`$ The conductance $`G`$ and magnetic moment of injected spins $`M`$ (per unit volume of the N wire) depend on mutual orientation of the ferromagnetic domains. In the case of the parallel orientation $`()`$ we obtain $`G=(2R_b+R_L)^1+(2R_b+R_L)^1`$ and $`M=0`$. In the case of the antiparallel orientation $`()`$ we have $`G=2/(R_b+R_b+R_L)`$ and $`M=4\mu (R_bR_b)V\beta /\sigma ^2(2R_b+R_L)(2R_b+R_L)(R_b+R_2+R_L).`$ Here $`R_{b,}`$ and $`R_L`$ are the F/N interface and the N wire resistances for electrons with up and down spins . Since the interface conductances $`R_{b,}^1`$ differ from each other (they are proportional to different density of states in the ferromagnets $`\nu _{F,}`$), spin injection takes place in the antiparallel $`()`$ configuration. However the critical current $`I_c`$ is the same for both configurations. Its dependence on $`V`$ is similar to that found for the structure under consideration with normal (nonmagnetic) reservoirs .
4.1. Semi-mesoscopic limit: $`L_ϵ<L<L_{sp}.`$
In this limit there is no spin relaxation in the N wire and the magnetization $`M`$ is determined again by formulae presented above. On the other hand the magnetization is related to a potential $`V_{sp}`$ which determines an imbalance between the spin subsystems in the N wire: $`M=2\mu \nu eV_{sp}`$. The function $`f_+(f_++f_+)/2`$ has an equilibrium form corresponding to the potential $`V_{sp}:f_+=[\mathrm{tanh}(ϵ+eV_{sp})\beta +\mathrm{tanh}(ϵeV_{sp})\beta ]/2`$. Therefore in the case of parallel $`()`$ orientation $`I_c`$ does not depend on $`V`$ as $`V_{sp}=0`$ ($`V`$ is proportional to $`V_{sp}`$). In the case of antiparallel $`()`$ orientation $`I_c`$ does depend on $`V_{sp}`$ in the same way as in the section 4.1. A similar effect of spin injection into a supereconductor on the energy gap $`\mathrm{\Delta }`$ was analysed in Ref. . The critical Josephson current is more sensitive to the form of the distribution function than the energy gap $`\mathrm{\Delta }`$. Therefore the effect considered here allows one to realise the $`\pi `$junction and to make certain conclusions about the energy and spin relaxation rate.
5. Giant thermoelectric effect
It is well known that if the terminals of a normal conductor are maintained at different temperatures $`\delta T=T_o\pm \delta T`$, then in the absence of the current an thermoemf $`V_{emf}`$ appears at the terminals. The magnitude of $`V_{emf}`$ is equal to: $`eV_{emf}=c_1(T/ϵ_F)\delta T,`$ where $`c_1`$ is a factor of the order 1 and $`ϵ_F`$ is the Fermi energy. In this report we study the thermoelectric effect in the mesoscopic structure shown in Fig.1a. It will be shown that a temperature difference between the normal (N) reservoirs leads to a voltage between normal and superconducting curcuits $`V_T`$. The magnitude of this voltage does not contain the small parameter $`(T/ϵ_F),`$ and besides it oscillates with a variation of the phase difference $`\phi `$ between the superconductors.
We assume that the superconductors are connected via a superconducting loop and the phase difference between them $`\phi `$ is controlled by an applied magneitic field. The N reservoirs are disconnected and maintained at different temperatures $`T(\pm L)=T_o\pm \delta T`$. We will calculate the electric potential in the N film and in particular the potential $`V_T`$ in the N reservoirs . Since we set the potential in the superconductors equal to zero, the potential $`V_T`$ is the voltage difference between the N reservoirs and superconductors which arises in the presence of the temperature difference $`\delta T.`$ In the limit of high interface resistance (see condition (5)) we obtain from Eq.(2) in the main approximation in $`r:`$ $`f_+(x)J_+x+f_+(0),`$ where $`f_+(0)=f_{eq}\mathrm{tanh}(\beta _oϵ),`$ and $`J_+=\delta \beta ϵ/L\mathrm{cosh}^2(\beta _oϵ),`$ $`\delta \beta =\beta \delta T/T_o.`$ We assumed the ratio $`\delta T/T_o`$ to be small and expant the distribution functions in the N reservoirs.
In order to find the function $`f_{},`$ we consider Eq.(2) and take into account that $`J_2`$ is zero (no current through the N reservoirs). In the absence of the temperature gradient we obtain from this equation
$$f_{}=0,\text{ }J_Sf_{eq}=J_{1eq}=r(g_z+g_{z+})f_{eq}$$
(7)
The temperature gradient leads to the non-zero function $`f_{}(x)`$ which determines the electrical potential $`V_T(x)`$. We find from Eq.(2)
$$f_{}(x)=f_{}(0)+\delta J_1x+_0^x𝑑x_1[J_{an}_x\delta f_+(x_1)J_s\delta f_+(x_1)]$$
(8)
where $`\delta f_+(x)f_+(x)f_{eq}`$. From Eq.(1) we obtain for $`\delta J_1J_1J_{1eq}`$
$$\delta J_1=\pm r[g_{z+}\delta f_+(\pm L_1)g_{1+}f_{}(\pm L_1)]$$
(9)
Here the signs $`\pm `$ relate to the points $`\pm L_1`$. From Eqs.(8) and (9) it follows that $`\delta J_1=0,`$ and in the main approximation in $`r`$ the distribution function is $`f_{}(0)\delta f_+(\pm L_1)g_{z+}/g_{1+},`$ where $`\delta f_+(\pm L_1)=\pm J_+L_1`$. The electric potential $`V_T`$ is related to the function $`f_{}`$ and is an even function of $`x`$
$$eV_T=_0^{\mathrm{}}𝑑ϵf_{}=\frac{L_1}{L}\frac{\delta T}{T_o}\beta _0^{\mathrm{}}𝑑ϵ\frac{g_{z+}}{g_{1+}}\frac{ϵ}{\mathrm{cosh}^2(\beta _oϵ)}$$
(10)
Here we substitute the function $`f_{}(0)`$ which is the main contribution to the function $`f_{}(x).`$ Therefore the potential $`V_T`$ arising in the presence of the temperature gradient is almost constant along the N wire. This potential equals approximately the voltage difference $`V_T`$ between the N reservoirs and superconducting loop. It is worth noting that $`V_T`$ determined by Eq.(10) does not depend on the small parameter $`r`$ (however one should have in mind that the condition (5) imposes limits on this parameter). We calculated the integrand using both analytical an numerical solutions of the Usadel equation. In Fig.3 we show the temperature dependence of $`V_T`$ for $`\phi =\pi /2`$. We see that the dependence $`V_T(T)`$ is nonmonotonic (reentrant behaviour). One can easily estimate $`V_T`$ on the order of magnitude. We obtain
$$eV_T=\delta T(L_1/L)\mathrm{sin}\phi \{\begin{array}{c}(T/ϵ_{Th})C_1(\phi ),\text{ }T<<ϵ_{Th}\\ (ϵ_{Th}/T)^2C_2(\phi ),\text{ }T>>ϵ_{Th}\end{array}$$
(11)
Here $`C_{1,2}(\phi )`$ are periodic functions of the phase difference $`\phi `$; they are of order 1 and are not zero at $`\phi =0.`$ If the ratio $`T/ϵ_{Th}`$ is of order 1, then the thermoemf is of the order $`\delta T`$, that is, several orders greater than the thermoemf in the normal metals.
The physical explanation of the effect is the following. The temperature gradient creates a deviation of the distrubtion function $`\delta f_+=\delta (n+p)`$ from equilibrium. On the other hand the superconductors do not affect this function because complete Andreev reflections conserve the total number of excess electrons and holes. The function $`\delta f_+`$ has different signs at $`\pm L_1`$ and leads to an additional Josephson current $`\delta J_S=rg_{z+}\delta f_+(L_1)`$ of the same sign at $`\pm L_1`$(sign of the function $`g_{z+}`$ is different at these points). In order to compensate these additional currents, the potential $`V_T`$ arises in the N wire producing a subgap current $`\delta J_{sg}=rg_{1+}f_{}(\pm L_1)`$ which compensates the current $`\delta J_S`$.
We have neglected the ordinary thermoelectric current $`I_T`$ in the N wire because it contains a small parameter $`(T/ϵ_T)`$ and leads to a small contribution to $`V_T.`$ Recently the influence of the proximity effect on the ordinary thermoelectric effect in the mesoscopic S/N structure was studied both theoretically and experimentally . The effect analysed here is much stronger.
In simplified models we have analysed three possible effects in the 4-terminals S/N mesoscopic structures; one of them (see section 3) has been observed experimentally. Further theoretical studies of these and other effects are needed to account for processes which were neglected in our analysis. It is interesting, for example, to investigate how the obtained results are changed by the energy relaxation in the N wire which seems to be faster than it was expected
We are grateful to the EPSRC for their financial support.
FIG. Phase-coherent effects in multiterminal superconductor/normal metal mesoscopic structures Schematic view of the 4-terminal S/N/S structure under consideration.
FIG. Phase-coherent effects in multiterminal superconductor/normal metal mesoscopic structures The measured critical current ($`I_m`$) vs $`V`$ for different temperatures: $`\beta =ϵ_{Th}/2T`$. The parameters are: $`\mathrm{\Delta }=10ϵ_{Th},L_1/L=0.3,r=0.3`$.
FIG. Phase-coherent effects in multiterminal superconductor/normal metal mesoscopic structures The dependence of the normalised voltage $`\stackrel{~}{V_T}`$ on inverse temperature $`\beta `$ at $`\phi =\pi /2`$ (the parameters are $`\mathrm{\Delta }/ϵ_{Th}=10,L_1/L=0.5`$).
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# On dispersive energy transport and relaxation in the hopping regime
## I Introduction
In recent years much attention has been devoted to the study of relaxation processes of non-equillibrium charge carriers in strongly localized systems, where transport proceeds via phonon-assisted hopping, like photoexited charge carriers in band tails (see e.g. -) and Anderson insulators (see e.g.-). In such systems particular small relaxation speeds are observed. Often the smallness of the relaxation speed is attributed to interaction effects. However, also in strongly localized non-interacting electron systems long lasting relaxation processes are known to be the rule, and not the exception .
The theoretical investigation of such relaxation processes is notoriously difficult, since the system is strongly disordered and always in a transient state. In most problems one is interested in time scales which are large as compared to the time a single hop needs. For such time scales all quantities depend strongly on frequency, also for very low frequencies, so that the consideration of dispersive transport is vital. On the other hand in most problems of interest both spatial and energetical disorder exists. Due to the latter fact the transitions are inelastic. The inelastic character of the transitions, the relaxation, leads to a flow of energy from the electron system to the phonon system. Due to disorder this transport of energy is also dispersive. Therefore, dispersive energy transport and relaxation are intimately connected. The investigation of the relaxation process requires both the consideration of dispersive particle transport and the consideration of dispersive energy transport.
The intricate physical situation manifests also in the equations, which have to be considered. Since the transport is inelastic already in the simplest approximation the investigation requires to find solutions to integral equations (see e.g.\- ). Due to the fact that the particle moves in an energy dependent density of states the kernel of these integral equations does not depend on the difference between the site energies only. The traditional method to handle the situation, the percolation theory, is not applicable here. The effective-medium methods by Movaghar and coworkers give, as pointed out by Movaghar and coworkers, wrong results for systems at low temperature . Furthermore, beyond the Markovian approximation the derivation of these integral equations represents itself an intricate problem. To our knowledge, beside the attempts by Movaghar and coworkers, numerical investigations (see e.g., ), and physical intuitive considerations (see e.g. ,), mainly Markovian equations have been used (see e.g. ).
It is the aim of the present paper to fill this gap and to provide a formalism that can be used for studying relaxation phenoma of strongly localized, charge carriers far from equillibrium in the hopping regime taking into account both dispersive particle transport and dispersive energy transport. To this end we focus, for the sake of definitness, on relaxation of photoxited, non-interacting, non-equillibrium charge carrires in band tails of, e.g. amorphous semiconductors, like amorphous Si:H. On the first glance, this problem seems to be rather special. If the number of charge carriers exited is small Fermi-correlation is negligible, so that one has to cope with linear rate equations. The linearity of the transport equations is, of course, the basic ingredient in solving the problem. A closer look on the problems of interest reveals, however, that most of the problems can be formulated in this way. Clearly, how to achieve linearization depends on the experiment chosen, but, provided the number of charge carrires exited is small, always the smallness of the particle number can be invoked. In this case, the number of particles at a site is certainly not small as compared to its equillibrium value, but small as compared to unity. After linearization the structure is, in principle, quite similar as for the case considered here.
Below we present the derivation of a framework for the consideration of relaxation phenomena due to phonon-assisted hopping at zero temperature. We focus on the limit of strong localization, where dispersive effects are expected to be most pronounced. Here the disorder manifests in a strong dependence of the transport coefficients on frequency for very low frequencies. Consequently, in this regime the diffusion propagator can not be calculated from Markovian transport equations, like those used e.g. in Ref.. To simplify our integral equations we use the concept of quasi-elasticity, the particle changes its energy only by a small amount by one hop.
The method is applied to the study of dispersive energy transport. We find that always two tendencies are present, if only the density of states not decreases with increasing energy. First, there is a widening of the packet due to the statistical spreading, the dispersive energy transport. Second, there is narrowing of the packet due to decrease of the density of states with decreasing energy. The overall evolution is determined by the interplay of these two tendencies. We have studied both tendencies in two limiting cases, for a particle moving in a constant density of states and for a particle moving in an exponential density of states. In the first case the impact of the density of states decrease is absent, so that only the statistical spreading is present. Here the packet evolves into a Gaussian packet moving with constant drift velocity in energy space. The packet width increases with time as $`\sqrt{t}`$. The other result is obtained for exponential density of states. Here both tendencies are present. We find, that the velocity of the packet is strongly slowing down with time. In that case the mean square deviation of energy eventually becomes time-independent. Consequently, the motion of the packet in energy space is “soliton-like”. The concrete results on the diffusion propagator, its time dependence and its moments for exponential densities of states are of relevance for photoluminiscense experiments on amorphous Si:H.
## II Basic equations
We consider photoexited, localized charge carriers in band tails at zero temperture. After exitation the charge carriers lower their energy by phonon-assisted hops between localized states. Since $`T=0`$, only hops from higher to lower energy occur. In this situation the charge carriers are in strong non-equilibrium. We assume, that the number of excited charge carriers is small, and their energies are far from the Fermi level, so that it is very unlikely that an electron jumps to a site already occupied. Consequently, we can neglect Fermi-correlation. In this case the electron transport can be described by the rate equation .
$$\frac{dp_m}{dt}=\underset{m^{}}{}\left(p_m^{}W_{m^{}m}p_mW_{mm^{}}\right).$$
(1)
In calculating the transition probabilities we assume, that the electron-phonon coupling strength is weak, so that only one-phonon processes have to be taken into account. Then, at zero temperature, the transition probabilities are given by
$$W_{m^{}m}=\mathrm{\Theta }\left(\omega \epsilon _m^{}+\epsilon _m\right)\mathrm{\Theta }\left(\epsilon _m^{}\epsilon _m\right)W\left(\left|𝐑_{mm^{}}\right|\right),$$
(2)
where
$$W\left(\left|𝐑_{mm^{}}\right|\right)=\nu e^{2\alpha \left|𝐑_{mm^{}}\right|}.$$
(3)
Here $`\alpha `$ is the inverse of localized state radius, and $`\nu `$ is the phonon frequency. The energy $`\omega `$ is the upper bound for the energy transfer by one hop. Note, that in the materials of interest $`\omega `$ can be much smaller as the Debye-frequency, since not all phonons can interact with localized electrons equally well. Short wave-length phonons are ineffective since the electron-phonon coupling constant tends to zero for momenta $`q`$ with $`q/2\alpha 1`$. Therefore, the effective upper phonon momentum is of the order $`2\alpha `$, and not of the order of the inverse lattice constant of the host material. Furthermore, in disordered systems the high-energetic phonons are localized, and need not contribute to transport by necessity.
The first step function in front of the transition probabilities restricts the transitions to transitions between sites separated by at most $`\omega `$ in energy space. Thus, it decreases the energy relaxation speed. In impurity conduction, and in nearly all papers on relaxation of charge carriers in band tails, this step function is usually replaced by unity. In impurity conduction this is quite reasonable, since it is assumed that hops are restricted within narrow stripe near the Fermi level, which is small compared to Debye energy. In the band-tail problem, however, we can see no reason for neglecting it in advance.
To calculate the transport quantities of interest, we have to calculate the diffusion propagator. In order to render the analytical calculations feasible we introduce continuous coordinates. The change of representation is defined by:
$$n(\rho )=\underset{m}{}p_m\delta \left(\rho \rho _m\right),$$
(4)
where $`\rho =(𝐑,\epsilon )`$ and $`\rho _m=(𝐑_m,\epsilon _m)`$. In this representation the rate equation (1) takes the form:
$$\frac{dn\left(\rho \right)}{dt}=𝑑\rho ^{}n\left(\rho ^{}\right)V(\rho ^{},\rho ),$$
(5)
where $`V`$ is determined by the equations:
$$V(\rho ^{},\rho )=𝑑\rho _1\eta \left(\rho _1\right)w_{\rho _1}(\rho ^{},\rho ),$$
(6)
$$w_{\rho _1}(\rho ^{},\rho )=W(\rho ^{},\rho _1)\left[\delta \left(\rho _1\rho \right)\delta \left(\rho ^{}\rho \right)\right],$$
(7)
$$\eta \left(\rho \right)=\underset{m}{}\delta \left(\rho \rho _m\right).$$
(8)
The Laplace-transformed equation is given by:
$$sn\left(\rho \right)n_o\left(\rho \right)=𝑑\rho ^{}n\left(\rho ^{}\right)V(\rho ^{},\rho ),$$
(9)
where $`n_o(\rho )=n(\rho ,t=0)`$ is the initial condition. We will assume that $`p_m(t=0)`$ is a function $`p_0(𝐑_m,\epsilon _m)`$, so that $`n_o(\rho )=p_0(\rho )\eta (\rho )`$.
Equation (9) can be solved using the Green’s function method. The solution is given by:
$$n\left(\rho \right)=𝑑\rho ^{}n_o\left(\rho ^{}\right)\mathrm{\Phi }(\rho ^{},\rho ).$$
(10)
The Green’s function satisfies the equation:
$$s\mathrm{\Phi }(\rho ^{},\rho )𝑑\rho _1V(\rho ^{},\rho _1)\mathrm{\Phi }(\rho _1,\rho )=\delta \left(\rho ^{}\rho \right).$$
(11)
Note, that, due to probability conservation, the Green’s function $`\mathrm{\Phi }`$ and the probability $`w_{\stackrel{~}{\rho }}`$ satisfy the relations:
$$𝑑\rho \mathrm{\Phi }(\rho ^{},\rho )=\frac{1}{s},$$
(12)
$$𝑑\rho w_{\stackrel{~}{\rho }}(\rho ^{},\rho )=0.$$
(13)
## III Configuration Average
In order to calculate the configuration average, we assume, that the sites are distributed homogeneously in space. The distribution of site energies is supposed to be given by some distribution function $`p(\{\epsilon _i\})`$. Accordingly, the average of any quantity, depending on the energy and positions of sites, is given by:
$$A=\mathrm{\Pi }_m\frac{d𝐑_m}{𝒱}𝑑\epsilon _mp\left(\left\{\epsilon _m\right\}\right)A\left(\{𝐑_m,\epsilon _m\}\right),$$
(14)
where $`𝒱`$ is the volume of the system. The application of the averaging procedure to the structural factor $`\eta `$ serves, in particular, as a definition for the density of states, i.e.
$$𝒩\left(\epsilon \right)=\eta \left(\rho \right).$$
(15)
Products of the structural factor $`\eta `$ are averaged according to:
$$\eta \left(\rho _1\right)\mathrm{}\eta \left(\rho _n\right)=𝒩\left(\rho _1\right)\delta \left(\rho _1\rho _2\right)\mathrm{}\delta \left(\rho _{n1}\rho _n\right).$$
(16)
Using these definitions, the configuration average of Eq.(10) can be calculated diagrammatically . The diagrammatic method leads to the following set of equations for the calculation of the configuration average $`n(\rho )`$ of the electron density $`n(\rho )`$ , (see Figs 14 for illustration):
$$n\left(\rho \right)=𝑑\rho _1𝑑\rho _2p_0\left(\rho _1\right)S(\rho _1,\rho _2)F(\rho _2,\rho ),$$
(17)
$$sF(\rho ^{},\rho )=\delta \left(\rho ^{}\rho \right)+𝑑\rho _1\mathrm{\Pi }(\rho ^{},\rho _1)F(\rho _1,\rho ),$$
(18)
$$\mathrm{\Pi }(\rho ^{},\rho )=𝑑\rho _1𝒩\left(\rho _1\right)\mathrm{\Pi }_{\rho _1}(\rho ^{},\rho ),$$
(19)
$`\mathrm{\Pi }_{\rho _1}(\rho ^{},\rho )=w_{\rho _1}(\rho ^{},\rho )`$ (20)
$`+`$ $`{\displaystyle 𝑑\rho _2𝑑\rho _3w_{\rho _1}(\rho ^{},\rho _2)F(\rho _2,\rho _3)\mathrm{\Pi }_{\rho _1}(\rho _3,\rho )},`$ (21)
$$S(\rho ^{},\rho )=𝒩\left(\epsilon ^{}\right)\left[\delta \left(\rho ^{}\rho \right)+𝑑\rho _1F(\rho ^{},\rho _1)\mathrm{\Pi }_\rho ^{}(\rho _1,\rho )\right].$$
(22)
Here $`F(\rho ^{},\rho )=\mathrm{\Phi }(\rho ^{},\rho )`$. A detailed derivation of the set of integral equations is given in Appendix A.
## IV Effective Medium
Given the system of integral equations (18)-(22), the main problem is to find an approximate self-consistent solution to it. The situation is quite similar to that of the calculation of the equilibrium conductivity in a disordered system . There an approximate solution of this system could be found by introducing a proper decomposition of the function $`F`$, the diffusion propagator, into short- and long-wavelength limit, according to:
$$F(\rho ^{},\rho )=f(s)C(\epsilon )\delta \left(\rho ^{}\rho \right)+\stackrel{~}{F}(\rho ^{},\rho ),$$
(23)
where $`f(s)`$ was a frequency dependent parameter that could be related to the critical hopping length $`R_c`$ via the equation:
$$f\nu =\mathrm{exp}\left(2\alpha R_c\right),$$
(24)
and $`C(\epsilon )`$ was an energy dependent function determined by the principle of detailed balance. This decompostion originated from the notion, that the integrals in the integral equations are governed by products of transition probabilities and diffusion propagators, and, for strongly localized electrons, the latter quantities are short ranged functions as compared to the transition probabilities. Using the decomposition (23), the effective medium approximation reduces to the repacement of $`F`$ by $`fC(\epsilon )\delta (\rho ^{}\rho )`$ in the calculation of the effective transition probability $`\mathrm{\Pi }`$ (Eq.(21)) and of the irreducible block $`S`$ (Eq.(22)).
Here in the band-tail problem we use the same philosophy. We first decompose the diffusion propagator in two parts, according to:
$$F(\rho ^{},\rho )=f(\epsilon ,s)\delta \left(\rho ^{}\rho \right)+\stackrel{~}{F}(\rho ^{},\rho ).$$
(25)
Then, to investigate the consequences of this decomposition, we insert Eq.(25) into Eq.(21). Performing a partial summation we obtain:
$`\mathrm{\Pi }_{\stackrel{~}{\rho }}(\rho ^{},\rho )=\stackrel{~}{w}_{\stackrel{~}{\rho }}(\rho ^{},\rho )`$ (26)
$`+`$ $`{\displaystyle 𝑑\rho _1𝑑\rho _2\stackrel{~}{w}_{\stackrel{~}{\rho }}(\rho ^{},\rho _1)\stackrel{~}{F}(\rho _1,\rho _2)\mathrm{\Pi }_{\stackrel{~}{\rho }}(\rho _2,\rho )},`$ (27)
where the renormalized transition probabilities $`\stackrel{~}{w}_{\stackrel{~}{\rho }}(\rho ^{},\rho )`$ are given by Eq.(7), with $`W`$ replaced by:
$$\stackrel{~}{W}(\rho ^{},\rho ;s)=\frac{W(\rho ^{},\rho )}{1+f(\epsilon ^{},s)W(\rho ^{},\rho )}.$$
(28)
Note that the renormalized transition probability depends now on $`s`$ via the function $`f(\epsilon ,s)`$.
At this point the introduction of the function $`f(\epsilon ,s)`$, the effective medium, still seems to be rather arbitrary. However, if we now choose the effective medium in such a way, that the integrals over $`\stackrel{~}{F}`$ vanish, the renormalized transition probabilities yield the exact solution of the diffusion problem. In that case the function $`S`$ turns into:
$$S(\rho ^{},\rho )=𝒩\left(\epsilon ^{}\right)\delta \left(\rho ^{}\rho \right),$$
(29)
so that Eq.(17) can be cast into the form:
$$n(\rho ;s)=𝑑\rho ^{}n_o\left(\rho ^{}\right)F(\rho ^{},\rho ).$$
(30)
Therefore the function $`F(\rho ^{},\rho )`$ can be identified with the diffusion propagator.
Note, that although the general philosophy is quite the same as in equilibrium, the situation is much more intricate here, — the principle of detailed balance is absent, since we are dealing with the situation at zero temperature, and so it can’t determine the energy dependence of the effective medium. Moreover, we expect the critical hopping length to depend somehow on the position of the electron in the tail. Thus, in contrast to equilibrium, $`f`$ is not only a frequency dependent parameter but also an energy dependent function, which has to be determined self-consistently.
## V Diffusion propagator in effective medium approximation
For the moment we put aside the question of the determination of the effective medium to elaborate further on the consequences of the renormalization of the transition probabilities. To this end we focus on the diffusion propagator.
The equation for the diffusion propagator is given by Eq.(14). In effective medium approximation, when calculating the irreducible part $`\mathrm{\Pi }`$, only the lowest order contribution to $`\mathrm{\Pi }`$ with respect to $`\stackrel{~}{F}`$ is taken into account, so that $`\mathrm{\Pi }_{\stackrel{~}{\rho }}=\stackrel{~}{w}_{\stackrel{~}{\rho }}`$. In that approximation the equation for the diffusion propagator in momentum representation reads:
$$sF\left(𝐪|\epsilon ^{},\epsilon \right)=\delta \left(\epsilon ^{}\epsilon \right)+𝑑\epsilon _1\left\{F\left(𝐪|\epsilon ^{},\epsilon _1\right)\stackrel{~}{W}\left(𝐪|\epsilon _1,\epsilon ;s\right)𝒩\left(\epsilon \right)F\left(𝐪|\epsilon ^{},\epsilon \right)\stackrel{~}{W}\left(\mathrm{𝟎}|\epsilon ,\epsilon _1;s\right)𝒩(\epsilon _1)\right\}.$$
(31)
Of course, as we don’t know how the effective medium looks like so far, and moreover, as the equation is a complicated integral equation, we can’t find a solution. The fact, that the effective medium is a function of energy, makes the problem much more complicated. Progress can only by achieved if we can find arguments to simplify the equation considerably. To this end we focus on the renormalized transition probability. According to the Eqs. (2) and (28), it is given by:
$$\stackrel{~}{W}\left(R|\epsilon ^{},\epsilon ;s\right)=\mathrm{\Theta }\left(\epsilon ^{}\epsilon \right)\mathrm{\Theta }\left(\omega \epsilon ^{}+\epsilon \right)\frac{W\left(R\right)}{1+f(\epsilon ^{},s)W(R)}=\mathrm{\Theta }\left(\epsilon ^{}\epsilon \right)\mathrm{\Theta }\left(\omega \epsilon ^{}+\epsilon \right)\stackrel{~}{W}\left(R|\epsilon ^{};s\right).$$
(32)
Owing to the step functions in front of the transition probability, the energy integrations are restricted to intervals of length $`\omega `$. Taking this fact into account the integral equation for the calculation of the diffusion propagator takes the form:
$$sF(𝐪|\epsilon ,\epsilon )=\delta (\epsilon ^{}\epsilon )+\underset{0}{\overset{\omega }{}}d\epsilon _1\{F(𝐪|\epsilon ^{},\epsilon +\epsilon _1)\stackrel{~}{W}(𝐪|\epsilon _1+\epsilon ;s)𝒩\left(\epsilon \right)F(𝐪|\epsilon ^{},\epsilon )\stackrel{~}{W}(0|\epsilon ;s)𝒩(\epsilon _1+\epsilon )\}.$$
(33)
Further we assume that the diffusion propagator, the transition probability $`\stackrel{~}{W}(R,\epsilon ;s)`$, and the density of states $`𝒩\left(\epsilon \right)`$ are slowly varying functions on intervals of length $`\omega `$, so that the integrand can be expanded with respect to $`\epsilon _1`$. Doing so we obtain:
$$sF(𝐪|\epsilon ,\epsilon )=\delta (\epsilon ^{}\epsilon )+\omega 𝒩\left(\epsilon \right)[\stackrel{~}{W}(𝐪|\epsilon ;s)\stackrel{~}{W}(\mathrm{𝟎}|\epsilon ;s)]F(𝐪|\epsilon ^{},\epsilon )+\frac{1}{2}\omega ^2\frac{}{\epsilon }\left[F(𝐪|\epsilon ^{},\epsilon )\stackrel{~}{W}(𝐪|\epsilon ;s)𝒩\left(\epsilon \right)\right].$$
(34)
We terminate the expansion after the first derivative with respect to energy. This term describes the biased motion of the particle from sites of higher to sites of lower energy. For finite temperatures one would have to replace $`F\left(𝐪|\epsilon ^{},\epsilon \right)`$ in the last term with $`F\left(𝐪|\epsilon ^{},\epsilon \right)+kTF\left(𝐪|\epsilon ^{},\epsilon \right)/\epsilon `$ which would describe the energy diffusion current. At zero temperature there is no input of energy from the phonon-system into the particle. Therefore, we expect the additional term to be rather small, and can be neglected.
To investigate the transport properties, we restrict our attention with the long-wavelength limit. In that case the elastic contribution yields the diffusion coefficient for particle diffusion. The term, containing the derivative with respect to energy is finite for $`q0`$. If we are interested only in the long-wavelength limit, we can put here $`q=0`$, since the remaining terms are of higher order with respect to ’$`q`$’ and ’$`\omega /\epsilon `$’. Since a nonzero momentum in this term couples particle diffusion to energy transport, this approximation corresponds to a decoupling of these two processes. Then, in the long-wavelength limit, we obtain:
$`sF\left(𝐪|\epsilon ^{},\epsilon \right)`$ $`=`$ $`\delta \left(\epsilon ^{}\epsilon \right)D(\epsilon ,s)q^2F\left(𝐪|\epsilon ^{},\epsilon \right)`$ (35)
$`+`$ $`{\displaystyle \frac{}{\epsilon }}\left[F\left(𝐪|\epsilon ^{},\epsilon \right)v(\epsilon ,s)\right],`$ (36)
where:
$$D(\epsilon ,s)=\frac{1}{2}\omega _q^2\stackrel{~}{W}\left(q|\epsilon ;s\right)|_{𝐪=\mathrm{𝟎}},$$
(37)
and:
$$v(\epsilon ,s)=\frac{1}{2}\omega ^2𝒩\left(\epsilon \right)\stackrel{~}{W}\left(0|\epsilon ;s\right).$$
(38)
Eq.(36) is easily solved. It’s solution is:
$$F\left(𝐪|\epsilon ^{},\epsilon \right)=\frac{\mathrm{\Theta }\left(\epsilon ^{}\epsilon \right)}{v(\epsilon ;s)}\mathrm{exp}\left[\underset{\epsilon }{\overset{\epsilon ^{}}{}}𝑑\epsilon _1\frac{s+D(\epsilon _1;s)q^2}{v(\epsilon _1;s)}\right].$$
(39)
Explicit expressions for the transport coefficients can be found in Appendix B.
Now, at this stage, the validity of our quasi-elastic approximation requires, that the second derivative terms are small as compared to terms with first derivatives with respect to energy. This requirement imposes the following restrictions on the transport coefficients $`v`$ and $`D`$, the frequency $`s`$ and the momentum $`q`$:
$$\omega \left|\frac{1}{v}\frac{dv}{d\epsilon }\right|1,$$
(40)
$$\omega \left|\frac{s}{v}\right|1,$$
(41)
$$q^2\omega \frac{D}{v}1.$$
(42)
The applicability of the quasi-elastic approximation was also discussed in Ref.. There it was concluded that this approximation should be inapplicable. To substantiate this statement numerical calculations were invoked. However, in interpreting these data it has to be taken into account, that they have been obtained using a model, neglecting a weighting of the transition probabilities according to the number of phonons emitted, so that hops between nearly isoenergetic sites were treated as likely as hops from the very top to the very bottom of the tail. Consequently the discussion in Ref. applies only to system with sufficient strong electron-phonon interaction, but not to systems with weak electron-phonon interaction. For extraordinary deep hops to contribute to the diffusion propagator they should be characteristic for an ensemble of electrons. This is, however, not expected and in the experimental data, e.g. on amorphous Si:H , not observed. For these reasons the discussion in Ref. does not apply to our model.
## VI Physical quantities and interpretation
Before establishing a self-consistency equation, we first elaborate further on the physical content of our diffusion equation. To this end let us first have a closer look on the coefficients $`D(\epsilon ,s)`$ and $`v(\epsilon ,s)`$. Imagine, that we have a particle initially located at ($`𝐑_0,\epsilon _0`$). Then the initial condition is $`n_o(𝐑,\epsilon )=\delta (𝐑𝐑_0)\delta (\epsilon \epsilon _0)`$. According to our approximation, the motion of the charge carrier is composed of two contributions: particle diffusion between isoenergetic sites, and relaxation in energy space. Characteristics of these two processes are the mean square displacement and the energy relaxation speed. It turns out, that both can be calculated from the function:
$`P_L(\epsilon _0,\epsilon ;s)={\displaystyle \frac{\mathrm{\Theta }\left(\epsilon _0\epsilon \right)}{s}}{\displaystyle \frac{}{\epsilon }}\mathrm{exp}\left[s{\displaystyle \underset{\epsilon }{\overset{\epsilon _0}{}}}{\displaystyle \frac{d\epsilon _1}{v(\epsilon _1,s)}}\right]`$ (43)
$`={\displaystyle \frac{\mathrm{\Theta }\left(\epsilon _0\epsilon \right)}{v(\epsilon ,s)}}\mathrm{exp}\left[s{\displaystyle \underset{\epsilon }{\overset{\epsilon _0}{}}}{\displaystyle \frac{d\epsilon _1}{v(\epsilon _1,s)}}\right],`$ (44)
which is just $`F(q=0|\epsilon _0,\epsilon )`$.
We define the mean energy of the particle by:
$$E\left(s\right)=𝑑\rho ^{}\epsilon ^{}n(\rho ^{},s).$$
(45)
Using our diffusion propagator, this equation can be rewritten in the form:
$$E\left(s\right)=𝑑\epsilon ^{}P_L(\epsilon _0,\epsilon ^{};s)\epsilon ^{},$$
(46)
and from the diffusion equation we obtain after integration by parts:
$$sE\left(s\right)\epsilon _0=𝑑\epsilon P_L(\epsilon _0,\epsilon ;s)v(\epsilon ,s).$$
(47)
Therefore, in general, the time dependence of the mean energy is given by complicated integrals. These integrals simplify considerably in two limiting cases: in the absence of dispersive energy transport, and for energy independent $`v(\epsilon ,s)`$. In the absence of dispersive energy transport, i.e. in the Markovian limit in which the transport coefficients are independent of $`s`$, the integrals can readly be calculated in time representation. In the latter situation we simply obtain:
$$P_L(\epsilon _0,\epsilon ;t)=\mathrm{\Theta }\left(\epsilon _0\epsilon \right)\delta \left(\epsilon _m\left(t\right)\epsilon \right),$$
(48)
where $`\epsilon _m(t)`$ is defined by:
$$t=\underset{\epsilon _m}{\overset{\epsilon _0}{}}\frac{d\epsilon _1}{v\left(\epsilon _1\right)}.$$
(49)
Therefore we obtain:
$$\frac{dE\left(t\right)}{dt}=\frac{d\epsilon _m\left(t\right)}{dt}=v\left(\epsilon _m\left(t\right)\right).$$
(50)
Consequently, $`v`$ is the velocity of energy relaxation, and $`\epsilon _m(t)`$ is the instantenous position of the particle in energy space. In general, however, in disordered systems the coefficient $`v(\epsilon ,s)`$ depends on $`s`$, so that energy transport is dispersive. In that case the integral in Eq.(47) can only be calculated easily if $`v`$ is independent of energy. Then, owing to probability conservation,
$$\left(\frac{dE}{dt}\right)\left(s\right)=\frac{v\left(s\right)}{s}$$
(51)
is obtained.
The mean square displacement, defined by:
$$R^2\left(t\right)=d\rho ^{}R^{}{}_{}{}^{2}n\left(\rho ^{}\right),$$
(52)
can be obtained from similar arguments. In Fourier representation this equation can be written in the form:
$$\left(\frac{dR^2}{dt}\right)\left(s\right)=s\mathrm{\Delta }_q|_{q=0}𝑑\epsilon ^{}F\left(q|\epsilon _0,\epsilon ^{}\right).$$
(53)
The derivative of the diffusion propagator can be calculated using the diffusion equation. Doing so, we obtain:
$$\left(\frac{dR^2}{dt}\right)\left(s\right)=2d𝑑\epsilon P_L(\epsilon _0,\epsilon ;s)D(\epsilon ,s).$$
(54)
Again Eq.(54) simplifies only in two special cases, in the Markovian limit and for energy independent diffusion coefficients. Whereas in the first case
$$\frac{dR^2}{dt}\left(t\right)=2dD\left(\epsilon _m\left(t\right)\right)$$
(55)
is obtained, we have
$$\left(\frac{dR^2}{dt}\right)\left(s\right)=2d\frac{D\left(s\right)}{s}$$
(56)
in the latter situation. Below we shall see, that energy independent transport coefficients are obtained for constant density of states only.
## VII The self-consistency equation
So far we have only investigated the consequences of the renormalization. In order to complete the approximation scheme, we still have to calculate the effective medium itself. Of course, we can not calculate the effective medium exactly, since this would amount to find an exact solution to the diffusion problem. Rather we shall try to calculate $`f(\epsilon ,s)`$ self-consistently.
The transport coefficients are properties of the Green’s function $`F`$, the diffusion propagator. Thus, in order to find an equation for $`f(\epsilon ,s)`$ we should relate the transport coefficients to $`\mathrm{\Pi }`$, the irreducible part of the diffusion propagator. The diffusion coefficient comprises only elastic contributions, however, in the relaxation problem a proper description of inelastic processes is vital. Moreover, as shown in the previuos section, the characteristics for particle and energy transport can already be calculated from the function $`P_L`$, the $`q0`$ limit of the diffusion propagator. Therefore, in establishing a self-consistency equation we should focus on those characteristics that are important in that limit, that is the energy relaxation speed $`v`$.
The equation for the diffusion propagator for $`q=0`$ is given by:
$`sP_L(\epsilon ^{},\epsilon ;s)=\delta \left(\epsilon ^{}\epsilon \right)+P_L(\epsilon ^{},\epsilon ,s){\displaystyle 𝑑\epsilon _1\mathrm{\Pi }(\epsilon _1,\epsilon ;s)}`$ (57)
$`+{\displaystyle \frac{P_L(\epsilon ^{},\epsilon ;s)}{\epsilon }}{\displaystyle 𝑑\epsilon _1\left(\epsilon _1\epsilon \right)\mathrm{\Pi }(\epsilon _1,\epsilon ;s)}.`$ (58)
If we compare Eq.(58) with Eq.(36) we deduce that
$$v(\epsilon ,s)=𝑑\epsilon _1\left(\epsilon _1\epsilon \right)\mathrm{\Pi }(\epsilon _1,\epsilon ;s).$$
(59)
Now we decompose $`v`$, as defined in Eq.(59), into two parts, one part that contains only the effective medium approximation of $`\mathrm{\Pi }`$:
$$v(\epsilon ,s)=𝑑\epsilon _1\left(\epsilon _1\epsilon \right)\mathrm{\Pi }|_{EMA}(\epsilon _1,\epsilon ;s),$$
(60)
and a part $`\delta v(\epsilon ,s)`$, that contains the deviations $`\stackrel{~}{F}`$. Self-consistency requires:
$$\delta v(\epsilon ,s)=0.$$
(61)
$`v(\epsilon )`$, as defined by Eq.(60), is in accordance with the definition (38), taking into account Eq.(B3) and the inequality $`\omega f^{}(\epsilon )/f(\epsilon )1`$ (Eq. (40)), owing to which contributions proportional to this parameter are negligible.
We now focus on the self-consistency equation (61). While Eq.(60) contains only the effective medium contribution to the diffusion propagator, $`\delta v`$ is a functional of the effective medium $`f`$ and the deviation $`\stackrel{~}{F}`$. By construction, it is at least linear in $`\stackrel{~}{F}`$. If this equation could be solved exactly, an exact solution to the diffusion problem could be found within quasi-elastic accuracy. In practise this is not possible, and, therefore, we depend on further approximations. To simplify this equation, we take into account only the lowest order contributions to this equation with respect to $`\stackrel{~}{F}`$, i.e. we linearize $`\delta v`$ with respect to $`\stackrel{~}{F}`$ and require that the first order contribution vanishes. This approach is quite close to the usual CPA-philosophy, in which vanishing of the t-matrix is required in its lowest order approximation. Using this procedure we obtain the following self-consistency equation:
$$\frac{1}{2}\omega ^2𝒩(ϵ)f(ϵ,s)\stackrel{~}{W}(0|ϵ;s)=a\omega \frac{sb}{\stackrel{~}{W}(0|ϵ;s)𝒩(ϵ)},$$
(62)
where $`a`$ and $`b`$ are simply numbers. A detailed derivation of this equation is given in Appendix C.
In deriving the self-consistency equations we have imposed further restrictions on the effective transition probabilities, which determine the range of its applicability. These restrictions can be formulated most conviniently using the dimensionless critical hopping length $`\rho _c(\epsilon ,s)`$, related to $`f(\epsilon ,s)`$ by (see Appendix B):
$$f(\epsilon ,s)\nu =\mathrm{exp}\rho _c(\epsilon ,s).$$
(63)
In terms of $`\rho _c`$, the inequalities (40) and (41), used in the derivation, read (prime is derivative with respect to $`\epsilon `$):
$$\left|\omega \rho _c^{}(\epsilon ,s)\right|1,$$
(64)
$$\left|\rho _c(\epsilon ,0)\rho _c(\epsilon ,s)\right|\rho _c(\epsilon ,0).$$
(65)
In addition, when calculating the integrals,
$$\rho _c(\epsilon ,s)1,$$
(66)
was used.
A closed solution to the self-consistency equation can only be found in the limit $`s=0`$. There we obtain:
$$\rho _c\left(\epsilon ,s=0\right)=\frac{2\alpha }{\left[\omega 𝒩\left(\epsilon \right)\right]^{1/d}}\left[\frac{2da}{S\left(d\right)}\right]^{1/d},$$
(67)
where $`S(d)`$ is the solid angle in $`d`$ dimensions.
For $`s`$ satisfying (65), Eq.(62) can be cast into the form:
$$\left[\rho _c(\epsilon ,0)\rho _c(\epsilon ,s)\right]\mathrm{exp}\left[\rho _c(\epsilon ,0)\rho _c(\epsilon ,s)\right]=\frac{s}{\mathrm{\Omega }\left(\epsilon \right)},$$
(68)
where:
$$\mathrm{\Omega }\left(\epsilon \right)=\frac{2da}{b\omega \rho _c(\epsilon ,0)}v(\epsilon ,0).$$
(69)
According to Eq.(68), the critical hopping length decreases with increasing frequency. Note, that the structure of the equation (68) for the calculation of the dispersion of the critical hopping length, obtained here, agrees completely with that obtained for the critical hopping length in calculating the equilibrium conductivity .
The dispersion of the transport coeffcicients is determined completely by the dispersion of the critical hopping length. For small $`s`$ the frequency dependence preexponetial factors can be ignored, so that from Eq.(68) explicit equations for the transport coefficients can be obtained. They are given by:
$$\frac{D(\epsilon ,s)}{D(\epsilon ,0)}\mathrm{ln}\frac{D(\epsilon ,s)}{D(\epsilon ,0)}=\frac{s}{\mathrm{\Omega }\left(\epsilon \right)},$$
(70)
$$\frac{v(\epsilon ,s)}{v(\epsilon ,0)}\mathrm{ln}\frac{v(\epsilon ,s)}{v(\epsilon ,0)}=\frac{s}{\mathrm{\Omega }\left(\epsilon \right)}.$$
(71)
The formal solution of these equations is given by the Lambert’s W-function $`𝒲(z)`$, defined by the equation $`z=𝒲(z)\mathrm{exp}𝒲(z)`$. Using the Lambert’s function we can write:
$$D(\epsilon ,s)=D(\epsilon ,0)\mathrm{exp}𝒲\left(s/\mathrm{\Omega }\left(\epsilon \right)\right),$$
(72)
$$v(\epsilon ,s)=v(\epsilon ,0)\mathrm{exp}𝒲\left(s/\mathrm{\Omega }\left(\epsilon \right)\right).$$
(73)
## VIII Constant density of states
A constant density of states, although of not much physical relevance, gives us the unique opportiunity to study pure dispersive energy transport. Here both $`\mathrm{\Omega }(\epsilon )`$ and $`v_0=v(\epsilon ,0)`$ are independent of energy. Consequently, Eq.(44) can readly be integrated. The integration yields:
$$P_L(\epsilon _0\epsilon ,s)=\frac{\mathrm{\Theta }\left(\epsilon _0\epsilon \right)}{v\left(s\right)}\mathrm{exp}\left[\frac{s\left(\epsilon _0\epsilon \right)}{v\left(s\right)}\right].$$
(74)
The time-dependence of this function can be obtained by the inverse Laplace-transformation. Using Eq.(71) to change the integration variable from $`s`$ to $`y=v/v_0`$, the inverse Laplace-transform of Eq.(74) may be written as:
$$P_L(\epsilon _0\epsilon ,t)=\frac{\mathrm{\Omega }}{v_0}\mathrm{\Theta }\left(\epsilon _0\epsilon \right)_C\frac{dy}{2\pi i}\frac{\mathrm{ln}y+1}{y}\mathrm{exp}\left[\mathrm{\Omega }ty\mathrm{ln}y\frac{\mathrm{\Omega }}{v_0}\left(\epsilon _0\epsilon \right)\mathrm{ln}y\right],$$
(75)
with the properly chosen integration contour $`C`$. At large enough $`t`$, that is for $`\mathrm{\Omega }t1`$, this expression, using the saddle-point method, gives simply the Gaussian packet:
$$P_L(\epsilon _0\epsilon ,t)\mathrm{\Theta }\left(\epsilon _0\epsilon \right)\frac{\mathrm{\Omega }}{v_0}\frac{1}{2\sqrt{\pi \mathrm{\Omega }t}}\mathrm{exp}\left[\frac{\left(\epsilon _0\epsilon v_0t\right)^2}{4v_0^2t/\mathrm{\Omega }}\right].$$
(76)
Now let us consider the time dependence of the energy relaxation. According to Eq.(51) the velocity of energy relaxation is:
$$\frac{dE(t)}{dt}=_C\frac{ds}{2\pi i}\frac{v\left(s\right)}{s}e^{st}.$$
(77)
Again, the time dependence of the energy relaxation speed can be calculated using asymptotics. Details of the calculations are presented in Appendix D, where it is shown that:
$$\frac{dE\left(t\right)}{dt}v_0\{\begin{array}{c}1+\sqrt{\frac{e}{2\pi }}\frac{1}{\left(\mathrm{\Omega }t\right)^{3/2}}\mathrm{exp}\left(\mathrm{\Omega }t/e\right),\mathrm{as}\mathrm{\Omega }t1,\hfill \\ \frac{\sqrt{2\pi }}{e\mathrm{\Omega }t}\left[\mathrm{ln}\left(e/(\mathrm{\Omega }t)\right)\right]^2,\mathrm{as}\mathrm{\Omega }t1.\hfill \end{array}$$
(78)
Note, that the problem of the energy relaxation in the case of constant density of states ($`v(\epsilon ,s)`$ is independent of $`\epsilon `$) is completely equivalent to one of the nonmarkovian charge transport in strong electric field $`E`$, when the diffusion (described by the second coordinate derivate) is totally neglected. One has only to replace $`v\left(s\right)u\left(s\right)E`$, $`u\left(s\right)`$ being the mobility. In the context of dispersive particle transport, the regimes $`\mathrm{\Omega }t1`$ and $`\mathrm{\Omega }t1`$ are the regimes of anomalous and normal “diffusion”, respectively. The main difference between both is in that while the sites are ususally distributed homogeneously in space, the density of states is usually an increasing function of energy.
## IX Exponential density of states
### A The saddle-point approximation and its break down for large times
The calculation of the time dependence of the energy distribution function for an arbitrary density of states on the basis of the Eqs. (36)-(38), (70) and (71) turns out to be a quite intricate problem. A tool that can be utilized in tackeling the problem is the saddle point approximation. How to apply the saddle point approximation for the calculation of the quantities of interest for an arbitrary density of states is shown in Appendix E. Below we focus on the exponential density of states, which is relevant, e.g., for amorphous Si:H.
We assume that the density of states is given by:
$$𝒩(\epsilon )=𝒩_0\mathrm{exp}\left(3\frac{\epsilon }{\mathrm{\Delta }}\right).$$
(79)
To simplify the notations we use the abbreviations:
$$\overline{\omega }=\frac{b\omega }{2da},$$
(80)
$$\overline{\nu }=\frac{2da^2}{b}\nu .$$
(81)
Then, the Eqs.(B3), (67) and (69) can be cast into the form:
$$v(\epsilon ,0)v\left(\epsilon \right)=\overline{\omega }\overline{\nu }\mathrm{exp}\left[\rho \left(\epsilon \right)\right],$$
(82)
$$\mathrm{\Omega }\left(\epsilon \right)=\frac{\overline{\nu }}{\rho \left(\epsilon \right)}\mathrm{exp}\left[\rho \left(\epsilon \right)\right],$$
(83)
$$\rho \left(\epsilon \right)=\rho _c(\epsilon ,0)=A\mathrm{exp}\left(\frac{\epsilon }{\mathrm{\Delta }}\right),$$
(84)
where:
$$A=\frac{2\alpha }{\left[\omega 𝒩\left(0\right)\right]^{1/d}}\left[\frac{2da}{S\left(d\right)}\right]^{1/d}.$$
(85)
Using these equations it follows from the formulas derived in Appendix E that the time dependence of the mean squared deviation and the time dependence of the mean energy are given by
$$\sigma ^2(\epsilon _m,\epsilon _0)\sigma ^2(t,\epsilon _0)\overline{\omega }\mathrm{\Delta }\left[1\frac{t_0^2}{\left(t+t_0\right)^2}\right],$$
(86)
$$\epsilon _m\left(t\right)\mathrm{\Delta }\mathrm{ln}\left[\frac{1}{A}\mathrm{ln}\left(\frac{\overline{\omega }}{\mathrm{\Delta }}\overline{\nu }\left(t+t_0\right)\right)\right],$$
(87)
where
$$t_0\left(\epsilon _0\right)=\frac{\mathrm{\Delta }}{\overline{\nu }\rho \left(\epsilon _0\right)v\left(\epsilon _0\right)}.$$
(88)
Note, that, according to the saddle-point approximation, the distribution is Gaussian. Furthermore, the dispersion is constant for $`tt_0`$. Consequently, for time scales in line with the applicability of the saddle-point approximation, the motion of the energy packet is “soliton-like”, that is the packet moves without distortion.
The applicability condition for the saddle-point (E13) requires:
$$\frac{\overline{\omega }}{\mathrm{\Delta }}\rho ^2(\epsilon _m)1.$$
(89)
Because of $`\rho (\epsilon _m)`$ grows with $`t`$, the saddle-point method, and, consequently, all the results of this chapter, becomes invalid, at least at sufficiently large $`t`$. The physical meaning of the condition (89) is the following: at $`tt_0`$ the width of the electron’s energies distribution is $`\delta \epsilon =\sqrt{\overline{\omega }\mathrm{\Delta }}`$. Because of $`\rho \mathrm{exp}\left(\epsilon /\mathrm{\Delta }\right)`$, the variation of $`\rho `$ is: $`\delta \rho /\rho =\sqrt{\overline{\omega }/\mathrm{\Delta }}`$, and, because $`v\mathrm{exp}\left(\rho \right)`$, we have: $`\delta v/v=\delta \rho =\sqrt{\overline{\omega }\rho ^2/\mathrm{\Delta }}`$. Thus, the condition (89) is just the one of the variation of electron’s “velocity” across the distribution is smaller, then the velocity itself.
### B Form of the distribution for large times
Even if the initial condition are such, that the condition (89) is fulfilled, and a gaussian distribution is formed, at some moment of time $`\rho _m\rho (\epsilon _m)`$ becomes of the order of $`\sqrt{\mathrm{\Delta }/\overline{\omega }}`$, which definitely should results in some deviation of the distribution function from its symmetric, gaussian form. For the exponential density of states, distribution function in the Laplace representation, Eq.(44), using Eqs.(73,82-84), may be written as:
$$P_L(\epsilon _0,\epsilon ;s)=\frac{1}{s}\rho \frac{}{\rho }\mathrm{exp}\left[\frac{1}{\overline{\omega }}\underset{\rho _0}{\overset{\rho }{}}𝑑x\frac{𝒲\left(sxe^x\right)}{x^2}\right]=\frac{𝒲\left(s\rho e^\rho \right)}{\overline{\omega }\rho s}\mathrm{exp}\left[\frac{1}{\overline{\omega }}\underset{\rho _0}{\overset{\rho }{}}𝑑x\frac{𝒲\left(sxe^x\right)}{x^2}\right],$$
(90)
where it was set $`\overline{\nu }=1`$ and $`\mathrm{\Delta }=1`$ by appropriate choice of time and energy units. Under the same conditions, which were used in the derivation of the above formula, the following approximation is valid (see Appendix F):
$$\underset{\rho _0}{\overset{\rho }{}}𝑑x\frac{𝒲\left(sxe^x\right)}{x^2}G(s,\rho )G(s,\rho _0);G(s,\rho )=\frac{1}{\rho ^2}𝒲\left(s\rho e^\rho \right)\left[1+\frac{1}{2}𝒲\left(s\rho e^\rho \right)\right]+O(\frac{𝒲}{\rho ^3},\frac{𝒲^3}{\rho ^3}).$$
(91)
When calculating distribution function at sufficiently large times, the characteristic values of $`s`$, giving the main contribution into inverse Laplace integral become so small, that one can suppose $`\left|s\right|\rho _0\mathrm{exp}(\rho _0)1`$, and: $`G(s,\rho _0)s\rho _0^1\mathrm{exp}\left(\rho _0\right)`$. Under this condition the distribution function approach initial conditions independent shape: $`P(\epsilon _0,\epsilon ;t)\varphi (\epsilon ,t+t_0(\epsilon _0))`$, $`t_0(\epsilon _0)=\left(\overline{\omega }\rho _0\right)^1\mathrm{exp}\rho _0`$,
$$\varphi (\epsilon ,t)=\underset{i\mathrm{}}{\overset{+i\mathrm{}}{}}\frac{ds}{2\pi i}\varphi _L(\epsilon ,s)e^{st}=\underset{i\mathrm{}}{\overset{+i\mathrm{}}{}}\frac{ds}{2\pi is}\rho \frac{}{\rho }\mathrm{exp}\left\{st\frac{1}{\overline{\omega }\rho ^2}𝒲\left(s\rho e^\rho \right)\left[1+\frac{1}{2}𝒲\left(s\rho e^\rho \right)\right]\right\}.$$
(92)
Let us consider, first, momenta of the distribution at the large times. In Laplace-representation the momenta are defined by the equation:
$$\chi _{nL}\left(s\right)=\underset{0}{\overset{\mathrm{}}{}}𝑑te^{st}\epsilon ^n\left(t\right)=\underset{\rho _0}{\overset{\mathrm{}}{}}\frac{d\rho }{\rho }\left(\mathrm{ln}\rho \right)^nP_L(\epsilon _0,\epsilon ,s)e^{st_0}\frac{d\rho }{\rho }\left(\mathrm{ln}\rho \right)^n\varphi _L(\epsilon ,s).$$
(93)
We shall omit further on irrelevant time shift multiple $`e^{st_0}`$. It turns out, that for $`\overline{\omega }\rho 1`$, the integrand of Eq.(93) has a maximum, located at:
$`\rho _m\left(1+\sqrt{\overline{\omega }}\right)\mathrm{ln}{\displaystyle \frac{\sqrt{\overline{\omega }}}{s}}.`$
However, if one would try to use the saddle-point method in evaluation of integrals like (90), the results would be incorrect at $`\sqrt{\overline{\omega }}\rho _m1`$, since the saddle point criterion is not fulfilled. However, at $`\overline{\omega }\rho ^21`$, momenta can be calculated as an expansion on the powers of small paramenter $`1/\left(\sqrt{\overline{\omega }}\rho \right)`$. Details of the calculations are presented in Appendix F.
As a result, we have for the first momentum, or the mean energy:
$`\chi _1\left(t\right)=\epsilon \left(t\right)\mathrm{\Delta }\mathrm{ln}{\displaystyle \frac{1}{A}}\mathrm{ln}\left(\stackrel{~}{b}\sqrt{\overline{\omega }}t\right)`$ (94)
$`\sqrt{{\displaystyle \frac{\pi \overline{\omega }}{2}}}+\sqrt{{\displaystyle \frac{\pi }{8\overline{\omega }}}}{\displaystyle \frac{1}{\mathrm{ln}^2\left(\stackrel{~}{b}\sqrt{\overline{\omega }}t\right)}}+\mathrm{},`$ (95)
where $`\stackrel{~}{b}=\sqrt{2}e^{1+\gamma /2}`$, and $`\gamma `$ is Eulers’s constant. For the distribution’s dispersion $`\sigma ^2(t)=\mu _2(t)=\chi _2(t)\chi _1^2(t)`$ we obtain:
$$\sigma ^2\left(t\right)=\mu _2\left(t\right)=\left(2\frac{\pi }{2}\right)\overline{\omega }+\left(1+\mathrm{ln}2\right)\frac{\sqrt{2\pi \overline{\omega }}}{\mathrm{ln}\left(\sqrt{\overline{\omega }}t\right)}+\mathrm{},$$
(96)
and for the third central momentum $`\mu _3(t)=\left(\epsilon \epsilon \left(t\right)\right)^3=\chi _3(t)3\chi _1\left(t\right)\chi _2\left(t\right)+2\chi _1^3\left(t\right)`$ we have:
$$\mu _3\left(t\right)=\left(\pi 3\right)\sqrt{\frac{\pi }{2}}\overline{\omega }^{3/2}+3\left[\pi \left(1\mathrm{ln}2\right)1\right]\frac{\overline{\omega }}{\mathrm{ln}\left(\sqrt{\overline{\omega }}t\right)},$$
(97)
Figures 5, 6, and 7 show mean energy, $`\chi _1`$, mean square deviation of energy $`\sigma `$, and dimensionless coefficient of asymmetry:
$$\mathrm{Asymmetry}=\frac{\mu _3\left(t\right)}{\sigma ^3\left(t\right)},$$
(98)
respectively, plotted versus $`\mathrm{ln}\mathrm{ln}t`$. There values, obtained numerically, using the distribution function (92), are compared with their corresponding asymptotic forms, following from Eqs.(95)-(97).
Moving along the same line, the whole distribution function may be reconstructed at large times, if to calculate all its momenta. However, it may be obtained in more simple way. Namely, when $`\overline{\omega }\rho ^21`$, one may replace $`s`$ within arguments of $`𝒲`$-functions in Eq.(92) with some $`s_0=c/t`$, $`c1`$ (a kind of argumentation may be found in Appendix G). Then the integral can easily be calculated to yield:
$`\varphi (\epsilon ,t)=`$ (99)
$``$ $`\rho {\displaystyle \frac{}{\rho }}\mathrm{exp}\left\{{\displaystyle \frac{1}{\overline{\omega }\rho ^2}}𝒲\left({\displaystyle \frac{c}{t}}\rho e^\rho \right)\left[1+{\displaystyle \frac{1}{2}}𝒲\left({\displaystyle \frac{c}{t}}\rho e^\rho \right)\right]\right\}.`$ (100)
One can easily show (see Appendix G), that asymptotic expressions for momenta (95-97) may be entirely reproduced from Eq.(100), if to set $`c=\mathrm{exp}(\gamma )`$ ($`\gamma `$ is the Euler’s constant). Figures 8-10 show some plots of the distribution function (92), compared with its asymptotic form (100). This latter can be further symplified, if to introduce $`\rho _t`$:
$$\frac{t}{c}=\rho _te^{\rho _t};\rho _t=𝒲\left(\frac{t}{c}\right).$$
(101)
As we shall see, the body of the distribution at a given time corresponds to values of $`\rho >\rho _t`$, $`\rho \rho _t\rho _t`$. Then, one can use approximation: $`𝒲\left(\left(c/t\right)\rho e^\rho \right)\rho \rho _t1`$. Then, we have in the exponent $`\left(1\rho _t/\rho \right)^2/2\overline{\omega }`$, which can be replaced by $`\left(\epsilon \epsilon _t\right)^2/2\overline{\omega }`$, where $`\epsilon _t\mathrm{ln}\mathrm{ln}t`$. As a result, at large time and small $`\overline{\omega }`$ the distribution function turns into:
$$\varphi (\epsilon ,t)=\frac{}{\epsilon }\{\begin{array}{c}\mathrm{exp}\left[\frac{\left(\epsilon \epsilon _t\right)^2}{2\overline{\omega }}\right],\epsilon <\epsilon _t\\ 1,\epsilon >\epsilon _t\end{array};\epsilon _t=\mathrm{ln}\mathrm{ln}t.$$
(102)
### C Initial stage of the evolution
At $`t=0`$ it was supposed $`P(\epsilon ,\epsilon _0;t=0)=\delta (\epsilon \epsilon _0)`$. To understand, what happens at the intitial stage of the distribution’s evolution, let us consider the behaviour of $`P`$ at $`\epsilon \epsilon _0`$. One can write from Eqs.(E1,73,82-84):
$`P_0(\epsilon _0,t)`$ $``$ $`P(\epsilon _0,\epsilon _0;t)={\displaystyle \frac{1}{\overline{\omega }\rho _0}}{\displaystyle \underset{i\mathrm{}}{\overset{+\mathrm{}}{}}}{\displaystyle \frac{ds}{2\pi is}}𝒲\left({\displaystyle \frac{s}{\overline{\nu }}}\rho _0e^{\rho _0}\right)e^{st}`$ (103)
$`=`$ $`{\displaystyle \frac{1}{\overline{\omega }\rho _0}}{\displaystyle _C}{\displaystyle \frac{dz}{2\pi i}}\left(1+z\right)\mathrm{exp}\left(\tau ze^z\right)`$ (104)
$`=`$ $`{\displaystyle \frac{1}{\overline{\omega }\rho _0\tau }}{\displaystyle _C}{\displaystyle \frac{dz}{2\pi i}}\mathrm{exp}\left(\tau ze^zz\right);`$ (106)
$`\tau ={\displaystyle \frac{\overline{\nu }t}{\rho _0}}e^{\rho _0},\rho _0=A\mathrm{exp}\left({\displaystyle \frac{\epsilon _0}{\mathrm{\Delta }}}\right),`$
where the new integration variable $`z=𝒲\left(\left(s/\overline{\nu }\right)\rho _0e^{\rho _0}\right)`$ was introduced instead of $`s`$, and the integration by parts was performed. This integral may be evaluated using the saddle-point method. The saddle-point equation, $`f^{}(z)=\tau (z+1)e^z1=0`$, gives us saddle point value $`z_c=𝒲\left(e/\tau \right)1`$, that is $`z_c\mathrm{ln}(1/\tau )\mathrm{ln}\mathrm{ln}(1/\tau )`$ at $`\tau 1`$, and $`z_ce/\tau 1`$ as $`\tau 1`$. To ensure the correctness of the saddle point approximation, one should require the parameter$`\left|f^{\prime \prime \prime }(z_c)\right|^2/\left|f^{\prime \prime }(z_c)\right|^3=\left|(3+z_c)^2/(2+z_c)^3\right|`$ to be small. While this is true at $`\tau 1`$, this parameter appears to be $`1`$ at $`\tau 1`$. Therefore, the result for $`\tau 1`$ is correct up to the multiple of the order of 1 only. We have:
$$P_0(\rho ,t)\{\begin{array}{c}\frac{1}{\sqrt{2\pi }\overline{\omega }\rho }\left(\frac{e}{\tau }\right)^{3/2}\mathrm{exp}\left(\frac{\tau }{e}\right)\text{ as }\tau 1,\\ \frac{e}{\sqrt{2\pi }\overline{\omega }\rho }\left(\mathrm{ln}\frac{e}{\tau }\mathrm{ln}\mathrm{ln}\frac{e}{\tau }\right)\text{as }\tau 1.\end{array}$$
(107)
In a similar fasion one can calculate the energy derivative at $`\epsilon =\epsilon _0`$:
$`P_1(\rho _0,t){\displaystyle \frac{P}{\epsilon }}|_{\epsilon =\epsilon _0}={\displaystyle \frac{\rho }{\mathrm{\Delta }}}{\displaystyle \frac{}{\rho }}P(\rho _0,\rho ;t)|_{\rho =\rho _0}`$ (108)
$`=`$ $`{\displaystyle \frac{1}{\overline{\omega }\mathrm{\Delta }\rho _0}}{\displaystyle \underset{i\mathrm{}}{\overset{+\mathrm{}}{}}}{\displaystyle \frac{ds}{2\pi is}}𝒲\left({\displaystyle \frac{\mathrm{\Delta }𝒲}{\overline{\omega }\rho _0}}+{\displaystyle \frac{𝒲\rho _0}{1+𝒲}}\right)e^{st}`$ (109)
$``$ $`{\displaystyle \frac{1}{\left(\overline{\omega }\rho _0\right)^2}}{\displaystyle _C}{\displaystyle \frac{dz}{2\pi i}}\left[z\left(1+z\right){\displaystyle \frac{\overline{\omega }}{\mathrm{\Delta }}}\rho _0^2\right]\mathrm{exp}\left(\tau ze^z\right),`$ (110)
where the term $`(\overline{\omega }/\mathrm{\Delta })\rho _0z`$ within the square brackets was neglected due to assumption $`\left(\overline{\omega }/\mathrm{\Delta }\right)\rho 1`$. At $`\tau 1`$ this integral may be readily evaluated in the saddle point approximation, which yields:
$$P_1(\rho ,t)\frac{1}{\overline{\omega }\mathrm{\Delta }}\sqrt{\frac{e}{2\pi \tau }}\mathrm{exp}\left(\frac{\tau }{e}\right).$$
(111)
When $`\tau 1`$, one have to perform integration by parts first, which gives:
$`P_1(\rho ,t)={\displaystyle \frac{1}{\overline{\omega }^2\rho ^2\tau }}`$
$`{\displaystyle _C}{\displaystyle \frac{dz}{2\pi i}}\left[z1{\displaystyle \frac{\overline{\omega }\rho ^2}{\mathrm{\Delta }\left(1+z\right)}}{\displaystyle \frac{\overline{\omega }\rho ^2}{\mathrm{\Delta }\left(1+z\right)^2}}\right]\mathrm{exp}\left(\tau ze^zz\right).`$
Then one have in the saddle point approximation:
$`P_1(\rho ,t)`$ $``$ $`{\displaystyle \frac{e}{\sqrt{2\pi }\overline{\omega }^2\rho ^2}}\left(z_c^21\overline{\omega }\rho ^2\right),`$ (112)
$`z_c`$ $``$ $`\mathrm{ln}{\displaystyle \frac{1}{\tau }}\mathrm{ln}\mathrm{ln}{\displaystyle \frac{1}{\tau }}.`$ (113)
At sufficiently small $`t`$ the above expression is positive, which means that the distribution is “sticked” near $`\epsilon =\epsilon _0`$, that is, it monotoneously decreases as $`\epsilon _0\epsilon >0`$ increases. At some $`\tau =\tau _0`$, $`P_1`$ changes its sign. If $`\left(\overline{\omega }/\mathrm{\Delta }\right)\rho _0^21`$, this corresponds to $`z_{c0}\mathrm{ln}\frac{1}{\tau _0}\sqrt{\overline{\omega }/\mathrm{\Delta }}\rho _0`$, otherwise $`\tau _01`$. This correspons just to the moment of time, when this distribution separates from the intial point, — its maximum is at $`\epsilon <\epsilon _0`$. For $`\tau \tau _0`$ the width of the distribution on the energy scale may be estimated as:
$$\mathrm{\Delta }\epsilon (t)\frac{P_0(\rho ,t)}{P_1(\rho ,t)}\frac{\overline{\omega }\rho }{z_c}\frac{\overline{\omega }\rho }{\mathrm{ln}(1/\tau )}.$$
(114)
This formula is valid if $`\mathrm{ln}(1/\tau )\rho `$, up to the times corresponding either $`\tau 1`$ if $`\left(\overline{\omega }/\mathrm{\Delta }\right)\rho _0^21`$, or $`\mathrm{ln}\left(1/\tau \right)\sqrt{\overline{\omega }/\mathrm{\Delta }}\rho _0`$ otherwise. In the former case the distribution width becomes of the order of $`\overline{\omega }\rho `$ just before it separates from the initial point. After this, the distribution, as it was shown in previous subsection, widens till $`\sqrt{\overline{\omega }\mathrm{\Delta }}`$, and the gaussian packet moves downward, untill its center reaches the value, corresponding to $`\rho \sqrt{\mathrm{\Delta }/\overline{\omega }}`$. In the latter case the packet’s width is $`\sqrt{\overline{\omega }\mathrm{\Delta }}\overline{\omega }\rho `$ at the very moment of its “unsticking”. In both cases, to consider subsequent evolution, one has to use some other approximation instead of saddle point one.
## X Conclusions
In the paper we have presented an effective method which permitts the investigation of relaxation phenomena of localized charge carriers far from equillibrium due to phonon-assisted hopping at zero temperature. From the point of view of the formalism the main equations are the equations for the calculation of the diffusion propagator (36), (37) and (38), and the equations (69), (70) and (71), which determine the dispersion of the transport coefficients. These equations show, that both particle transport and energy transport are dispersive. The equations (70) and (71) lead to a strong dependence of the diffusion constant $`D(ϵ;s)`$ and the velocity of energy relaxation $`v(ϵ;s)`$ on frequency $`s`$ already for low frequencies. To our knowledge, these equations have not been dervied in the literature so far, for systems far from equillibrium. In fact, in the literature mainly frequency independent transport coefficients can be found (see e.g. and references therein). The strong dependence of the transport coefficients on $`s`$ results in a non-Markovian equation for the calculation of the diffusion propagator (36). This discriminates our equations from the Markovian integral equations used in the Refs. -. The latter equations can, in principle, be obtained from our leading equation (11) by neglecting statistical correlation, so as to average every factor independently.
Using our effective-medium method, we have investigated the relaxation of charge carriers in band tails at $`T=0`$. According to our results, energy relaxation is connected with dispersional transport. Even if we have no real diffusion in the system, since the temperature is zero, we have some spreading of the energy distribution with time. For a constant density of states the situation is completely equivalent to the one in the problem of the electron’s motion in a disordered system, subjected to an electric field . The main difference between the above mentioned problem and the energy relaxation is in that, while in the former case the sites are distributed homogeneously in space, in the energy relaxation problem the density of states is a decaying function of energy for most physical systems. In this situation we arrive at the picture that the particles of a package with lower energies move slower than particles with higher energy. This leads to the opposite tendency: at first, there is a slowing of the package spreading. The time dependence of the dispersion, being linear at the first stage of the evolution, is later slowing down. Later on, different possibilities exist, dependend on the particular type of the energy dependence of the density of states.
We performed a detailed investigation for the exponential density of states for two time regimes. If the variation of the velocity $`v`$ in the energy space across the distribution is smaller than the velocity itself, that is if $`\omega \rho ^2/\mathrm{\Delta }1`$, we arrive at the situation of a package, moving steadily down along the energy axis without any deformation. Dispersion becomes time independent (see Eq.(86)). However, since the particles are sinking down this condition becomes violated when time goes by. The parameter $`\omega \rho ^2/\mathrm{\Delta }`$, being small at the first stage of the evolution, is getting large. When this parameter is larger than one, the steady motion condition becomes violated again. The package, being of gaussian form before, undergoes some restructuring to another, non-gaussian stable form, with its width lower than before by some numeric factor of the order of 1 (Eq.(96)). In our pictures, Figs 7-10, it is clearly seen, that the package becomes non-gaussian. This result remains valid until the very moment the quasi-elasticity conditions break down ($`\omega \rho /\mathrm{\Delta }1`$).
For the exponential density of states we have found the packet to move as $`\mathrm{ln}\mathrm{ln}t`$, i.e., its motion is strongly slowing down with time, the packet becomes, roughly speaking, almost stopped. This type of behaviour may be called “glassy”, because of the overall time scale for the packet evolution (governed by the exponential function of the large parameter $`\mathrm{\Delta }/\omega `$), becomes huge. For example, when $`\omega /\mathrm{\Delta }=10^4`$, this time scale approaches $`10^{1000}`$ and more (see Fig. 10)!
The main simplification used in our paper is the quasi-elastic approximation. This approximation relies on the smallness of the upper bound of the energy transfer to the phonon system by one hop. For localized electrons, this upper bound can be much smaller as the Debye energy of the host material, since not all phonons can interact with localized electrons equally well. For localized electrons the electron phonon coupling constant approaches zero, for phonons with wave-vector $`q>2\alpha `$. Thus high energetic phonons are less effective. Only phonons with energies $`\omega <\omega _D2\alpha a`$, where $`a`$ is the lattice constant of the host material, are effective. Furthermore, in disordered systems the high energetic acustical phonons are localized, and thus need not contribute to transport by necessity. Nevertheless, the question whether the quasi-elastic approximation is applicable depends much on the material of interest. If, however, we compare our result for the time dependence of the mean energy (87) in an exponential density of states with those existing in the literature , we also find agreement for $`\omega /\mathrm{\Delta }1`$, which indicates that our results are, at least qualitatively, of wider validity. Unfortunately, to our knowledge, in the literature there are no further results on the width of the energy distribution available to compare with.
From our point of view the main question remained open is how the results, obtained for the exponential density of states, may be generalized for other types of energy dependencies. One can imagine, e.g., that for densities of states decaying with decreasing energy slower than exponential, the dispersion is growing, but under some sublinear law. Also, in the opposite case, in a density of states decaying faster than exponential, we would possibly have a dispersion tending to zero for large times.
It should, however, also be mentioned that it is not completely clear under which situation zero temperature results can be applied to systems at finite temperature. In our approach we assumed that $`T=0`$. However, in a real system, when the carriers are sinking down, the criterion for the temperature to be treated as zero, is violated at every finite value of $`T`$ at some moment of time, since the contribution of hops to sites with higher energy values becomes more and more comparable with one of downward hops. Therefore, the consideration of temperature becomes vital. Work in this direction is in progress.
## ACKNOWLEDGMENTS
On of the authors (O.B.) would like to thank Prof. P. Thomas for drawing his attention to the photoluminscense problem and stimulating discussions.
## A On the configuration average
Here we derive the set of integral equations (17)-(22). The same method has already been used in Ref..
According to Eq.(10), the configuration average of the electron density $`n(\rho ,s)`$ is given by:
$$n(\rho ,s)=𝑑\rho _1p_0\left(\rho _1\right)P(\rho _1,\rho ),$$
(A1)
where
$$P(\rho ^{},\rho )=\eta \left(\rho ^{}\right)\mathrm{\Phi }(\rho ^{},\rho )$$
(A2)
is the diffusion function. According to Eq.(11), the diffusion function is given by:
$$P=\frac{1}{s}\underset{n=0}{\overset{\mathrm{}}{}}\eta \left(\frac{V}{s}\right)^n.$$
(A3)
Diagrammatically this series can be represented as depicted in Fig.1. In this picture every full dot represents a potential $`V`$, which depends on disorder via the structur factor $`\eta `$. The expansion starts with a single structural factor $`\eta `$, depicted by the empty dot. Note that all disorder is comprised into the structure factor $`\eta `$.
In order to calculate the configuration average, we average the series (A3) term by term. The resulting expansion can be depicted in the usual form, as in impurity scattering problem (see Fig.2). One can see from the picture, that the set of all diagrams can be decomposed into two subsets, the set $`S`$ containing all diagrams connected with the empty point and the set $`F`$ containing all other diagrams. Owing to this decomposition configuration averaged diffusion function can be written as:
$$P(\rho ^{},\rho )=𝑑\rho _1S(\rho ^{},\rho _1)F(\rho _1,\rho ),$$
(A4)
from which Eq.(17) follows. Note that $`F`$ is nothing but $`\mathrm{\Phi }`$.
Again the class of diagrams contributing to $`F`$ can be decomposed into reducible and irreducible diagrams. All irreducible diagrams can be comprised into a function $`\mathrm{\Pi }`$. Doing so, we obtain Eq.(18). Note, that, although the introduction of the irreducible part $`\mathrm{\Pi }`$ is parallel to the introduction of the self-energy it has to be stressed, that its physical interpretation is completely different. In the present context $`\mathrm{\Pi }(\rho ^{},\rho )`$ is the true transition probability from $`\rho ^{}`$ to $`\rho `$. It is $`\mathrm{\Pi }(\rho ^{},\rho )`$ that has to be calculated and compared with the experimental situation and not the bar transition probability, as suggested in many papers (see e.g. ). Diagrams, contributing to $`\mathrm{\Pi }`$ are depicted in Fig.3. The set of diagrams depicted in Fig. 3 can be generated by means of the propagator $`\mathrm{\Pi }_{\stackrel{~}{\rho }}`$, defined by the equation (21), using Eq.(19).
The propagator $`\mathrm{\Pi }_{\stackrel{~}{\rho }}`$ can also be used to generate the irreducible block $`S`$. Diagrams contributing to $`S`$ are depicted in Fig.4. From the picture it follows, that $`S`$ is given by Eq.(22).
## B Transport coefficients
According to Eq.(38), $`v(\epsilon ,s)`$ is given by:
$$v(\epsilon ,s)=\frac{1}{2}\omega ^2𝒩\left(\epsilon \right)\stackrel{~}{W}\left(q=0|\epsilon \right)=\frac{1}{2}\omega ^2𝒩\left(\epsilon \right)\nu \mathrm{exp}\left[\rho _c(\epsilon ,s)\right]𝑑𝐑\frac{1}{1+\mathrm{exp}\left[2\alpha R\rho _c(\epsilon ,s)\right]},$$
(B1)
where:
$$f(\epsilon ,s)\nu =\mathrm{exp}\rho _c(\epsilon ,s).$$
(B2)
We assume, that $`\rho _c`$ is large, so the integrals can be calculated in the limit $`\rho _c\mathrm{}`$. This assumption is justified a posteriory by explicit calculation of $`\rho _c`$. In limit $`\rho _c\mathrm{}`$, the integrand reduces to a step function, so that the integrals can be calculated easily. Doing so, we find:
$$v(\epsilon ,s)=\frac{1}{2}\frac{S\left(d\right)}{d}\omega ^2𝒩\left(\epsilon \right)\left[\frac{\rho _c(\epsilon ,s)}{2\alpha }\right]^d\nu \mathrm{exp}\left[\rho _c(\epsilon ,s)\right],$$
(B3)
where $`S(d)`$ is the solid angle and $`d`$ is the spatial dimension. From the calculation and the result it is to be seen, that $`\rho _c(\epsilon ,s)`$ is to be identified with the dimensionless critical hopping length at energy $`\epsilon `$.
The same procedure is applied to the calculation of $`D(\epsilon ,s)`$, defined by Eq.(37). Using the same approximations, we obtain:
$$D(\epsilon ,s)=\frac{1}{2}\frac{S\left(d\right)}{d\left(d+2\right)}\omega 𝒩\left(\epsilon \right)\left[\frac{\rho _c(\epsilon ,s)}{2\alpha }\right]^{d+2}\nu \mathrm{exp}\left[\rho _c(\epsilon ,s)\right].$$
(B4)
Thus $`D(\epsilon ,s)`$ and $`v(\epsilon ,s)`$ differ only in preexponential factors.
We anticipate that for $`s=0`$ we obtain:
$$\rho _c(\epsilon ,0)=\frac{2\alpha }{\left[\omega 𝒩\left(\epsilon \right)\right]^{1/d}}\left[\frac{2da}{S\left(d\right)}\right]^{1/d},$$
(B5)
where $`a`$ is a number and $`S(d)`$ is the solid angle. A different form for the diffusion coefficient can be found in Ref.. The expression given in this paper differs from our expression essentially only in that $`\omega `$ is replaced by the tailing parameter. In looking on this difference one should, however, take into account, that the model used in that paper is different. Furthermore, in Ref. the transport coefficient are independent of $`s`$. It is therefore not completely clear to which range of time they apply. Clearly, they can not determine the evolution for all times.
## C Derivation of self-consistency equation
Here we derive the self-consistency equation (62). In order to work out explicitely the first order contribution to the self-consistency equation (61) we first need the first order correction to $`\mathrm{\Pi }`$. It is given by:
$$\mathrm{\Pi }^{(1)}(\rho ^{},\rho ;s)=𝑑\rho _1𝑑\rho _2𝑑\rho _3𝒩\left(\epsilon _3\right)\stackrel{~}{w}_{\rho _3}(\rho ^{},\rho _1;s)\stackrel{~}{F}(\rho _1,\rho _2)\stackrel{~}{w}_{\rho _3}(\rho _2,\rho ;s).$$
(C1)
If we insert the expression for $`\stackrel{~}{w}_{\rho _3}`$, we obtain:
$`\mathrm{\Pi }^{(1)}(\rho ,\rho ;s)={\displaystyle }d\rho _1`$ (C2)
$`\{𝒩\left(\epsilon \right)\stackrel{~}{W}(\rho ,\rho ;s)[\stackrel{~}{F}(\rho ,\rho _1)\stackrel{~}{F}(\rho ,\rho _1)]\stackrel{~}{W}(\rho _1,\rho ;s)𝒩\left(\epsilon _1\right)\stackrel{~}{W}(\rho ^{},\rho _1;s)[\stackrel{~}{F}(\rho _1,\rho )\stackrel{~}{F}(\rho ^{},\rho )]\stackrel{~}{W}(\rho ,\rho _1;s)\}.`$ (C3)
Thus, at this stage the self-consistency equation takes the form:
$`0`$ $`=`$ $`{\displaystyle }d𝐪d\epsilon d\epsilon _1(\epsilon \epsilon )\{𝒩\left(\epsilon \right)\stackrel{~}{W}(0|\epsilon ,\epsilon ;s)\stackrel{~}{F}(q|\epsilon ,\epsilon _1)\stackrel{~}{W}(q|\epsilon _1,\epsilon ;s)𝒩\left(\epsilon \right)\stackrel{~}{W}(q|\epsilon ^{},\epsilon ;s)\stackrel{~}{F}(q|\epsilon ^{},\epsilon _1)\stackrel{~}{W}(q|\epsilon _1,\epsilon ;s)`$ (C5)
$`𝒩\left(\epsilon _1\right)\stackrel{~}{W}(0|\epsilon ^{},\epsilon _1;s)\stackrel{~}{F}(q|\epsilon _1,\epsilon )\stackrel{~}{W}(q|\epsilon ,\epsilon _1;s)+𝒩\left(\epsilon _1\right)\stackrel{~}{W}(q|\epsilon ^{},\epsilon _1;s)\stackrel{~}{F}(q|\epsilon ^{},\epsilon )\stackrel{~}{W}(q|\epsilon ,\epsilon _1;s)\}.`$
In order to eliminate $`\stackrel{~}{F}`$, we have to replace $`\stackrel{~}{F}(q|\epsilon ^{},\epsilon )=F(q|\epsilon ^{},\epsilon )f(\epsilon ,s)\delta (\epsilon ^{}\epsilon )`$, where $`F`$ is the effective medium approximation of the diffusion propagator. Let us first work out the local contribution. Replacing $`\stackrel{~}{F}`$ by $`f\delta `$, the first and the third term of Eq.(C5) cancel each other. The 4th term is zero, taken into account that the delta function of the effective medium is multiplied by its argument. Thus, when $`\stackrel{~}{F}`$ is replaced by $`f(\epsilon ,s)\delta (\epsilon ^{}\epsilon )`$, only the 2nd term of Eq.(C5) survives. To simplify this term we take into account that
$$\stackrel{~}{W}(q|ϵ;s)=\stackrel{~}{W}(0|ϵ,s)\varphi (\frac{q\rho _c(ϵ,s)}{2\alpha }),$$
(C6)
where $`\varphi `$ is a dimensionless function. Consequently, we obtain
$``$ $`{\displaystyle 𝑑𝐪𝑑\epsilon ^{}\left(\epsilon ^{}\epsilon \right)𝒩\left(\epsilon \right)f(\epsilon ^{},s)\left[\stackrel{~}{W}\left(q|\epsilon ^{},\epsilon ;s\right)\right]^2}{\displaystyle 𝑑𝐪\left[\stackrel{~}{W}\left(𝐪|\epsilon ,s\right)\right]^2\frac{\omega ^2}{2}f(\epsilon ,s)𝒩\left(\epsilon \right)}`$ (C7)
$`=`$ $`{\displaystyle \frac{1}{2}}\omega ^2\left({\displaystyle \frac{2\alpha }{\rho _c(ϵ,s)}}\right)^d𝒩(ϵ)f(ϵ,s)\stackrel{~}{W}^2(0|ϵ,s){\displaystyle d^dx\varphi ^2(x)}.`$ (C8)
Here terms proportional to $`\omega f^1(\epsilon ,s)df(\epsilon ,s)/d\epsilon 1`$ have been neglected.
Let us now focus on the contribution of the regular part of the diffusion propagator to the self-consistency equation. Owing to the step functions in $`F`$ an $`\stackrel{~}{W}`$, the first and third term in Eq.(C5) are zero, when $`\stackrel{~}{F}`$ is replaced by $`F`$. Thus we are left with:
$$d𝐪d\epsilon d\epsilon _1(\epsilon \epsilon )\{𝒩\left(\epsilon \right)\stackrel{~}{W}(q|\epsilon ,\epsilon ;s)F(q|\epsilon ,\epsilon _1)\stackrel{~}{W}(q|\epsilon _1,\epsilon ;s)+𝒩\left(\epsilon _1\right)\stackrel{~}{W}(q|\epsilon ^{},\epsilon _1;s)F(q|\epsilon ^{},\epsilon )\stackrel{~}{W}(q|\epsilon ,\epsilon _1;s)\}.$$
(C9)
The range of integration in (C9) is determined by the step functions in $`F`$ and $`\stackrel{~}{W}(q|\epsilon ^{},\epsilon ;s)=\theta (\epsilon ^{}\epsilon )\theta (\omega \epsilon ^{}+\epsilon )\stackrel{~}{W}(q|\epsilon ^{};s)`$. Owing to these step functions, the energy integrations extend at most over intervals of length $`\omega `$, so that the quasi-elastic approximation can again be applied to the effective-transition probabilities and the density of states. The diffusion propagator enetering Eq.(C9), however, can not be dealt with in this way since for arbitrary $`q`$ the derivatives of the diffusion propagator are not small as compared to the diffusion propagator itself. Therefore, another procedure is needed. To simplify this expression further we consider Eq.(31). If terms small with respect to $`\omega N^{}/N`$ and $`\omega \stackrel{~}{W}^{}(q|ϵ,s))/\stackrel{~}{W}(q|ϵ,s)`$ are neglected the function
$$\mathrm{\Phi }(x|y^{},y)=\omega ^2F(x\frac{2\alpha }{\rho _c(\omega y,s)}|\omega y^{},\omega y)\stackrel{~}{W}(\omega y,s)𝒩(\omega y),$$
(C10)
can be introduced, that satisfies the equation
$`{\displaystyle \frac{s}{\stackrel{~}{W}(\omega y,s)𝒩(\omega y)\omega }}\mathrm{\Phi }(x|y^{},y)=\delta (y^{}y)+{\displaystyle \underset{0}{\overset{1}{}}}𝑑y_1[\mathrm{\Phi }(x|y^{},y_1+y)\varphi (x)\mathrm{\Phi }(x|y^{},y)\varphi (0)].`$ (C11)
Then the expression (C9) can be cast into the form
$$\omega \stackrel{~}{W}(\omega y)(\frac{2\alpha }{\rho _c(\omega y,s)})^dd^dx\varphi ^2(x)\left[\underset{0}{\overset{1}{}}𝑑y^{}\underset{0}{\overset{y^{}}{}}𝑑y_1y^{}\mathrm{\Phi }(x|y^{}+\frac{ϵ}{\omega },y_1+\frac{ϵ}{\omega })+\underset{0}{\overset{1}{}}\underset{0}{\overset{1y_1}{}}𝑑y^{}y^{}\mathrm{\Phi }(x|y^{}+\frac{ϵ}{\omega },\frac{ϵ}{\omega })\right]$$
(C12)
For $`s=0`$ Eq.(C11) does not contain any physical parameter. It only leads to the determination of the function
$`\mathrm{\Phi }_0(x|y^{},y)=\theta (y^{}y)\mathrm{\Phi }_0(x|y^{}y),`$
which satisfies the equation
$`0=\delta (y^{}y)+{\displaystyle \underset{0}{\overset{1}{}}}𝑑y_1[\mathrm{\Phi }_0(x|y^{},y_1+y)\varphi (x)\mathrm{\Phi }_0(x|y^{},y)\varphi (0)].`$ (C13)
Provided, we restrict our consideration to small frequencies in deriving the self-consistency equation we only need to take into account the linear contribution of the function $`\mathrm{\Phi }`$ with respect to $`s/(\stackrel{~}{W}(\omega y)𝒩(\omega y)\omega )`$. Then, using again the smallness of the variation of the effective transition probabilities and the density of states with respect to changes of energy over intervals of length $`\omega `$, $`\mathrm{\Phi }`$ can be approximated as
$$\mathrm{\Phi }(x|y^{},y)=\mathrm{\Phi }_0(x|y^{}y)\frac{s}{\stackrel{~}{W}(\omega y)𝒩(\omega y)\omega }\underset{y}{\overset{y^{}}{}}𝑑y_1\mathrm{\Phi }_0(x|yy_1)\mathrm{\Phi }_0(x|y_1y).$$
(C14)
Using this expression the self-consistency equation takes the form (62), where the coefficients $`a`$ and $`b`$ are given by
$$a=\frac{1}{2}\frac{d^dx\varphi ^2(x)\underset{0}{\overset{1}{}}(1y)^2\mathrm{\Phi }_0(x|y)}{d^dx\varphi ^2(x)},$$
(C15)
$$b=\frac{1}{2}\frac{d^dx\varphi ^2(x)_0^1𝑑y(1y)^2_0^y𝑑y_1\mathrm{\Phi }_0(x|yy_1)\mathrm{\Phi }_0(x|y_1)}{d^dx\varphi ^2(x)}.$$
(C16)
## D Calculation of the energy relaxation speed for constant DOS
To calculate the integral (77), we again use Eq.(71) to change the integration variable from $`s`$ to $`y=v/v_0`$. Furthermore, for convinience, we introduce the parameter $`\tau =\mathrm{\Omega }t`$. Doing so, we obtain:
$$\frac{dE\left(t\right)}{dt}=v_0F\left(\tau \right),$$
(D1)
where $`F(\tau )`$ is given by:
$$F\left(\tau \right)=_C\frac{dy}{2\pi i}\frac{1+\mathrm{ln}y}{\mathrm{ln}y}\mathrm{exp}\left(\tau y\mathrm{ln}y\right).$$
(D2)
After an integration by parts we obtain:
$$\frac{dF}{d\tau }=\frac{1}{\tau }_C\frac{dy}{2\pi i}\mathrm{exp}\left(\tau y\mathrm{ln}y\right)\frac{1}{\tau }S\left(\tau \right).$$
(D3)
To calculate $`F(\tau )`$, Eq. (D3) should be supplemented with the condition $`F(+\mathrm{})=1`$.
To evaluate $`S\left(\tau \right)`$, it is convenient to use the integration contour $`Im(y\mathrm{ln}y)=0`$, or, introducing polar coordinates $`y=r\mathrm{exp}\left(i\theta \right)`$:
$$r=\mathrm{exp}\left(\theta \mathrm{cot}\theta \right).$$
(D4)
Substituting Eq.(D4) into Eq.(D3), we have:
$$S\left(\tau \right)=_\pi ^\pi \frac{d\theta }{2\pi }V\left(\theta \right)e^{\tau V\left(\theta \right)},V\left(\theta \right)=\frac{\theta }{\mathrm{sin}\theta }e^{\theta \mathrm{cot}\theta }.$$
(D5)
The value of the integral may be estimated by the saddle-point method, looking for the maxima of the expression $`\mathrm{ln}V\left(\theta \right)\tau V\left(\theta \right)`$. The stationary point equation is:
$$V^{}\left(\theta \right)\left[\tau \frac{1}{V\left(\theta \right)}\right]=0.$$
(D6)
As $`\tau 1`$, the stationary point is $`\theta _s=0`$, and the asymptote of Eq.(D5) is:
$$S\left(\tau \right)\left(2\pi e\tau \right)^{1/2}\mathrm{exp}\left(\tau /e\right).$$
(D7)
On the other hand, if $`\tau 1`$, we have two stationary points $`\pm \theta _s`$, $`V\left(\pm \theta _s\right)=1/\tau `$, $`\theta _s=\pi \delta `$, and for $`\delta 1`$ the following equation may be obtained:
$`{\displaystyle \frac{\pi }{\delta }}\mathrm{exp}\left({\displaystyle \frac{\pi }{\delta }}1\right)={\displaystyle \frac{1}{\tau }},`$
the solution for which at small $`\tau `$ is: $`\pi /\delta =𝒲(e/\tau )\mathrm{ln}\left(e/\tau \right)\mathrm{ln}\mathrm{ln}(e/\tau ).`$The asymptotics of $`S\left(\tau \right)`$ turns out to be:
$$S\left(\tau \right)\frac{\sqrt{2\pi }}{e\tau }\left[\mathrm{ln}\left(e/\tau \right)\right]^2.$$
(D8)
## E The saddle-point approximation
For energy dependent density of states it turns out to be difficult to obtain explicite expressions for the time dependence of the energy distribution function. Here, according to Eq.(44), the time-dependence has to be calculated from the equation:
$$P(\epsilon ,\epsilon _0,t)=\underset{i\mathrm{}}{\overset{+i\mathrm{}}{}}\frac{ds}{2\pi iv(\epsilon ,s)}\mathrm{exp}\left[sts\underset{\epsilon }{\overset{\epsilon _0}{}}\frac{d\epsilon ^{}}{v(\epsilon ^{},s)}\right],$$
(E1)
where $`v(\epsilon ,s)`$ is given by Eq. (71) or (73). Let us try to integrate over $`s`$, using the saddle-point approximation. The saddle-point position $`s_0(\epsilon ,\epsilon _0,t)`$ may be obtained from the equation:
$$t=\underset{\epsilon }{\overset{\epsilon _0}{}}\frac{d\epsilon ^{}}{v\left(\epsilon ^{}\right)g(\epsilon ^{},s_0)}s_0\underset{\epsilon }{\overset{\epsilon _0}{}}\frac{d\epsilon ^{}}{v\left(\epsilon ^{}\right)g^2(\epsilon ^{},s_0)}\frac{g(\epsilon ,s_0)}{s_0},$$
(E2)
where $`g=v(\epsilon ,s)/v(\epsilon ,0)=\mathrm{exp}𝒲(s/\mathrm{\Omega }(\epsilon ))`$ was introduced. The expression for the diffusion propagator in the saddle-point approximation may be written as:
$$P_L(\epsilon ,\epsilon _0,t)=\frac{1}{\sqrt{4\pi D(\epsilon ,\epsilon _0,s_0)}}\frac{1}{v_0\left(\epsilon \right)g(\epsilon ,s_0)}\mathrm{exp}\left[s_0^2\underset{\epsilon }{\overset{\epsilon _0}{}}\frac{d\epsilon ^{}}{v_0\left(\epsilon ^{}\right)g^2(\epsilon ^{},s_0)}\frac{g(\epsilon ,s_0)}{s_0}\right],$$
(E3)
where:
$$D(\epsilon ,\epsilon _0,s_0)=\frac{^2}{s_0^2}\underset{\epsilon }{\overset{\epsilon _0}{}}\frac{d\epsilon ^{}\mathrm{\Omega }\left(\epsilon ^{}\right)\mathrm{ln}g(\epsilon ^{},s_0)}{2v_0\left(\epsilon ^{}\right)}>0.$$
(E4)
Assuming $`s_0/\mathrm{\Omega }1`$, we have: $`g(\epsilon ,s_0)1+s_0/\mathrm{\Omega }s_0^2/2\mathrm{\Omega }^2`$,
$$s_0=\frac{tT(\epsilon ,\epsilon _0)}{2D(\epsilon ,\epsilon _0)},D(\epsilon ,\epsilon _0)=\underset{\epsilon }{\overset{\epsilon _0}{}}\frac{d\epsilon ^{}}{v_0\left(\epsilon ^{}\right)\mathrm{\Omega }\left(\epsilon ^{}\right)},T(\epsilon ,\epsilon _0)=\underset{\epsilon }{\overset{\epsilon _0}{}}\frac{d\epsilon ^{}}{v_0\left(\epsilon ^{}\right)}$$
(E5)
and, finally, the diffusion propagator:
$$P_L(\epsilon ,\epsilon _0,t)\frac{1}{\sqrt{4\pi D(\epsilon ,\epsilon _0)}}\frac{1}{v\left(\epsilon \right)}\mathrm{exp}\left[\frac{\left(tT(\epsilon ,\epsilon _0)\right)^2}{4D(\epsilon ,\epsilon _0)}\right].$$
(E6)
At a given time $`t`$ the exponent in the above distribution function is maximal at $`\epsilon =\epsilon _m`$, where $`\epsilon _m`$ is given by the condition $`t=T(\epsilon _m,\epsilon _0)`$, or:
$$t=\underset{\epsilon _m}{\overset{\epsilon _0}{}}\frac{d\epsilon ^{}}{v_0\left(\epsilon ^{}\right)}.$$
(E7)
Expanding the expression (E6) around $`\epsilon =\epsilon _m`$, we have gaussian distribution:
$$P_L(\epsilon ,\epsilon _0,t)\frac{1}{\sigma (\epsilon _m,\epsilon _0)\sqrt{2\pi }}\mathrm{exp}\left[\frac{\left(\epsilon \epsilon _m(\epsilon _0,t)\right)^2}{2\sigma ^2(\epsilon _m,\epsilon _0)}\right],$$
(E8)
with the dispersion:
$$\sigma ^2(\epsilon _m,\epsilon _0)=2v_0^2\left(\epsilon _m\right)D(\epsilon _m,\epsilon _0)=2v^2\left(\epsilon _m\right)\underset{\epsilon _m}{\overset{\epsilon _0}{}}\frac{d\epsilon ^{}}{v_0\left(\epsilon ^{}\right)\mathrm{\Omega }\left(\epsilon ^{}\right)}.$$
(E9)
Note, that according to Eq.(E7), the dispersion is time dependent. If the density of states $`𝒩(\epsilon )`$, and, consequently, $`\mathrm{\Omega }(\epsilon )`$ and $`v_0(\epsilon )`$, are strongly varying functions of $`\epsilon `$, the integrals can be further simplified. In this case we have:
$$\underset{\epsilon _m}{\overset{\epsilon _0}{}}\frac{d\epsilon ^{}}{v\left(\epsilon ^{}\right)}\left[\frac{dv\left(\epsilon _0\right)}{d\epsilon _0}\right]^1\left[\frac{dv\left(\epsilon _m\right)}{d\epsilon _m}\right]^1,$$
(E10)
$$\underset{\epsilon _m}{\overset{\epsilon _0}{}}\frac{d\epsilon ^{}}{v\left(\epsilon ^{}\right)\mathrm{\Omega }\left(\epsilon ^{}\right)}\left[\frac{d\left[v\left(\epsilon _0\right)\mathrm{\Omega }\left(\epsilon _0\right)\right]}{d\epsilon _0}\right]^1\left[\frac{d\left[v\left(\epsilon _m\right)\mathrm{\Omega }\left(\epsilon _m\right)\right]}{d\epsilon _m}\right]^1,$$
(E11)
$$_{\epsilon _m}^{\epsilon _0}\frac{d\epsilon ^{}}{v\left(\epsilon ^{}\right)\mathrm{\Omega }\left(\epsilon ^{}\right)}\left[\frac{dv\left(\epsilon _0\right)\mathrm{\Omega }\left(\epsilon _0\right)}{d\epsilon _0}\right]^1\left[\frac{dv\left(\epsilon _m\right)\mathrm{\Omega }\left(\epsilon _m\right)}{d\epsilon _m}\right]^1.$$
(E12)
The applicability condition for the saddle-point method is:
$$D^3(\epsilon ,\epsilon _0,s_0)\left[\frac{^3}{s_0^3}\underset{\epsilon }{\overset{\epsilon _0}{}}𝑑\epsilon ^{}\frac{\mathrm{\Omega }\left(\epsilon ^{}\right)\mathrm{ln}g(\epsilon ^{},s_0)}{2v\left(\epsilon ^{}\right)}\right]^2.$$
(E13)
Note, that this condition restricts the applicablity of the saddle-point approximation to times, that are not too large. Also, from Eqs.(E5) one can see, that this approximation is invalid when $`\epsilon `$ is close enough to the initial point $`\epsilon _0`$, — here $`\left|s_0\right|`$ becomes large, and it is not possible to expand over $`s_0`$. So, the initial stage of the evolution, when the entire distribution is concentrated near $`\epsilon _0`$, has to be investigated separately too.
## F Calculation of the momenta for large times
Here we present details on the calculation of the momenta of the distribution function for large times and for the exponential density of states. In our calculation we take into account that the density of states decays with decreasing energy.
Let us change the integration variable in (93), $`\rho z=𝒲\left(s\rho e^\rho \right)`$, $`\rho =𝒲\left(s^1ze^z\right)`$. As $`\rho 1`$, we have, introducing $`\rho _s=𝒲\left(s^1\right)`$, $`s=\rho _s^1\mathrm{exp}\left(\rho _s\right)`$:
$`\rho =𝒲\left(\rho _sze^{\rho s+z}\right)\rho _s+z+\mathrm{ln}z,`$
as $`\left|z\right|\left|\rho _s\right|`$ and $`\left|\mathrm{ln}z\right|\left|\rho _s\right|`$. Both inequalities holds within actual region of integration. We can write now:
$$\chi _{nL}(s)=\frac{1}{s}\underset{0}{\overset{\mathrm{}}{}}𝑑z\left[\mathrm{ln}\rho _s\frac{z+\mathrm{ln}z}{\rho _s}\right]^n\frac{}{z}\mathrm{exp}\left[\frac{1}{\overline{\omega }\rho _s^2}z\left(1+\frac{1}{2}z\right)\right].$$
(F1)
Something more about the approximation we have chosen: every extra power of $`z/\rho _s`$ means an extra power of $`\sqrt{\overline{\omega }}`$, which can be neglected. An extra power of $`z`$ means the multiplier $`\sqrt{\overline{\omega }}\rho _s`$. All our considerations are valid if $`\overline{\omega }\rho _s1`$ only. Therefore, we have the reason to omit extra powers of $`z/\rho _s\sqrt{\overline{\omega }}`$ in the following, but to keep some extra powers of $`z^11/\sqrt{\overline{\omega }}\rho _s\sqrt{\overline{\omega }}`$.
Thus we have, for example:
$`\chi _1(s)={\displaystyle \frac{1}{s}}\mathrm{ln}\rho _s+{\displaystyle \frac{1}{s\rho _s}}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑z\left(z+\mathrm{ln}z\right){\displaystyle \frac{d}{dz}}\mathrm{exp}\left[{\displaystyle \frac{1}{\overline{\omega }\rho _s^2}}z\left(1+{\displaystyle \frac{1}{2}}z\right)\right]`$ (F2)
$`=`$ $`{\displaystyle \frac{1}{s}}\mathrm{ln}\rho _s+{\displaystyle \frac{\sqrt{\overline{\omega }}}{s}}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑y\left(y+{\displaystyle \frac{\mathrm{ln}\left(\sqrt{\overline{\omega }}\rho _s\right)}{\sqrt{\overline{\omega }}\rho _s}}+{\displaystyle \frac{\mathrm{ln}y}{\sqrt{\overline{\omega }}\rho _s}}\right){\displaystyle \frac{d}{dy}}\mathrm{exp}\left({\displaystyle \frac{y^2}{2}}{\displaystyle \frac{y}{\sqrt{\overline{\omega }}\rho _s}}\right),`$ (F3)
where
$`y={\displaystyle \frac{z}{\sqrt{\overline{\omega }}\rho _s}}.`$
Within three leading orders on the parameter $`1/\sqrt{\overline{\omega }}\rho _s`$ we have:
$$\chi _{1L}(s)\frac{\mathrm{ln}\rho _s}{s}\sqrt{\frac{\pi \overline{\omega }}{2}}\frac{1}{s}\frac{\mathrm{ln}\left(\sqrt{\overline{\omega }}\rho _s\right)}{s\rho _s}+\left(1\frac{\mathrm{ln}2\gamma }{2}\right)\frac{1}{s\rho _s}+\sqrt{\frac{\pi }{8\overline{\omega }}}\frac{1}{s\rho _s^2},$$
(F4)
which can be written, using $`\rho _s+\mathrm{ln}\rho _s=\mathrm{ln}\left(1/s\right)`$ also as:
$$\chi _{1L}(s)=\frac{1}{s}\mathrm{ln}\mathrm{ln}\frac{b\sqrt{\overline{\omega }}}{s}\sqrt{\frac{\pi \overline{\omega }}{2}}\frac{1}{s}+\sqrt{\frac{\pi }{8\overline{\omega }}}\frac{1}{s\rho _s^2},b=\sqrt{2}e^{1\gamma /2},$$
(F5)
where $`\gamma `$ is the Euler’s constant. In the time representation we have:
$`\chi _1(t)={\displaystyle _C}{\displaystyle \frac{ds}{2\pi i}}\chi _{1L}(s)e^{st}={\displaystyle _C}{\displaystyle \frac{ds}{2\pi is}}\mathrm{ln}\mathrm{ln}{\displaystyle \frac{b\sqrt{\overline{\omega }}}{s}}e^{st}\sqrt{{\displaystyle \frac{\pi \overline{\omega }}{2}}}+\sqrt{{\displaystyle \frac{\pi }{8\overline{\omega }}}}{\displaystyle _C}{\displaystyle \frac{ds}{2\pi is\mathrm{ln}^2\frac{\sqrt{\overline{\omega }}}{s}}}e^{st}.`$
At large $`t`$, $`b\sqrt{\overline{\omega }}t1`$, the first integral, after the substitution: $`st=y`$, may be expanded as:
$`{\displaystyle _C}{\displaystyle \frac{ds}{2\pi is}}\mathrm{ln}\mathrm{ln}{\displaystyle \frac{b\sqrt{\overline{\omega }}}{s}}e^{st}={\displaystyle _C}{\displaystyle \frac{dy}{2\pi iy}}\mathrm{ln}\left[\mathrm{ln}b\sqrt{\overline{\omega }}t\mathrm{ln}y\right]e^y=\mathrm{ln}\mathrm{ln}b\sqrt{\overline{\omega }}t`$
$`{\displaystyle \frac{1}{\mathrm{ln}b\sqrt{\overline{\omega }}t}}{\displaystyle _C}{\displaystyle \frac{dy}{2\pi iy}}e^y\mathrm{ln}y{\displaystyle \frac{1}{2\mathrm{ln}^2b\sqrt{\overline{\omega }}t}}{\displaystyle _C}{\displaystyle \frac{dy}{2\pi iy}}e^y\mathrm{ln}^2y\mathrm{}`$
$`=`$ $`\mathrm{ln}\mathrm{ln}b\sqrt{\overline{\omega }}t+{\displaystyle \frac{\gamma }{\mathrm{ln}b\sqrt{\overline{\omega }}t}}+{\displaystyle \frac{\pi ^2/6\gamma ^2}{2\mathrm{ln}^2b\sqrt{\overline{\omega }}t}}.`$
As for the second integral, one may restrict oneself with the leading term only. As a result, we have:
$$\chi _1(t)=\mathrm{ln}\mathrm{ln}\stackrel{~}{b}\sqrt{\overline{\omega }}t\sqrt{\frac{\pi \overline{\omega }}{2}}+\sqrt{\frac{\pi }{8\overline{\omega }}}\frac{1}{\mathrm{ln}^2\stackrel{~}{b}\sqrt{\overline{\omega }}t}+\mathrm{},\stackrel{~}{b}=\sqrt{2}e^{1+\gamma /2}.$$
(F6)
Let us consider now the second momentum:
$$\chi _{2L}(s)=\underset{0}{\overset{\mathrm{}}{}}\frac{d\rho }{\rho }\mathrm{ln}^2\rho \varphi _L(\rho ;s)\mathrm{ln}^2\rho .$$
(F7)
Again, we change the integration variable to $`z=𝒲\left(s\rho e^\rho \right)`$, write $`\rho 1`$ as $`\rho =\rho _s+z+\mathrm{ln}z`$, and set the lower limit of the integration to be zero. Restricting ourselves with the lowest order in $`z/\rho _s\sqrt{\overline{\omega }}`$, we set $`\mathrm{ln}\rho =\mathrm{ln}\rho _s+(z+\mathrm{ln}z)/\rho _s`$, and:
$`{\displaystyle \frac{d\rho }{\rho }}\varphi _L={\displaystyle \frac{dz}{s}}{\displaystyle \frac{}{z}}\mathrm{exp}\left[{\displaystyle \frac{1}{\overline{\omega }\rho _s^2}}z\left(1+{\displaystyle \frac{1}{2}}z\right)\right].`$
Taking into account:
$`\chi _{1L}(s)={\displaystyle \frac{1}{s}}\mathrm{ln}\rho _s{\displaystyle \frac{z+\mathrm{ln}z}{\rho _s}},`$
we can write:
$$\chi _{2L}(s)=s\chi _{1L}^2(s)+\left(\frac{z+\mathrm{ln}z}{\rho _s}\right)^2\frac{z+\mathrm{ln}z}{\rho _s}^2.$$
(F8)
Then we have:
$`\left({\displaystyle \frac{z+\mathrm{ln}z}{\rho _s}}\right)^2`$ $`=`$ $`{\displaystyle \frac{1}{s\rho _s^2}}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑z(z+\mathrm{ln}z)^2{\displaystyle \frac{d}{dz}}\mathrm{exp}\left[{\displaystyle \frac{1}{\overline{\omega }\rho _s^2}}z\left(1+{\displaystyle \frac{1}{2}}z\right)\right]`$
$`=`$ $`{\displaystyle \frac{\overline{\omega }}{s}}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑y\left(y+{\displaystyle \frac{\mathrm{ln}\left(\sqrt{\overline{\omega }}\rho _s\right)}{\sqrt{\overline{\omega }}\rho _s}}+{\displaystyle \frac{\mathrm{ln}y}{\sqrt{\overline{\omega }}\rho _s}}\right)^2\mathrm{exp}\left({\displaystyle \frac{y^2}{2}}{\displaystyle \frac{y}{\sqrt{\overline{\omega }}\rho _s}}\right).`$
Expanding this expressions in power series on $`1/\sqrt{\overline{\omega }}\rho _s`$, and keeping three leading powers, it may be easily obtained:
$`\left({\displaystyle \frac{z+\mathrm{ln}z}{\rho _s}}\right)^2={\displaystyle \frac{2\overline{\omega }}{s}}+{\displaystyle \frac{\sqrt{2\pi \overline{\omega }}\mathrm{ln}\left(\sqrt{\overline{\omega }}\rho _s\right)}{s\rho _s}}{\displaystyle \frac{\mathrm{ln}2+\gamma }{2}}{\displaystyle \frac{\sqrt{2\pi \overline{\omega }}}{s\rho _s}}+{\displaystyle \frac{\mathrm{ln}^2\left(\sqrt{\overline{\omega }}\rho _s\right)}{s\rho _s^2}}`$
$`2\left(1{\displaystyle \frac{\mathrm{ln}2\gamma }{2}}\right){\displaystyle \frac{\mathrm{ln}\left(\sqrt{\overline{\omega }}\rho _s\right)}{s\rho _s^2}}+\left[{\displaystyle \frac{\pi ^2}{24}}\mathrm{ln}2+\gamma +\left({\displaystyle \frac{\mathrm{ln}2\gamma }{2}}\right)^2\right]{\displaystyle \frac{1}{s\rho _s^2}}.`$
¿From Eq. (F4) we have:
$`{\displaystyle \frac{z+\mathrm{ln}z}{\rho _s}}=\sqrt{{\displaystyle \frac{\pi \overline{\omega }}{2}}}{\displaystyle \frac{1}{s}}+{\displaystyle \frac{\mathrm{ln}\left(\sqrt{\overline{\omega }}\rho _s\right)}{s\rho _s}}\left(1{\displaystyle \frac{\mathrm{ln}2\gamma }{2}}\right){\displaystyle \frac{1}{s\rho _s}}\sqrt{{\displaystyle \frac{\pi }{8\overline{\omega }}}}{\displaystyle \frac{1}{s\rho _s^2}}.`$
So, from Eq. (F8) we obtain:
$$\chi _{2L}(s)=\chi _{1L}^2(s)+\left(2\frac{\pi }{2}\right)+\left(1\mathrm{ln}2\right)\frac{\sqrt{2\pi \overline{\omega }}}{s\rho _s}+\left(\frac{\pi ^2}{24}+\frac{\pi }{2}1\right)\frac{1}{s\rho _s^2}.$$
(F9)
However, we want to calculate central momentum:
$$\mu _2(t)=\chi _2(t)\chi _1^2(t),$$
(F10)
or, in the Laplace representation:
$`\mu _{2L}(s)`$ $`=`$ $`\chi _{2L}(s)\left(\chi _1^2\right)_L(s),`$ (F11)
$`\left(\chi _1^2\right)_L(s)`$ $``$ $`{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑te^{st}\chi _1^2(t)={\displaystyle \underset{i\mathrm{}+0}{\overset{+i\mathrm{}+0}{}}}{\displaystyle \frac{ds_1}{2\pi i}}\chi _{1L}(s_1)\chi _{1L}(ss_1).`$ (F12)
Taking the expression for $`\chi _{1L}`$ from Eq. (F5), we have:
$`\left(\chi _1^2\right)_L(s)={\displaystyle \underset{i\mathrm{}+0}{\overset{+i\mathrm{}+0}{}}}{\displaystyle \frac{ds_1}{2\pi i}}{\displaystyle \frac{\mathrm{ln}\mathrm{ln}\frac{b\sqrt{\overline{\omega }}}{s}\mathrm{ln}\mathrm{ln}\frac{b\sqrt{\overline{\omega }}}{ss_1}}{s_1\left(ss_1\right)}}+{\displaystyle \frac{\sqrt{2\pi \overline{\omega }}}{s}}\mathrm{ln}\mathrm{ln}{\displaystyle \frac{b\sqrt{\overline{\omega }}}{s}}+{\displaystyle \frac{\pi \overline{\omega }}{2s}}`$ (F13)
$`{\displaystyle \underset{i\mathrm{}+0}{\overset{+i\mathrm{}+0}{}}}{\displaystyle \frac{ds_1}{2\pi i}}{\displaystyle \frac{\mathrm{ln}\mathrm{ln}\frac{b\sqrt{\overline{\omega }}}{s_1}}{s_1\left(ss_1\right)\mathrm{ln}^2\frac{b\sqrt{\overline{\omega }}}{ss_1}}}{\displaystyle \frac{\pi }{2s\mathrm{ln}^2\frac{b\sqrt{\overline{\omega }}}{s}}}+\mathrm{}.`$ (F14)
Let us consider the first integral. After changing the integration variable, $`s_1=sx`$, we have:
$`{\displaystyle \frac{1}{s}}\mathrm{ln}^2\mathrm{ln}{\displaystyle \frac{b\sqrt{\overline{\omega }}}{s}}+{\displaystyle \frac{2}{s}}\mathrm{ln}{\displaystyle \frac{b\sqrt{\overline{\omega }}}{s}}{\displaystyle _C}{\displaystyle \frac{dx}{2\pi i}}{\displaystyle \frac{\mathrm{ln}\left(1\mathrm{ln}(1x)/\mathrm{ln}\left(\frac{b\sqrt{\overline{\omega }}}{s}\right)\right)}{x(1x)}}`$
$`+{\displaystyle _c}{\displaystyle \frac{dx}{2\pi i}}{\displaystyle \frac{\mathrm{ln}\left(1\mathrm{ln}x/\mathrm{ln}\left(\frac{b\sqrt{\overline{\omega }}}{s}\right)\right)\mathrm{ln}\left(1\mathrm{ln}(1x)/\mathrm{ln}\left(\frac{b\sqrt{\overline{\omega }}}{s}\right)\right)}{x(1x)}}.`$
Closing the integration contour to the left, one can easily prove, that the first intehral in the above expression is zero. The second one at small $`s`$ may be expanded into powers of $`1/\mathrm{ln}s`$. In the leading order we have:
$`{\displaystyle \frac{1}{s}}\mathrm{ln}^2\mathrm{ln}{\displaystyle \frac{b\sqrt{\overline{\omega }}}{s}}+{\displaystyle \frac{1}{s\mathrm{ln}^2\frac{b\sqrt{\overline{\omega }}}{s}}}{\displaystyle _C}{\displaystyle \frac{dx}{2\pi i}}{\displaystyle \frac{\mathrm{ln}x\mathrm{ln}(1x)}{x(1x)}}.`$
Performing integration by part in the integral, and then deforming the integration contour to make it run along the lower and upper shores of the cut $`(\mathrm{},0)`$ of the integration function:
$`{\displaystyle _C}{\displaystyle \frac{dx}{2\pi i}}{\displaystyle \frac{\mathrm{ln}x\mathrm{ln}(1x)}{x(1x)}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _C}{\displaystyle \frac{dx}{2\pi i}}\mathrm{ln}^2x{\displaystyle \frac{d}{dx}}{\displaystyle \frac{\mathrm{ln}(1x)}{1x}}={\displaystyle \frac{1}{2\pi }}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑x{\displaystyle \frac{\mathrm{ln}(1+x)1}{(1+x)^2}}Im\mathrm{ln}^2(x+i0)`$
$`=`$ $`{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑x{\displaystyle \frac{\left[\mathrm{ln}(1+x)1\right]\mathrm{ln}x}{(1+x)^2}}={\displaystyle \frac{\pi ^2}{6}}.`$
The evaluation of the second integral in Eq. (F14) gives:
$`2{\displaystyle \frac{\mathrm{ln}\mathrm{ln}\frac{b\sqrt{\overline{\omega }}}{s}}{s\mathrm{ln}^2\frac{b\sqrt{\overline{\omega }}}{s}}}+O\left({\displaystyle \frac{1}{s\mathrm{ln}^4\frac{b\sqrt{\overline{\omega }}}{s}}}\right)`$
As a result, one can write:
$$\left(\chi _1^2\right)_L(s)=s\chi _{1L}^2(s)+\frac{\pi ^2}{6s\rho _s^2}+\mathrm{}$$
(F15)
Thus, one have the second central momentum in the Laplace representation to be:
$$\mu _{2L}(s)=(2\frac{\pi }{2})\frac{\overline{\omega }}{s}+(1\mathrm{ln}2)\frac{\sqrt{2\pi \overline{\omega }}}{s\rho _s}+\left(\frac{\pi }{2}\frac{\pi ^2}{8}1\right)\frac{1}{s\rho _s^2},$$
(F16)
and the time dependence of the second central momentum:
$$\mu _2(t)=\left(2\frac{\pi }{2}\right)\overline{\omega }+\left(1+\mathrm{ln}2\right)\frac{\sqrt{2\pi \overline{\omega }}}{\mathrm{ln}t}.$$
(F17)
Expression for the third central momentum $`\mu _3\left(t\right)`$ (97) may be obtained quite analogously.
## G Asymptotic form of the distribution function at large times
One can write the large times distribution function (92), after the replacement of the integration variable $`sz=𝒲\left(s\rho e^\rho \right)`$as follows:
$$\varphi (\epsilon ,t)=\frac{}{t}\psi (\epsilon ,t)$$
(G1)
$$\psi (\epsilon ,t)=\rho e^\rho \underset{i\mathrm{}+\delta }{\overset{+i\mathrm{}+\delta }{}}\frac{dz}{2\pi iz^2}(1+z)\frac{}{\epsilon }\mathrm{exp}\left[z+t\frac{z}{\rho }e^{z\rho }\frac{1}{\overline{\omega }\rho ^2}z\left(1+\frac{1}{2}z\right)\right].$$
(G2)
One can try to calculate the latter integral within the saddle point approximation. The saddle point equation is:
$`f^{}(z)=1+t{\displaystyle \frac{z+1}{\rho }}e^{z\rho }{\displaystyle \frac{1}{\overline{\omega }\rho ^2}}(z+1)=0.`$
Assuming $`\left|z\right|\overline{\omega }\rho ^2`$ (remind that $`\omega \rho ^21`$), we can neglect the third term in the above equation, and to obtain:
$`z_c+1=𝒲\left({\displaystyle \frac{\rho }{t}}e^{\rho +1}\right).`$
As we shall see later, the actual values of parameters $`\rho `$, $`t`$, corresponding to the body of the distribution function, obeys the inequality $`\rho e^\rho /t1`$, and, therefore, $`z_c1`$. But this leads to the saddle point parameter, $`\left|f^{\prime \prime \prime }(z_c)/\left[f^{\prime \prime }(z_c)\right]^{3/2}\right|=\left|(z_c+3)/(z_c+2)^{3/2}\right|1`$, to be of the order of 1, instead of being much less. Let us note, however, that the function in the exponent in Eq. (G2) may be represented as a sum of two terms, of which the first one, $`z+tz/\rho \mathrm{exp}\left(z\rho \right)`$, ensures the function in the integral to have a maximum at $`z=z_c`$ of the width $`1`$, while the second one, $`\left(\overline{\omega }\rho ^2\right)^1z\left(1+z/2\right)`$, may be supposed to be almost constant within this peak, — it varies essentially at scales $`\mathrm{\Delta }z\sqrt{\overline{\omega }}\rho 1`$ within the actual interval of system’s parameters. Thus, one can replace that second term in the exponent by its value at some $`z_0=z_c+a`$, $`\left|a\right|1`$. After this, the integration may be easily performed, returning back to the initial integration variable $`s`$:
$`\psi (\epsilon ,t)\rho {\displaystyle \frac{}{\rho }}\mathrm{exp}\left\{{\displaystyle \frac{1}{\overline{\omega }\rho ^2}}𝒲\left(s_0\rho e^\rho \right)\left[1+{\displaystyle \frac{1}{2}}𝒲\left(s_0\rho e^\rho \right)\right]\right\}{\displaystyle \underset{i\mathrm{}+\delta }{\overset{+i\mathrm{}+\delta }{}}}{\displaystyle \frac{ds}{2\pi is^2}}e^{st},`$
where the last integral is equal to simply $`t`$, and, taking into account $`z_01`$, $`s_0=z_0/\rho \mathrm{exp}\left(z_0\rho \right)c/t`$, where $`c1`$ will be obtained later. Performing differentiation with respect to $`t`$, see Eq. (G2), and neglecting $`z_0/\overline{\omega }\rho ^21`$, we obtain finally:
$$\varphi (\epsilon ,t)\rho \frac{}{\rho }\mathrm{exp}\left\{\frac{1}{\overline{\omega }\rho ^2}𝒲\left(\frac{c}{t}\rho e^\rho \right)\left[1+\frac{1}{2}𝒲\left(\frac{c}{t}\rho e^\rho \right)\right]\right\}.$$
(G3)
To obtain $`c`$, let us calculate $`\chi _1(t)=\overline{\epsilon }(t)`$ with the distribution function (G3). We have:
$`\chi _1(t)={\displaystyle \frac{d\rho }{\rho }\varphi (\rho ,t)\mathrm{ln}\rho }={\displaystyle 𝑑z\mathrm{ln}\rho (z)\frac{}{z}\mathrm{exp}\left[\frac{1}{\overline{\omega }\rho ^2(z)}z\left(1+\frac{1}{2}z\right)\right]}`$
$`\rho (z)=𝒲\left({\displaystyle \frac{t}{c}}ze^z\right),`$
where the replacement of the integration variable $`\rho z=𝒲\left(c\rho e^\rho /t\right)`$ was done. Assuming $`\rho 1`$, one can write, using $`𝒲\left(x\right)x\mathrm{ln}x`$ as $`\left|x\right|1`$, that $`\rho (z)𝒲\left(t/c\right)+z+\mathrm{ln}z\rho _t+z+\mathrm{ln}z`$ (the condition $`z\sqrt{\overline{\omega }}\rho _t\rho _t=𝒲\left(t/c\right)`$ was used). In the leading order on $`\sqrt{\overline{\omega }}`$, one can replace $`\rho (z)`$ within the exponent with $`\rho _t`$, and $`\mathrm{ln}\rho (z)\mathrm{ln}\rho _t+(z+\mathrm{ln}z)/\rho _t`$. The limits of integration by $`z`$ are to be set $`(0,+\mathrm{})`$. Then we have:
$$\chi _1(t)=\mathrm{ln}\rho _t\underset{0}{\overset{\mathrm{}}{}}𝑑z\frac{z+\mathrm{ln}z}{\rho _t}\frac{}{z}\mathrm{exp}\left[\frac{1}{\overline{\omega }\rho _t^2}z\left(1+\frac{1}{2}z\right)\right].$$
(G4)
This latter integral, analogously to one in Eq. (F), may be evaluated as an expansion on small parameter $`1/\sqrt{\overline{\omega }}\rho `$:
$$\chi _1\left(t\right)=\mathrm{ln}\mathrm{ln}\left(\sqrt{2}e^{1\gamma /2}\frac{\sqrt{\overline{\omega }}t}{c}\right)\sqrt{\frac{\pi \overline{\omega }}{2}}+\mathrm{}$$
(G5)
Comparing Eqs.(G5) and (95), one can see, that $`c=\mathrm{exp}\left(\gamma \right)`$.
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# Macaulay 2 and the Geometry of Schemes
## Distinguished open sets
We begin with a simple example involving the definition of an affine scheme; see section I.1.4 in Eisenbud and Harris (1999). This example also indicates some of the subtleties involved in working with arithmetic schemes.
###### Question 1.
Let $`R=[x,y,z]`$ and $`X=\mathrm{Spec}(R)`$; in other words, $`X`$ is affine $`3`$-space over the integers. Let $`f=x`$ and consider the basic open subset $`X_fX`$.
1. If $`e_1=x+y+z`$, $`e_2=xy+xz+yz`$ and $`e_3=xyz`$ are the elementary symmetric functions then the set $`\left\{X_{e_i}\right\}_{1i3}`$ is an open cover of $`X_f`$.
2. If $`p_1=x+y+z`$, $`p_2=x^2+y^2+z^2`$ and $`p_3=x^3+y^3+z^3`$ are the power sum symmetric functions then $`\left\{X_{p_i}\right\}_{1i3}`$ is NOT an open cover of $`X_f`$.
###### Solution.
By Lemma I-6 in Eisenbud and Harris (1999), it suffices to show that $`e_1`$, $`e_2`$ and $`e_3`$ generate the unit ideal in $`R_f`$. This is equivalent to showing that $`x^m`$ belongs to the $`R`$-ideal $`e_1,e_2,e_3`$ for some $`m`$. In particular, the saturation $`\left(e_1,e_2,e_3:x^{\mathrm{}}\right)`$ is the unit ideal if and only if $`\left\{X_{e_i}\right\}_{1i3}`$ is an open cover of $`X_f`$. Macaulay2 allows us to work with homogenous ideals over $``$ and we obtain:
```
i1 : R = ZZ[x, y, z];
```
```
i2 : elementaryBasis = ideal(x+y+z, x*y+x*z+y*z, x*y*z);
```
```
o2 : Ideal of R
```
```
i3 : saturate(elementaryBasis, x)
```
```
o3 = ideal 1
```
```
o3 : Ideal of R
```
Similarly, to prove that $`\left\{X_{p_i}\right\}_{1i3}`$ is not an open cover of $`X_f`$, it is enough to show that $`\left(p_1,p_2,p_3:x^{\mathrm{}}\right)`$ is not the unit ideal. We compute this saturation:
```
i4 : powerSumBasis = ideal(x+y+z, x^2+y^2+z^2, x^3+y^3+z^3);
```
```
o4 : Ideal of R
```
```
i5 : saturate(powerSumBasis, x)
```
```
2 2
o5 = ideal (6, x + y + z, 2y - y*z + 2z , 3y*z)
```
```
o5 : Ideal of R
```
However, working over the field $``$, we find that $`\left(p_1,p_2,p_3:x^{\mathrm{}}\right)`$ is the unit ideal.
```
i6 : S = QQ[x, y, z];
```
```
i7 : powerSumBasis = ideal(x+y+z, x^2+y^2+z^2, x^3+y^3+z^3);
```
```
o7 : Ideal of S
```
```
i8 : saturate(powerSumBasis, x)
```
```
o8 = ideal 1
```
```
o8 : Ideal of S
```
## Irreducibility
The study of complex semisimple Lie algebras gives rise to an important family of algebraic varieties called nilpotent orbits. To illustrate one of the properties appearing in section I.2.1 of Eisenbud and Harris (1999), we examine the irreducibility of a particular nilpotent orbit.
###### Question 2.
Let $`X`$ be the set of complex $`3\times 3`$ matrices which are nilpotent. Show that $`X`$ is an irreducible algebraic variety.
###### Solution.
A $`3\times 3`$ matrix $`M`$ is nilpotent if and only if its minimal polynomial divides $`𝖳^k`$, for some $`k`$. Since each irreducible factor of the characteristic polynomial of $`M`$ is also a factor of the minimal polynomial, we conclude that the characteristic polynomial of $`M`$ is $`𝖳^3`$. It follows that the coefficients of the characteristic polynomial (except for the leading coefficient which is $`1`$) of a generic $`3\times 3`$ matrix define the algebraic variety $`X`$.
To show $`X`$ is irreducible over $``$, it is enough to construct a rational parameterization of $`X`$; see Proposition 4.5.6 in Cox, Little, and O’Shea (1996). To achieve this, observe that $`\mathrm{GL}_n()`$ acts on $`X`$ by conjugation. Jordan’s canonical form theorem implies that there are exactly three orbits; one for each of the following matrices:
$$N_0=\left[\begin{array}{ccc}0& 0& 0\\ 0& 0& 0\\ 0& 0& 0\end{array}\right],N_1=\left[\begin{array}{ccc}0& 1& 0\\ 0& 0& 0\\ 0& 0& 0\end{array}\right]\text{ and }N_2=\left[\begin{array}{ccc}0& 1& 0\\ 0& 0& 1\\ 0& 0& 0\end{array}\right].$$
Each orbit is defined by a rational parameterization, so it suffices to show that the closure of the orbit containing $`N_2`$ is the entire variety $`X`$. In Macaulay2, this calculation can be done as follows:
```
i1 : S = QQ[t, y_0..y_8, a..i, MonomialOrder => Eliminate 10];
```
```
i2 : N2 = (matrix {{0,1,0},{0,0,1},{0,0,0}}) ** S
```
```
o2 = {0} | 0 1 0 |
{0} | 0 0 1 |
{0} | 0 0 0 |
```
```
3 3
o2 : Matrix S <--- S
```
```
i3 : G = genericMatrix(S, y_0, 3, 3);
```
```
3 3
o3 : Matrix S <--- S
```
To determine the entries in $`GN_2G^1`$, we use the classical adjoint to construct the inverse of the generic matrix $`G`$.
```
i4 : adj = (G,i,j) -> (
n := degree target G;
m := degree source G;
(-1)^(i+j)*det(submatrix(G, {0..(i-1), (i+1)..(n-1)},
{0..(j-1), (j+1)..(m-1)}))
);
```
```
i5 : classicalAdjoint = (G) -> (
n := degree target G;
matrix table(n, n, (i, j) -> adj(G, j, i))
);
```
```
i6 : numerators = G*N2*classicalAdjoint(G);
```
```
3 3
o6 : Matrix S <--- S
```
```
i7 : D = det(G);
```
```
i8 : M = genericMatrix(S, a, 3, 3);
```
```
3 3
o8 : Matrix S <--- S
```
The entries in $`GN_2G^1`$ give a rational parameterization of the orbit generated by $`N_2`$. We give an “implicit representation” of this variety by using elimination theory; see section 3.3 in Cox, Little, and O’Shea (1996).
```
i9 : elimIdeal = minors(1, (D*id_(S^3))*M-numerators) + ideal(1 - D*t);
```
```
o9 : Ideal of S
```
```
i10 : closureOfOrbit = ideal(selectInSubring(1, gens gb elimIdeal));
```
```
o10 : Ideal of S
```
Finally, we check that the closure of this orbit is equal to $`X`$ scheme-theoretically.
```
i11 : X = ideal submatrix((coefficients({0},
det(M - t*id_(S^3))))_1, {1,2,3})
```
```
o11 = ideal (a + e + i, b*d - a*e + c*g + f*h - a*i - e*i,
- c*e*g + b*f*g + c*d*h - a*f*h - b*d*i + a*e*i)
```
```
o11 : Ideal of S
```
```
i12 : closureOfOrbit == X
```
```
o12 = true
```
## Singular Points
Section I.2.2 in Eisenbud and Harris (1999) provides the definition of a singular point of a scheme. In our third question, we study the singular locus of a family of elliptic curves. Section V.3 in Eisenbud and Harris (1999) also contains related material.
###### Question 3.
Consider a general form of degree $`3`$ in $`[x,y,z]`$;
$$F=ax^3+bx^2y+cx^2z+dxy^2+exyz+fxz^2+gy^3+hy^2z+iyz^2+jz^3.$$
Give necessary and sufficient conditions in terms of $`a,\mathrm{},j`$ for the cubic curve $`\mathrm{Proj}\left([x,y,z]/F\right)`$ to have a singular point.
###### Solution.
A time consuming elimination gives the degree $`12`$ polynomial which defines the singular locus of a general form of degree $`3`$. This can be done in Macaulay2 as follows. We have not displayed the output o6, as this discriminant has $`2040`$ terms in the $`10`$ variables $`a,\mathrm{},j`$.
```
i1 : S = QQ[x, y, z, a..j, MonomialOrder => Eliminate 2];
```
```
i2 : F = a*x^3+b*x^2*y+c*x^2*z+d*x*y^2+e*x*y*z+
f*x*z^2+g*y^3+h*y^2*z+i*y*z^2+j*z^3;
```
```
i3 : partials = submatrix(jacobian matrix{{F}}, {0..2}, {0})
```
```
o3 = {1} | 3x2a+2xyb+y2d+2xzc+yze+z2f |
{1} | x2b+2xyd+3y2g+xze+2yzh+z2i |
{1} | x2c+xye+y2h+2xzf+2yzi+3z2j |
```
```
3 1
o3 : Matrix S <--- S
```
```
i4 : singularities = ideal(partials) + ideal(F);
```
```
o4 : Ideal of S
```
```
i5 : elimDiscr = ideal selectInSubring(1, gens gb singularities);
```
```
o5 : Ideal of S
```
```
i6 : elimDiscr = substitute(elimDiscr, {z => 1});
```
```
o6 : Ideal of S
```
There is also a simple and more useful determinantal formula for this discriminant. It is a specialization of the formula (2.8) in section 3.2 in Cox, Little, and O’Shea (1998):
```
i7 : hessian = det submatrix(jacobian ideal partials, {0..2}, {0..2});
```
```
i8 : A = (coefficients({0,1,2}, submatrix(
jacobian matrix{{F}}, {0..2}, {0})))_1;
```
```
3 6
o8 : Matrix S <--- S
```
```
i9 : B = (coefficients({0,1,2}, submatrix(
jacobian matrix{{hessian}}, {0..2}, {0})))_1;
```
```
3 6
o9 : Matrix S <--- S
```
```
i10 : detDiscr = ideal det (A || B);
```
```
o10 : Ideal of S
```
```
i11 : detDiscr == elimDiscr
```
```
o11 = true
```
## Fields of Definition
Schemes over non-algebraically closed fields arise in number theory. Our solution to Exercise II-6 in Eisenbud and Harris (1999) indicates one technique for working over a number field in Macaulay2.
###### Question 4.
An inclusion of fields $`KL`$ induces a map $`𝔸_L^n𝔸_K^n`$. Find the images in $`𝔸_{}^2`$ of the following points of $`𝔸_\overline{}^2`$ under this map.
1. $`x\sqrt{2},y\sqrt{2}`$;
2. $`x\sqrt{2},y\sqrt{3}`$;
3. $`x\zeta ,y\zeta ^1`$ where $`\zeta `$ is a $`5`$-th root of unity ;
4. $`\sqrt{2}x\sqrt{3}y`$;
5. $`\sqrt{2}x\sqrt{3}y1`$.
###### Solution.
The images can be determined by (1) replacing coefficients not belonging to $`K`$ with indeterminates, (2) adding the minimal polynomials of these coefficients to the given ideal in $`𝔸_\overline{}^2`$ and (3) eliminating the new indeterminates. Here are the five examples:
```
i1 : S = QQ[a, b, x, y, MonomialOrder => Eliminate 2];
```
```
i2 : Ia = ideal(x-a, y-a, a^2-2);
```
```
o2 : Ideal of S
```
```
i3 : ideal selectInSubring(1, gens gb Ia)
```
```
2
o3 = ideal (x - y, y - 2)
```
```
o3 : Ideal of S
```
```
i4 : Ib = ideal(x-a, y-b, a^2-2, b^2-3);
```
```
o4 : Ideal of S
```
```
i5 : ideal selectInSubring(1, gens gb Ib)
```
```
2 2
o5 = ideal (y - 3, x - 2)
```
```
o5 : Ideal of S
```
```
i6 : Ic = ideal(x-a, y-a^4, a^4+a^3+a^2+a+1);
```
```
o6 : Ideal of S
```
```
i7 : ideal selectInSubring(1, gens gb Ic)
```
```
2 2 3 2
o7 = ideal (x*y - 1, x + y + x + y + 1, y + y + x + y + 1)
```
```
o7 : Ideal of S
```
```
i8 : Id = ideal(a*x+b*y, a^2-2, b^2-3);
```
```
o8 : Ideal of S
```
```
i9 : ideal selectInSubring(1, gens gb Id)
```
```
2 3 2
o9 = ideal(x - -*y )
2
```
```
o9 : Ideal of S
```
```
i10 : Ie = ideal(a*x+b*y-1, a^2-2, b^2-3);
```
```
o10 : Ideal of S
```
```
i11 : ideal selectInSubring(1, gens gb Ie)
```
```
4 2 2 9 4 2 3 2 1
o11 = ideal(x - 3x y + -*y - x - -*y + -)
4 2 4
```
```
o11 : Ideal of S
```
## Multiplicity
The multiplicity of a zero-dimensional scheme $`X`$ at a point $`pX`$ is defined to be the length of the local ring $`𝒪_{X,p}`$. Unfortunately, we cannot work directly in the local ring in Macaulay2. What we can do, however, is to compute the multiplicity by computing the degree of the component of $`X`$ supported at $`p`$; see page 66 in Eisenbud and Harris (1999).
###### Question 5.
What is the multiplicity of the origin $`(0,0,0)`$ as a zero of the polynomial equations
$$x^5+y^3+z^3=x^3+y^5+z^3=x^3+y^3+z^5=0\mathrm{?}$$
###### Solution.
If $`I`$ is the ideal generated by $`x^5+y^3+z^3`$, $`x^3+y^5+z^3`$ and $`x^3+y^3+z^5`$ in $`[x,y,z]`$, then the multiplicity of the origin is
$$dim_{}\frac{[x,y,z]_{x,y,z}}{I[x,y,z]_{x,y,z}}.$$
It follows that the multiplicity is the vector space dimension of the ring $`[x,y,z]/\phi ^1(I[x,y,z]_{x,y,z})`$ where $`\phi :[x,y,z][x,y,z]_{x,y,z}`$ is the natural map. Moreover, we can express this using ideal quotients:
$$\phi ^1(I[x,y,z]_{x,y,z})=(I:(I:x,y,z^{\mathrm{}})).$$
Carrying out this calculation in Macaulay2, we obtain:
```
i1 : S = QQ[x, y, z];
```
```
i2 : I = ideal(x^5+y^3+z^3, x^3+y^5+z^3, x^3+y^3+z^5);
```
```
o2 : Ideal of S
```
```
i3 : multiplicity = degree(I : saturate(I))
```
```
o3 = 27
```
## Flat Families
Non-reduced schemes arise naturally as the flat limit of a family of reduced schemes. Exercise III-68 in Eisenbud and Harris (1999) illustrates how a family of skew lines in $`^3`$ gives a double line with an embedded point.
###### Question 6.
Let $`L`$ and $`M`$ be the lines in $`_{k[t]}^3`$ given by $`x=y=0`$ and $`xtz=y+t^2w=0`$ respectively. Show that the flat limit as $`t0`$ of the union $`LM`$ is the double line $`x^2=y=0`$ with an embedded point of degree $`1`$ located at the point $`(0:0:0:1)`$.
###### Solution.
We find the flat limit by saturating the intersection ideal:
```
i1 : PP3 = QQ[x, y, z, w];
```
```
i2 : S = QQ[t, x, y, z, w];
```
```
i3 : phi = map(PP3, S, 0 | vars PP3 );
```
```
o3 : RingMap PP3 <--- S
```
```
i4 : L = ideal(x, y);
```
```
o4 : Ideal of S
```
```
i5 : M = ideal(x-t*z, y-t^2*w);
```
```
o5 : Ideal of S
```
```
i6 : X = intersect(L, M);
```
```
o6 : Ideal of S
```
```
i7 : Xzero = trim phi substitute(saturate(X, t), {t => 0})
```
```
2 2
o7 = ideal (y*z, y , x*y, x )
```
```
o7 : Ideal of PP3
```
This is the union of a double line and an embedded point of degree $`1`$.
```
i8 : use PP3;
```
```
i9 : intersect(ideal(x^2, y), ideal(x, y^2, z))
```
```
2 2
o9 = ideal (y*z, y , x*y, x )
```
```
o9 : Ideal of PP3
```
```
i10 : degree( ideal(x^2, y) / ideal(x, y^2, z))
```
```
o10 = 1
```
## Bézout’s Theorem
Bézout’s theorem (Theorem III-78 in Eisenbud and Harris, 1999) fails without the Cohen-Macaulay hypothesis. Following Exercise III-81 in Eisenbud and Harris (1999), we illustrate this in Macaulay2.
###### Question 7.
Find irreducible closed subvarieties $`X`$ and $`Y`$ in $`^4`$ with
$`\mathrm{codim}(XY)`$ $`=`$ $`\mathrm{codim}(X)+\mathrm{codim}(Y)\text{ and}`$
$`\mathrm{deg}(XY)`$ $`>`$ $`\mathrm{deg}(X)\mathrm{deg}(Y).`$
###### Solution.
We show that the assertion holds when $`X`$ is the cone over the nonsingular rational quartic curve in $`^3`$ and $`Y`$ is a two-plane passing through the vertex of the cone. The computation is done as follows:
```
i1 : S = QQ[a, b, c, d, e];
```
```
i2 : quarticCone = trim minors(2,
matrix{{a, b^2, b*d, c}, {b, a*c, c^2, d}})
```
```
3 2 2 2 3 2
o2 = ideal (b*c - a*d, c - b*d , a*c - b d, b - a c)
```
```
o2 : Ideal of S
```
```
i3 : plane = ideal(a, d);
```
```
o3 : Ideal of S
```
```
i4 : codim quarticCone + codim plane == codim (quarticCone + plane)
```
```
o4 = true
```
```
i5 : (degree quarticCone) * (degree plane)
```
```
o5 = 4
```
```
i6 : degree (quarticCone + plane)
```
```
o6 = 5
```
## Constructing Blow-ups
The blow-up of a scheme $`X`$ along a subscheme $`Y`$ can be constructed from the Rees algebra associated to the ideal sheaf of $`Y`$ in $`X`$; see Theorem IV-22 in Eisenbud and Harris (1999). Gröbner basis techniques allow one to express the Rees algebra in terms of generators and relations. We demonstrate this by solving Exercise IV-43 in Eisenbud and Harris (1999).
###### Question 8.
Find the blow-up of the affine plane $`𝔸^2=\mathrm{Spec}\left(k[x,y]\right)`$ along the subscheme defined by $`x^3,xy,y^2`$.
###### Solution.
We first provide a general function which given an ideal and a list of variables returns the ideal of relations for the Rees algebra.
```
i1 : blowUpIdeal = method();
```
```
i2 : blowUpIdeal(Ideal, List) := (I, L) -> (
r := numgens I;
S := ring I;
kk := coefficientRing S;
n := numgens S;
y := symbol y;
St := kk[t, gens S , y_1..y_r, MonomialOrder => Eliminate 1];
phi := map(St, S, submatrix(vars St, {1..n}));
F := phi gens I;
local J;
J = ideal apply(1..r, j -> y_j - t*F_(0, (j-1)));
J = ideal selectInSubring(1, gens gb J);
if (#L < r) then error "not enough variables";
R := kk[gens S, L];
theta := map(R, St, 0 | vars R);
theta J
);
```
Applying the function to our specific case yields:
```
i3 : S = QQ[x, y];
```
```
i4 : I = ideal(x^3, x*y, y^2);
```
```
o4 : Ideal of S
```
```
i5 : blowUpIdeal(I, {A, B, C})
```
```
2 2 3 2
o5 = ideal (y*B - x*C, x*B - A*C, x B - y*A, x C - y A)
```
```
o5 : Ideal of QQ [x, y, A, B, C]
```
## A Classic Blow-up
We consider the blow-up of the projective plane $`^2`$ at a point. Many related examples appear in section IV.2.2 of Eisenbud and Harris (1999).
###### Question 9.
Show that the following varieties are isomorphic.
1. the image of the rational map from $`^2`$ to $`^4`$ given by
$$(r:s:t)(r^2:s^2:rs:rt:st);$$
2. the blow-up of the plane $`^2`$ at the point $`(0:0:1)`$;
3. the determinantal variety defined by the $`2\times 2`$ minors of the matrix
$$\left[\begin{array}{ccc}a& c& d\\ b& d& e\end{array}\right]$$
where $`^4=\mathrm{Proj}\left(k[a,b,c,d,e]\right)`$.
This surface is called the cubic scroll in $`^4`$.
###### Solution.
We find the ideal in part (a) by elimination theory.
```
i1 : PP4 = QQ[a, b, c, d, e];
```
```
i2 : S = QQ[r, s, t, A..E, MonomialOrder => Eliminate 3 ];
```
```
i3 : I = ideal(A - r^2, B - s^2, C - r*s, D - r*t, E - s*t);
```
```
o3 : Ideal of S
```
```
i4 : phi = map(PP4, S, (matrix{{0, 0, 0}}**PP4) | vars PP4)
```
```
o4 = map(PP4,S,{0, 0, 0, a, b, c, d, e})
```
```
o4 : RingMap PP4 <--- S
```
```
i5 : surfaceA = phi ideal selectInSubring(1, gens gb I)
```
```
2
o5 = ideal (c*d - a*e, b*d - c*e, a*b - c )
```
```
o5 : Ideal of PP4
```
We determine the surface in part (b) by constructing the blow-up of $`^2`$ in $`^2\times ^1`$ and then projecting its Segre embedding from $`^5`$ into $`^4`$. Notice that its image under the Segre map lies on a hyperplane in $`^5`$.
```
i6 : R = QQ[t, x, y, z, u, v, MonomialOrder => Eliminate 1];
```
```
i7 : blowUpIdeal = ideal selectInSubring(1, gens gb ideal(u-t*x, v-t*y))
```
```
o7 = ideal(y*u - x*v)
```
```
o7 : Ideal of R
```
```
i8 : PP2xPP1 = QQ[x, y, z, u, v];
```
```
i9 : psi = map(PP2xPP1, R, 0 | vars PP2xPP1);
```
```
o9 : RingMap PP2xPP1 <--- R
```
```
i10 : blowUp = PP2xPP1 / psi(blowUpIdeal);
```
```
i11 : PP5 = QQ[A, B, C, D, E, F];
```
```
i12 : segre = map(blowUp, PP5, matrix{{x*u, y*u, z*u, x*v, y*v, z*v}});
```
```
o12 : RingMap blowUp <--- PP5
```
```
i13 : ker segre
```
```
2
o13 = ideal (B - D, C*E - D*F, D - A*E, C*D - A*F)
```
```
o13 : Ideal of PP5
```
```
i14 : theta = map( PP4, PP5, matrix{{a, c, d, c, b, e}})
```
```
o14 = map(PP4,PP5,{a, c, d, c, b, e})
```
```
o14 : RingMap PP4 <--- PP5
```
```
i15 : surfaceB = trim theta ker segre
```
```
2
o15 = ideal (c*d - a*e, b*d - c*e, a*b - c )
```
```
o15 : Ideal of PP4
```
Finally, we compute the surface in part (c) and apply a permutation of the variables to obtain the desired isomorphisms
```
i16 : determinantal = minors(2, matrix{{a, c, d},{b, d, e}})
```
```
2
o16 = ideal (- b*c + a*d, - b*d + a*e, - d + c*e)
```
```
o16 : Ideal of PP4
```
```
i17 : sigma = map(PP4, PP4, matrix{{d, e, a, c, b}});
```
```
o17 : RingMap PP4 <--- PP4
```
```
i18 : surfaceC = sigma determinantal
```
```
2
o18 = ideal (c*d - a*e, b*d - c*e, a*b - c )
```
```
o18 : Ideal of PP4
```
```
i19 : surfaceA == surfaceB
```
```
o19 = true
```
```
i20 : surfaceB == surfaceC
```
```
o20 = true
```
## Fano Schemes
Our last example concerns the family of Fano schemes associated to a flat family of quadrics. We solve Exercise IV-69 in Eisenbud and Harris (1999).
###### Question 10.
Consider the one-parameter family of quadrics
$$Q=V(tx^2+ty^2+tz^2+w^2)_{k[t]}^3=\mathrm{Proj}\left(k[t][x,y,z,w]\right).$$
As the fiber $`Q_t`$ tends to the double plane $`Q_0`$, what is the flat limit of the Fano scheme $`F_1(Q_t)`$ of lines lying on these quadric surfaces ?
###### Solution.
We first make the homogeneous coordinate ring of the ambient projective $`3`$-space and the ideal of our family of quadrics.
```
i1 : PP3overBase = QQ[t, x, y, z, w];
```
```
i2 : Qt = ideal(t*x^2+t*y^2+t*z^2+w^2);
```
```
o2 : Ideal of PP3overBase
```
We construct an indeterminate line in $`_{[t]}^3`$ by adding parameters $`u,v`$ and two points $`\left(A:B:C:D\right)`$ and $`\left(E:F:G:H\right)`$. The map $`\varphi `$ sends the variables to the coordinates of the general point on this line.
```
i3 : S = QQ[t, u, v, A..H];
```
```
i4 : phi = map(S, PP3overBase, matrix{{t}} |
u*matrix{{A, B, C, D}}+v*matrix{{E, F, G, H}});
```
```
o4 : RingMap S <--- PP3overBase
```
The indeterminate line is contained in our family of quadrics if and only if $`\varphi (tx^2+ty^2+tz^2+w^2)`$ vanishes identically in $`u`$ and $`v`$. Thus, we extract the coefficients of $`u`$ and $`v`$.
```
i5 : imageOfQt = phi Qt;
```
```
o5 : Ideal of S
```
```
i6 : coeff = (coefficients({1,2}, gens imageOfQt))_1;
```
```
1 3
o6 : Matrix S <--- S
```
We no longer need the variables $`u`$ and $`v`$.
```
i7 : Sprime = QQ[t, A..H];
```
```
i8 : coeff = substitute(coeff, Sprime);
```
```
1 3
o8 : Matrix Sprime <--- Sprime
```
```
i9 : Sbar = Sprime / (ideal coeff);
```
To move to the Grassmannian over $`[t]`$, we introduce a polynomial ring in $`6`$ new variables corresponding to the minors of the matrix $`\left[\begin{array}{cccc}A& B& C& D\\ E& F& G& H\end{array}\right]`$. The map $`\psi `$ sends the new variables $`a,\mathrm{}f`$ to the appropriate minor, regarded as elements of $`\mathrm{𝚂𝚋𝚊𝚛}`$.
```
i10 : PP5overBase = QQ[t, a..f];
```
```
i11 : psi = matrix{{Sbar_"t"}} | substitute(
exteriorPower(2, matrix{{A, B, C, D}, {E, F, G, H}}), Sbar);
```
```
1 7
o11 : Matrix Sbar <--- Sbar
```
```
i12 : fanoOfQt = trim ker map(Sbar, PP5overBase, psi);
```
```
o12 : Ideal of PP5overBase
```
We next determine the limit as $`t`$ tends to $`0`$.
```
i13 : fanoOfQ0 = trim substitute(saturate(
fanoOfQt, t), {t => 0})
```
```
2 2 2
o13 = ideal (e*f, d*f, d*e, a*e + b*f, d , f , e , c*d - b*e + a*f,
```
```
2 2 2
b*d + c*e, a*d - c*f, a + b + c )
```
```
o13 : Ideal of PP5overBase
```
```
i14 : degree( ideal(d, e, f, a^2+b^2+c^2) / fanoOfQ0)
```
```
o14 = 2
```
We see that $`F_1(Q_0)`$ is supported on the plane conic $`d,e,f,a^2+b^2+c^2`$ and $`F_1(Q_0)`$ is not reduced — it has multiplicity two.
From section IV.3.2 in Eisenbud and Harris (1999), we know that $`F_1(Q_1)`$ is the union of two conics lying in complementary planes. We verify this as follows:
```
i15 : fanoOfQ1 = trim substitute(saturate(
fanoOfQt, t), {t => 1})
```
```
2 2 2
o15 = ideal (a*e + b*f, d + e + f , c*d - b*e + a*f, b*d + c*e,
```
```
2 2 2 2 2
a*d - c*f, c + e + f , b*c + d*e, a*c - d*f, b - e ,
```
```
2 2
a*b + e*f, a - f )
```
```
o15 : Ideal of PP5overBase
```
```
i16 : fanoOfQ1 == intersect(ideal(c-d, b+e, a-f, d^2+e^2+f^2),
ideal(c+d, b-e, a+f, d^2+e^2+f^2))
```
```
o16 = true
```
Thus, $`F_1(Q_0)`$ is the double conic obtained when the two conics in $`F_1(Q_1)`$ move together.
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# Implementation of the Projector Augmented Wave LDA+U Method: Application to the Electronic Structure of NiO
## I Introduction
For many materials, the density functional theory (DFT) in the local-spin-density approximation (LSDA) provides a good description of their ground-state properties. However, problems arise when the DFT-LSDA approach is applied to materials with ions that contain incomplete $`d`$\- or $`f`$-shells, such as transition-metal oxides or heavy fermion systems. For example, most transition-metal oxides are wide-gap antiferromagnetic insulators, and the DFT-LSDA predicts them to be either metals (FeO and CoO) or small-gap semiconductors (MnO and NiO). The failure of the DFT-LSDA can be traced to the mean-field character of the Kohn-Sham equations as well as to the poor description of strong correlation effects within the homogeneous electron gas. The strong correlation effects are responsible for the breakdown of the DFT-LSDA description of the electronic structure of these compounds. In order to provide a better description of these effects, the Mott-Hubbard picture has been introduced.
In the Mott-Hubbard picture of NiO, the $`d`$-$`d`$ Coulomb interaction splits the Ni $`d`$ sub-bands into the so-called lower and upper Hubbard bands. The upper Hubbard band has mostly Ni $`3d^9`$ character, while the top of the valence band is of $`3d^8`$ character, leading to a Mott-Hubbard gap of $`d`$-$`d`$ type. However, O $`1s`$ x-ray absorption as well as x-ray photoemission and bremsstrahlung isochromat spectroscopies on Li<sub>x</sub>Ni<sub>1-x</sub>O have shown that the additional hole has mainly oxygen character. In contrast to the Mott-Hubbard model, the energy band gap caused by the Ni $`3d`$ correlations is therefore of the charge-transfer type between the occupied oxygen $`2p`$ and the Ni $`3d`$ empty states.
On the other hand, localized approaches, such as the local cluster scheme based on the configuration interaction method or the Anderson impurity model, in which transition-metal ions are treated like an impurity in an oxygen $`2p`$ host, predict a well-defined band gap of 5.0 eV. However, the oxygen $`2p`$ band dispersion observed in angle-resolved photoemission spectroscopy cannot be described by these methods because the lattice effects are neglected.
Several attempts have been made to include the missing correlation effects in DFT-LSDA. The generalized gradient approximation (GGA), which takes into account the radial and angular gradient corrections, can only open a small band gap. The self-interaction correction (SIC) eliminates the spurious interaction of an electron with itself from the conventional DFT-LSDA method. Compared to LSDA, the band gap and the magnetic moments are significantly increased. However, the band gap still is too small, and the SIC-LSDA method predicts a larger energy band gap for NiO than for FeO and CoO, in contradiction to experiment. The crystal-field orbital polarization introduced by Norman to determine the magnetism and insulating band gap of transition-metal oxides is promising but underestimates both the spin magnetic moment and the band gap.
Another promising approach for correlated materials is the so-called local density approximation (LDA) plus the multi-orbital mean-field Hubbard model (LDA+U) which includes the on-site Coulomb interaction in the LSDA Hamiltonian. After adding the on-site Coulomb interaction to the LSDA Hamiltonian, the potential becomes spin- and orbital-dependent. Because a larger energy cost is associated with fluctuations of the $`d`$-occupancy, the orbital-dependent potential reduces the fluctuations of the $`d`$-occupancy, resulting in a better justification of a mean-field approach. LDA+U, although it is a mean-field approach, has the advantage of describing both the chemical bonding and the electron-electron interaction.
The question regarding the best value for the Coulomb repulsion parameter $`U`$ is, however, still under debate. The $`U`$ parameter for NiO obtained from a constrained LDA calculation is about 7 to 8 eV, and this is the value generally used in LDA+U calculations. A similar value of $`U`$ has been obtained from a constrained LDA calculation for bulk Fe, even though a much smaller value had been expected because of the metallic screening in Fe. The authors argued that the higher value could be an artifact due to the poor screening within the atomic sphere approximation (ASA), and that within a full-potential calculation a much smaller value of less than 4 eV would be expected. In contrast, an unpublished full-potential linear muffin-tin orbital method calculation by Alouani and Wills clearly shows that the value of $`U`$ for bulk Fe is even slightly larger than the ASA value. Furthermore, Bulut, Scalapino, and White showed that the renormalization of the Coulomb interaction depends on the type of model used. As LDA is not a diagrammatic method, it is not known which type of renormalization is the most appropriate for the LDA+U model.
In this work we shed light on this problem by treating the Coulomb repulsion parameter $`U`$ as adjustable parameter, and by investigating how the electronic and optical properties depend on its value. We show that for an intermediate value of $`U=5`$ eV, good agreement with the measured ground-state antiferromagnetic magnetic moment and optical properties is obtained. We also show that the O $`2p`$ character near the top of the valence states is enhanced for a larger value of $`U`$. Our calculation seems to indicate that the nature of the band gap at intermediate $`U`$ is a mixture of charge transfer and Mott-Hubbard type, and that it becomes almost purely of the charge-transfer type for higher values of $`U`$.
Our calculations are based on the projector augmented wave (PAW) method, an efficient all-electron method without shape approximations on the potential or electron density to avoid uncertainties due to the ASA approach. Based on a Car-Parrinello-like formalism, the PAW method allows complex relaxations and dynamical properties in strongly correlated systems to be studied. Our implementation of LDA+U within the PAW method is described in detail. Furthermore we discuss possible extensions of the existing method that will enhance its applicability.
The paper is organized as follows: In Section II we present those aspects of the PAW formalism that are needed for the implementation of the LDA+U method. In Section III we present and discuss the LSDA and LDA+U ground-state properties of NiO, and in Section IV we study the optical properties of NiO, namely the imaginary part of the dielectric function, and compare the results to experiment.
## II Formalism
### A PAW method
The PAW method developed by one of us combines ideas of the pseudo-potential (PP) and the linear augmented plane wave (LAPW) methods. It is applicable to all elements of the periodic table. The nodal behavior of the wave function is correctly described and, as in the PP method, the forces on the ions are easily expressed.
In the PAW method, the all-electron (AE) crystal wave function is constructed from a pseudo (PS) wave function and atom-like functions localized near the nuclei. The PS wave function $`|\stackrel{~}{\mathrm{\Psi }}`$ coincides with the crystal AE wave function $`|\mathrm{\Psi }`$ in the interstitial region, i.e. outside the atomic regions. Inside the atomic regions $`\mathrm{\Omega }_t`$, called augmentation regions, the wave function is almost atom-like because the effect of the surrounding crystal is small. Therefore, a natural choice is to use solutions $`|\varphi _\mathrm{\Lambda }`$ of Schrödinger’s equation for the isolated atom, the so-called AE partial waves, as a basis set for the augmentation region. Here $`\mathrm{\Lambda }=\{t,\alpha ,\mathrm{},m\}`$ is a global index for the atom $`t`$, the angular momentum $`\mathrm{}`$, the magnetic quantum number $`m`$, and the index $`\alpha `$, the energy for which Schrödinger’s equation is solved.
To link the expansion in atom-like functions near the nuclei to the PS wave function, we introduce a set of auxiliary functions $`|\stackrel{~}{\varphi }_\mathrm{\Lambda }`$, so-called PS partial waves, which are centered on the atom and coincide per construction with the corresponding AE partial waves $`|\varphi _\mathrm{\Lambda }`$ outside their augmentation regions:
$$\varphi _\mathrm{\Lambda }(𝐫)=\stackrel{~}{\varphi }_\mathrm{\Lambda }(𝐫)\text{for}r\mathrm{\Omega }_t.$$
(1)
The coefficients $`c_\mathrm{\Lambda }`$ of the expansions in AE and PS partial waves are chosen such that the PS partial wave expansion $`_\mathrm{\Lambda }|\stackrel{~}{\varphi }_\mathrm{\Lambda }c_\mathrm{\Lambda }`$ cancels out the PS wave function $`|\stackrel{~}{\mathrm{\Psi }}`$ inside the augmentation region. For this purpose we introduce so-called projector functions $`\stackrel{~}{p}|`$ such that
$$\underset{\mathrm{\Lambda }}{}|\stackrel{~}{\varphi }_\mathrm{\Lambda }\stackrel{~}{p}_\mathrm{\Lambda }|=1,$$
(2)
and therefore
$$|\stackrel{~}{\mathrm{\Psi }}=\underset{\mathrm{\Lambda }}{}|\stackrel{~}{\varphi }_\mathrm{\Lambda }\stackrel{~}{p}_\mathrm{\Lambda }|\stackrel{~}{\mathrm{\Psi }}$$
(3)
for the Hilbert space spanned by the PS partial waves $`|\stackrel{~}{\varphi }_\mathrm{\Lambda }`$. Thus we identify the expansion coefficients with $`c_\mathrm{\Lambda }=\stackrel{~}{p}_\mathrm{\Lambda }|\stackrel{~}{\mathrm{\Psi }}`$. Equation (2) results in the biorthogonality condition
$$\stackrel{~}{p}_\mathrm{\Lambda }|\stackrel{~}{\varphi }_\mathrm{\Lambda }^{}=\delta _{\mathrm{\Lambda },\mathrm{\Lambda }^{}}$$
(4)
for the projector functions, which moreover are chosen to be localized within the corresponding augmentation region.
With these conditions, the AE Bloch wave function $`\mathrm{\Psi }(𝐫)`$ can be obtained from the PS wave function $`\stackrel{~}{\mathrm{\Psi }}(𝐫)`$ as
$$\mathrm{\Psi }(𝐫)=\stackrel{~}{\mathrm{\Psi }}(𝐫)+\underset{\mathrm{\Lambda }}{}[\varphi _\mathrm{\Lambda }(𝐫)\stackrel{~}{\varphi }_\mathrm{\Lambda }(𝐫)]\stackrel{~}{p}_\mathrm{\Lambda }|\stackrel{~}{\mathrm{\Psi }}.$$
(5)
The first term represents the PS wave function defined over the entire space, which is equal to the AE wave function in the interstitial region, and which is expanded in plane waves. The second term is the AE partial wave expansion, which describes the correct nodal behavior of the wave function in the augmentation region $`\mathrm{\Omega }_R^t`$ ($`rr_c^t`$). The third term eliminates the spurious contribution of the PS wave function in the augmentation region.
Note that Eq. (3) holds only approximately if the set of PS partial waves is not entirely within the augmentation regions. However, this has the advantage that only those contributions of the PS wave function will be removed that are also replaced by AE partial waves. As a result, the AE wave function converges rapidly with the number of partial waves used and, moreover, it is continuous and differentiable for every truncation of the partial-wave expansion.
Expectation values of any sufficiently local operator $`A`$ are obtained as
$$\mathrm{\Psi }|A|\mathrm{\Psi }=\stackrel{~}{\mathrm{\Psi }}|A|\stackrel{~}{\mathrm{\Psi }}+\underset{\mathrm{\Lambda },\mathrm{\Lambda }^{}}{}\stackrel{~}{\mathrm{\Psi }}|\stackrel{~}{p}_\mathrm{\Lambda }\left(\varphi _\mathrm{\Lambda }|A|\varphi _\mathrm{\Lambda }^{}\stackrel{~}{\varphi }_\mathrm{\Lambda }|A|\stackrel{~}{\varphi }_\mathrm{\Lambda }^{}\right)\stackrel{~}{p}_\mathrm{\Lambda }^{}|\stackrel{~}{\mathrm{\Psi }}.$$
(6)
Note that the double sum is diagonal in the site indices $`t,t^{}`$. This equation is exact for a complete set of PS partial waves and rapidly attains the converged result if incomplete. The PAW method provides the freedom to represent a zero operator in the form
$$0=\stackrel{~}{\mathrm{\Psi }}|B|\stackrel{~}{\mathrm{\Psi }}\underset{\mathrm{\Lambda },\mathrm{\Lambda }^{}}{}\stackrel{~}{\mathrm{\Psi }}|\stackrel{~}{p}_\mathrm{\Lambda }\stackrel{~}{\varphi }_\mathrm{\Lambda }|B|\stackrel{~}{\varphi }_\mathrm{\Lambda }^{}\stackrel{~}{p}_\mathrm{\Lambda }^{}|\stackrel{~}{\mathrm{\Psi }}$$
(7)
by any operator $`B`$ entirely localized within the augmentation region. Equation (7) has the same range of validity as Eq. (6) does and allows a further acceleration of the convergence by using a well-chosen operator $`B`$ and adding the corresponding zero operator to the expression for the expectation value.
### B The LDA+U total energy functional
For transition-metal oxides, the $`d`$-orbitals are well localized and keep a strong atom-like character. Even though LDA provides a good approximation for the average Coulomb energy of the $`d`$-$`d`$ interactions, it fails to describe correctly the strong Coulomb and exchange interaction between electrons in the same $`d`$-shell. The main intention of LDA+U is to identify these atomic orbitals and to describe their electronic interactions as strongly correlated states. The other orbitals are delocalized and considered to be properly described by the LDA. The procedure is to eliminate the averaged LDA energy contribution of these atom-like orbitals from the LDA total energy functional $`E^{\mathrm{LDA}}`$, and to add an orbital- and spin-dependent correction. The total energy within the LDA+U method then has the form
$`E`$ $`=`$ $`E^{\mathrm{LDA}}+{\displaystyle \frac{1}{2}}{\displaystyle \underset{t,\sigma }{}}{\displaystyle \underset{i,j,k,l}{}}\chi _i^t;\chi _k^t|V_{ee}|\chi _j^t;\chi _l^tn_{i,j}^{t,\sigma }n_{k,l}^{t,\sigma }`$ (8)
$`+`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{t,\sigma }{}}{\displaystyle \underset{i,j,k,l}{}}\left(\chi _i^t;\chi _k^t|V_{ee}|\chi _j^t;\chi _l^t\chi _i^t;\chi _k^t|V_{ee}|\chi _l^t;\chi _j^t\right)n_{i,j}^{t,\sigma }n_{k,l}^{t,\sigma }`$ (9)
$``$ $`{\displaystyle \underset{t}{}}\left[{\displaystyle \frac{1}{2}}U{\displaystyle \underset{\sigma ,\sigma ^{}}{}}N^{t,\sigma }(N^{t,\sigma ^{}}\delta _{\sigma ,\sigma ^{}}){\displaystyle \frac{1}{2}}J{\displaystyle \underset{\sigma }{}}N^{t,\sigma }(N^{t,\sigma }1)\right],`$ (10)
where $`N^{t,\sigma }=_in_{i,i}^{t,\sigma }`$ is the average occupation of the $`d`$-shell for each spin direction as obtained from the $`d`$-orbital occupancies $`n_{i,j}^{t,\sigma }`$. $`U`$ and $`J`$ are the Coulomb self-energy and the exchange parameter, respectively. The expressions $`\chi _i^t;\chi _k^t|V_{ee}|\chi _j^t;\chi _l^t`$ are the four-center matrix elements of the screened Coulomb interaction $`V_{ee}`$. An additional requirement of the LDA+U approach is that the additional energy is applied only to the valence electrons, which are re-optimized while constrained to remain orthogonal to the core states.
### C Orbital occupations
The rationale behind the definition of the orbital occupations is to project the density matrix onto one consisting of a restricted set of localized orbitals $`|\chi _i^t`$, where each orbital index $`i`$ stands for an angular momentum and $`t`$ for an atomic site. Here we focus in particular on the $`d`$-orbitals of Ni.
The projected density matrix operator has the form
$$\overline{\rho }^\sigma =\underset{i,j}{}|\chi _i^tn_{i,j}^{t,\sigma }\chi _j^t|.$$
(11)
The density matrix elements $`n_{i,j}^t`$ in the localized basis set can be obtained from the condition that the expectation values of a set of projection operators $`P_s`$ are optimally reproduced by the localized basis set, i.e. the weighted square deviation
$`F`$ $`=`$ $`{\displaystyle \underset{s}{}}w_s\left[{\displaystyle \underset{n,𝐤}{}}\mathrm{\Psi }_n^𝐤|P_s^t|\mathrm{\Psi }_n^{𝐤,\sigma }{\displaystyle \underset{i,j}{}}n_{i,j}^{t,\sigma }\chi _i^t|P_s^t|\chi _j^t\right]^2`$
of the projections should be minimized.
Here we choose the projection operators acting on the Ni $`d`$-orbitals in analogy to previous implementations as
$$P_{m,m^{}}^t(𝐫,𝐫^{})=\theta _{\mathrm{\Omega }_t}(𝐫)\delta (|𝐫^{}𝐑_𝐭||𝐫𝐑_𝐭|)Y_{d,m}^t(𝐫)Y_{d,m^{}}^{}(𝐫^{}),$$
(12)
where the site index $`t`$ refers to a particular Ni site, and $`|Y_{d,m}^t`$ is the spherical harmonic centered at site $`t`$. The step functions $`\theta _{\mathrm{\Omega }_t}(r)`$ are unity for $`|rR_t|<r_c^t`$ and zero otherwise. $`Y_{d,m}`$ are the spherical harmonics for $`\mathrm{}=2`$. (Note that $`\chi |P|\chi ^{}=𝑑𝐫𝑑𝐫^{}\chi (𝐫)P(𝐫,𝐫^{})\chi ^{}(𝐫^{})`$.) The radius $`r_c^t`$ for Ni corresponds to an atomic sphere radius, and we have chosen $`r_c^t=2.1`$ a<sub>0</sub>. The localized orbitals $`|\chi _m^t`$ have been chosen to be identical to those of the spheridized, non-spin-polarized atoms. We used an atom in the $`3d^84s^2`$ configuration of Ni. For such a localized Ni $`d`$-orbital $`|\chi _m^t`$ with magnetic quantum number $`m`$, we obtain
$$I=\chi _m^t|P_{m,m}^t|\chi _m^t=\chi _m^t|\theta _{\mathrm{\Omega }_t}(r)|\chi _m^t,$$
(13)
independent of $`m`$ and $`t`$. Thus we obtain the orbital occupations
$$n_{m,m^{}}^{t,\sigma }=\frac{1}{I}\underset{n,𝐤}{}\mathrm{\Psi }_n^{𝐤,\sigma }|P_{m,m^{}}^t|\mathrm{\Psi }_n^{𝐤,\sigma }$$
(14)
from the minimal condition for $`F`$, which in this case are independent of the weights $`w_s`$.
In the PAW method we obtain the orbital occupations directly from the PS wave functions $`|\stackrel{~}{\mathrm{\Psi }}_n^{𝐤,\sigma }`$ as
$$n_{m,m^{}}^{t,\sigma }=\frac{1}{I}\underset{𝐤,n}{}\stackrel{~}{\mathrm{\Psi }}_n^{𝐤,\sigma }|\stackrel{~}{P}_{m,m^{}}^t|\stackrel{~}{\mathrm{\Psi }}_n^{𝐤,\sigma },$$
(15)
using the pseudo-version
$$\stackrel{~}{P}_{m,m^{}}^t=\underset{\mathrm{\Lambda },\mathrm{\Lambda }^{}}{}|\stackrel{~}{p}_\mathrm{\Lambda }\varphi _\mathrm{\Lambda }|P_{m,m^{}}^t|\varphi _\mathrm{\Lambda }^{}\stackrel{~}{p}_\mathrm{\Lambda }^{}|+\mathrm{\Delta }\stackrel{~}{P}_{m,m^{}}^t$$
(16)
of the projection operator $`P_{m,m^{}}^t`$. Consistent with the PAW formalism, the small correction
$$\mathrm{\Delta }\stackrel{~}{P}_{m,m^{}}^t=P_{m,m^{}}^t\underset{\mathrm{\Lambda },\mathrm{\Lambda }^{}}{}|\stackrel{~}{p}_\mathrm{\Lambda }\stackrel{~}{\varphi }_\mathrm{\Lambda }|P_{m,m^{}}^t|\stackrel{~}{\varphi }_\mathrm{\Lambda }^{}\stackrel{~}{p}_\mathrm{\Lambda }^{}|$$
(17)
is ignored in the present calculations because it can be considered the pseudo-version of the zero operator.
Note that the present choice for the projection operator may not be the most elegant form. Here we suggest another version using projector functions $`q_i^t|`$. The construction is analogous that of the projector functions in the PAW method, but the purpose is different: namely, the goal is to decompose a wave function into local orbitals such that
$$\underset{i,t}{}|\chi _i^tq_i^t|\mathrm{\Psi }=|\mathrm{\Psi }$$
(18)
holds for any wave function $`|\mathrm{\Psi }`$ that can be represented exactly as a superposition of the local orbitals $`|\chi _i^t`$. This requirement is fulfilled whenever $`q_i^t|\chi _j^t^{}=\delta _{i,j}\delta _{t,t^{}}`$. The projector functions can be obtained from a set of functions $`|f_i^t`$, which are linearly independent and localized within an augmentation region as
$$\underset{j,t^{}}{}f_i^t|\chi _j^t^{}q_j^t^{}|=f_i^t|$$
(19)
after multiplication with the inverse of $`f_i^t|\chi _j^t^{}`$. If we now use
$$\overline{P}_{i,j}^t=|q_i^tq_j^t|,$$
(20)
we obtain the matrix elements in the representation of localized orbitals directly as
$$n_{i,j}^t=\underset{k,l,t^{\prime \prime },t^{\prime \prime \prime }}{}n_{k,l}^{t^{\prime \prime },t^{\prime \prime \prime }}\chi _k^{t^{\prime \prime }}|q_i^tq_j^t|\chi _l^{t^{\prime \prime \prime }}=\underset{𝐤,n}{}\mathrm{\Psi }_n^𝐤|\overline{P}_{i,j}^t|\mathrm{\Psi }_n^𝐤.$$
(21)
It is now a simple matter to obtain the corresponding pseudo-projection operator
$$\stackrel{~}{\overline{P}}_{i,j}^t=\underset{\mathrm{\Lambda },\mathrm{\Lambda }^{}}{}|\stackrel{~}{p}_\mathrm{\Lambda }\varphi _\mathrm{\Lambda }|q_i^tq_j^t|\varphi _\mathrm{\Lambda }^{}\stackrel{~}{p}_\mathrm{\Lambda }^{}|.$$
(22)
This expression has a number of advantages: (i) It is the most general expression because any operator can be expressed in a fully separable expression, and, unlike Eq. (12), it is not limited to a semi-local form. (ii) The step function in Eq. (12) introduces a potential that is discontinuous, whereas the operator suggested here can be constructed such that the resulting potential is continuous and differentiable, while still being fully localized in a well-defined region. (iii) The separable form of the operator renders the evaluation numerically convenient. In particular, it is conceivable to evaluate $`\mathrm{\Delta }\stackrel{~}{P}`$ explicitly if necessary, whereas the corresponding operation using the semi-local projection operator would be computationally prohibitive. (iv) We can systematically construct projection operators that specifically act on a well-defined portion of the Hilbert space: By enlarging the set of functions $`|\chi _i`$ to include also orbitals that should not be represented in the restricted set, they are excluded by simply setting the corresponding matrix elements of the density matrix $`N_{i,j}`$ to zero.
### D Coulomb and exchange parameters
The four-center integrals used in the expression of the LDA+U total energy are defined as
$$\chi _i^t;\chi _j^t|V_{ee}|\chi _k^t;\chi _l^t=𝑑𝐫_1𝑑𝐫_2\chi _i^t(𝐫_1)\chi _j^t(𝐫_2)v_{ee}(𝐫_\mathrm{𝟏},𝐫_\mathrm{𝟐})\chi _k^t(𝐫_1)\chi _l^t(𝐫_2),$$
(23)
where $`v_{ee}(𝐫,𝐫^{})`$ is the screened Coulomb interaction between two electrons.
If we choose localized $`d`$-orbitals that are described by an atomic $`d`$-wave function $`\chi _m^t(𝐫)=\chi _d(|𝐫𝐑_t|)Y_{d,m}(𝐫𝐑_t)`$ with magnetic quantum number $`m`$ (where $`Y_{d,m}`$ are the spherical harmonics for $`\mathrm{}=2`$) and furthermore assume that the static dielectric function $`ϵ`$ is constant in space, we can exploit the multipole expansion of $`1/|𝐫_1𝐫_2|`$:
$$v_{ee}(𝐫_1,𝐫_2)=\frac{1}{ϵ\left|𝐫_1𝐫_2\right|}=\frac{1}{ϵ}\underset{\mathrm{}=0}{\overset{\mathrm{}}{}}\frac{4\pi }{2\mathrm{}+1}\frac{r_<^{\mathrm{}}}{r_>^{\mathrm{}+1}}\underset{m=\mathrm{}}{\overset{+\mathrm{}}{}}Y_\mathrm{}m(𝐫_\mathrm{𝟏})Y_\mathrm{}m^{}(𝐫_\mathrm{𝟐}).$$
(24)
Here $`r_<`$ and $`r_>`$ denote the smallest and largest values of $`r_1`$ and $`r_2`$, respectively. Under these assumptions we can transform Eq. (23) into
$$\chi _1^t;\chi _3^t|V_{ee}|\chi _2^t;\chi _4^t=\underset{\mathrm{}=0}{\overset{\mathrm{}}{}}\frac{4\pi }{2\mathrm{}+1}\underset{m=\mathrm{}}{\overset{+\mathrm{}}{}}\mathrm{}_1,m_1|Y_{\mathrm{},m}|\mathrm{}_2,m_2\mathrm{}_3,m_3|Y_{\mathrm{},m}^{}|\mathrm{}_4,m_4F^{\mathrm{}},$$
(25)
where $`(\mathrm{}_i,m_i)`$ are the angular momenta quantum numbers of $`|\chi _i^t`$, $`\mathrm{},m|Y_{\mathrm{}^{\prime \prime },m^{\prime \prime }}|\mathrm{},m^{}`$ the Gaunt coefficients, and $`F^{\mathrm{}}`$ the so-called screened Slater’s integrals. Because of the special properties of the Gaunt coefficients, only $`F^0,F^2,`$ and $`F^4`$ contribute to the Coulomb integrals:
$`F^{\mathrm{}}={\displaystyle \frac{1}{ϵ}}{\displaystyle _0^{\mathrm{}}}𝑑r_1{\displaystyle _0^{\mathrm{}}}𝑑r_2r_1^2r_2^2\chi _d^2(r_1)\chi _d^2(r_2){\displaystyle \frac{r_<^{\mathrm{}}}{r_>^{\mathrm{}+1}}}.`$ (26)
The parameters $`U`$ and $`J`$ are identified with averages of the Coulomb and exchange integrals, which are related to the Slater integrals $`F^0,F^2,`$ and $`F^4`$ by the properties of Clebsch-Gordan coefficients, Eq. (25):
$`U`$ $`=`$ $`{\displaystyle \frac{1}{(2\mathrm{}+1)^2}}{\displaystyle \underset{m,m^{}}{}}\chi _m^t;\chi _m^{}^t|V_{ee}|\chi _m^t;\chi _m^{}^t=F^0,`$ (27)
$`J`$ $`=`$ $`{\displaystyle \frac{1}{(2\mathrm{})(2\mathrm{}+1)}}{\displaystyle \underset{mm^{},m^{}}{}}\chi _m^t;\chi _m^{}^t|V_{ee}|\chi _m^{}^t;\chi _m^t=(F^2+F^4)/14,`$ (28)
The dielectric constant and therefore the Coulomb and exchange parameters $`U`$ and $`J`$ are not known a priori. Usually, they are obtained from a constrained DFT calculation. Here, we are interested in how sensitively the results depend on the choice of Coulomb parameters, and which Coulomb parameters will provide the best agreement with reality as probed by optical absorption. Therefore, we adopt the general form for the four-center integrals as function of the $`U`$ and $`J`$ suggested by the arguments provided above, and perform calculations for different $`U`$ values, namely $`U=5`$ eV$`\times I^2`$ and $`U=8`$ eV$`\times I^2`$. In the following we suppose that $`I^2=1`$. Because the results are fairly insensitive to the exchange parameter $`J`$, we have adopted $`J=0.95`$ eV$`\times I^2`$ from previous constrained LDA calculations. The third mandatory relation is obtained from the work of DeGroot et al., who determined the ratio $`F^4/F^2`$ for transition-metal oxides to be between 0.62 and 0.63. We therefore adopt a ratio $`F^4/F^2=0.625`$.
### E Hamiltonian
The pseudo-Hamiltonian operator
$$\stackrel{~}{H}_\sigma =\stackrel{~}{H}_\sigma ^{LDA}+\stackrel{~}{H}_\sigma ^U,$$
(29)
which acts on the PS wave functions, is obtained as the derivative of the total energy functional with respect to the two-center pseudo-density matrix operator $`\stackrel{~}{\rho }=_{n,𝐤}|\stackrel{~}{\mathrm{\Psi }}_n^{𝐤,\sigma }f_n^{𝐤,\sigma }\stackrel{~}{\mathrm{\Psi }}_n^{𝐤,\sigma }|`$. The non-LDA contribution of the LDA+U Hamiltonian is then obtained as the product of the derivative $`V_{m,m^{}}^{t,\sigma }`$ of the non-LDA contribution to the total energy and the projection operator $`\stackrel{~}{P}`$, which is the derivative of the occupation with respect to the two-center pseudo-density matrix operator. Thus we obtain
$$\stackrel{~}{H}_\sigma ^U=\underset{t,m,m^{}}{}\frac{1}{I}\stackrel{~}{P}_{m,m^{}}^tV_{m^{},m}^{t,\sigma },$$
(30)
where
$`V_{m_1,m_2}^{t,\sigma }`$ $`=`$ $`{\displaystyle \underset{m_3,m_4}{}}\chi _{m_1}^t;\chi _{m_3}^t|V_{ee}|\chi _{m_2}^t;\chi _{m_4}^tn_{m_3,m_4}^{t,\sigma }`$ (31)
$`+`$ $`{\displaystyle \underset{m_3,m_4}{}}\left[\chi _{m_1}^t;\chi _{m_3}^t|V_{ee}|\chi _{m_2}^t;\chi _{m_4}^t\chi _{m_1}^t;\chi _{m_3}^t|V_{ee}|\chi _{m_4}^t;\chi _{m_2}^t\right]n_{m_3,m_4}^{t,\sigma }`$ (32)
$``$ $`{\displaystyle \underset{\sigma ^{}}{}}\left[U(N^{t,\sigma ^{}}{\displaystyle \frac{1}{2}}\delta _{\sigma ,\sigma ^{}})\delta _{\sigma ,\sigma ^{}}J(N^{t,\sigma ^{}}{\displaystyle \frac{1}{2}})\right]\delta _{m,m^{}}.`$ (33)
The LDA contribution of the pseudo-Hamiltonian has the usual form:
$`\stackrel{~}{H}_\sigma ^{LDA}={\displaystyle \frac{^2}{2}}+\stackrel{~}{v}+{\displaystyle \underset{\mathrm{\Lambda }_1,\mathrm{\Lambda }_2}{}}|\stackrel{~}{p}_{\mathrm{\Lambda }_1}\left[\varphi _{\mathrm{\Lambda }_1}\left|{\displaystyle \frac{^2}{2}}+v^1\right|\varphi _{\mathrm{\Lambda }_2}\stackrel{~}{\varphi }_{\mathrm{\Lambda }_1}\left|{\displaystyle \frac{^2}{2}}+\stackrel{~}{v^1}\right|\stackrel{~}{\varphi }_{\mathrm{\Lambda }_2}\right]\stackrel{~}{p}_{\mathrm{\Lambda }_2}|.`$ (34)
For the LDA+U calculation, we first performed a self-consistent LSDA calculation using the all-electron PAW method, and then used the self-consistent potential to construct the LSDA Hamiltonian for a large number of $`𝐤`$-points in the Brillouin zone. Next, the Hubbard correction is added to the LSDA Hamiltonian as given by Eq. (29), and the new Hamiltonian iterated until the occupation numbers $`n_{m,m^{}}^{t,\sigma }`$ have converged. We did not use the so-called second-variation procedure for self-consistent LDA+U, in which the LSDA potential is updated, because it failed to yield any improvement on calculations similar to ours.
## III Ground state of NiO
### A LSDA ground state
The ground state of NiO has been calculated using the PAW method, and the density of states (DOS) is calculated from the self-consistent PAW potential using the tetrahedron method for the Brillouin-zone integration. Figure 1 presents the atom-resolved DOS in the augmentation region. As can be seen, LSDA produced an antiferromagnetic insulating ground state with a small band gap. The LSDA magnetic moment is 0.95 $`\mu _B`$, and mainly due to the Ni$`_{e_g}`$ band splitting. This value is much smaller than the experimental value (1.64–1.90 $`\mu _B`$, Refs. ). The calculated band gap of about 0.1 eV is also much smaller than the experimental one (3.0–4.0 eV). The most interesting feature of our LSDA DOS is that the $`d`$-states of Ni dominate the region in the vicinity of the band gap, and that the top of the valence state is of Ni$`_{e_g}`$ type for the first and Ni$`_{t_2g}`$ for the second spin. This electronic structure suggests that the band gap is of Mott-Hubbard type. Hence, this LSDA picture of NiO disagrees completely with experiment. It is surprising that the quasi-particle calculation within the so-called GW approximation, performed by Aryasetiawan and Gunnarsson, produced results qualitatively similar to LSDA except for an increased band gap of 6 eV and an increased magnetic moment of 1.6 $`\mu _B`$. The quantitative change is the reduction of the O $`2p`$ bandwidth by almost 1 eV. However, a recent self-consistent model GW calculation by Massida et al., in which the dynamic effects were neglected, produced other results than those of Aryasetiawan and Gunnarsson. It was argued by Massida al. that the results of the former GW calculation are not quite self-consistent, presumably because of the additional non-local ad-hoc potential that is adjusted to the GW calculation in each self-consistent step. The main difference between the two reported GW model calculations is that the latter calculation produced (i) a spreading of the Ni $`d`$-states over the entire valence bandwidth, (ii) a vanishing gap between the O $`2p`$ and Ni $`3d`$ and, most importantly, (iii) an enhancement of O $`2p`$ states at the valence band maximum. The latter effect attributes the origin of the band gap mainly to a charge transfer gap because this gap is now between the O $`2p`$ and Ni 3$`d`$ conduction states. Next, we will show that this latter finding is in agreement with the results of the LDA+U model.
### B LDA+U ground state
We have used our implementation of the LDA+U model to determine the ground-state electronic structure of NiO. Although it is common practice to use the $`U`$ extracted from a constrained LDA calculation, we adopt a different point of view here. As stated in the introduction and in agreement with recent results reported in the literature, we believe that the value of $`U`$ extracted from constrained LDA is not the best possible choice. Therefore we have determined the electronic structure of NiO for an intermediate $`U`$ of 5 eV as well as for a larger value of 8 eV.
Figure 2 shows our calculated LDA+U DOS for $`U=5`$ and 8 eV. The energy band gap is found to be 2.8 and 4.1 eV, respectively. The total antiferromagnetic spin moment is 1.73 and 1.83 $`\mu _B`$, respectively. Our DOS obtained for $`U=8`$ eV agrees well with previous LDA+U calculations. However, the DOS for $`U`$ = 5 eV is in better agreement with the GW model calculation of Massida et al. The top of valence band is reinforced by the O $`2p`$ states, rendering the band gap a mixture of charge-transfer type and Ni $`dd`$-like excitations. In agreement with previous LDA+U calculations, the spin majority Ni $`e_g`$ states are pushed towards lower energies, and the energy difference between $`e_g^{}`$ and $`e_g^{}`$ is about 11 eV for $`U=8`$ eV and 8.6 eV for $`U=5`$ eV. Here again the GW model yields a value of about 9 eV for this splitting, in good agreement with our results for $`U`$ = 5 eV.
## IV Dielectric function
For insulators, the imaginary part of the macroscopic dielectric function is obtained within the random phase approximation (RPA) in the long wave length limit without local-field effects as
$$ϵ_2(\omega )=\underset{q0}{lim}\underset{v,c}{}\underset{𝐤}{}\frac{8\pi ^2}{\mathrm{\Omega }q^2}\left|M_{v,𝐤𝐪}^{c,𝐤}\right|^2f_{v,𝐤}(1f_{c,𝐤})\delta (E_c^𝐤E_v^𝐤\mathrm{}\omega ).$$
(35)
Here $`M_{v,𝐤𝐪}^{c,𝐤}`$ are the interband transition matrix elements, $`f_{v,𝐤}`$ is the zero-temperature Fermi distribution, $`\mathrm{\Omega }`$ is the cell volume, $`c`$ denotes the conduction-band and $`v`$ the valence-band index. In the case of a local potential, the interband transition matrix elements are given by
$$\underset{q0}{lim}M_{v,𝐤𝐪}^{c,𝐤}=\frac{𝐪}{ϵ_{c,𝐤}ϵ_{v,𝐤}}\mathrm{\Psi }_v^𝐤|𝐩|\mathrm{\Psi }_c^𝐤,$$
(36)
where the matrix elements $`\mathrm{\Psi }_v^𝐤|𝐩|\mathrm{\Psi }_c^𝐤`$ are calculated using the PAW crystal wave functions $`\mathrm{\Psi }^𝐤`$ described by Eq. (5):
$$\mathrm{\Psi }_v^𝐤|𝐩|\mathrm{\Psi }_c^𝐤=\stackrel{~}{\mathrm{\Psi }}_v^𝐤|𝐩|\stackrel{~}{\mathrm{\Psi }}_c^𝐤+\underset{\mathrm{\Lambda },\mathrm{\Lambda }^{}}{}\stackrel{~}{\mathrm{\Psi }}_v^𝐤|\stackrel{~}{p}_\mathrm{\Lambda }\left[\varphi _\mathrm{\Lambda }|𝐩|\varphi _\mathrm{\Lambda }^{}\stackrel{~}{\varphi }_\mathrm{\Lambda }|𝐩|\stackrel{~}{\varphi }_\mathrm{\Lambda }^{}\right]\stackrel{~}{p}_\mathrm{\Lambda }^{}|\stackrel{~}{\mathrm{\Psi }}_c^𝐤,$$
(37)
In the most general case, i.e. where the potential is nonlocal as in LDA+U, a nonlocal contribution has to be added to the interband transition matrix elements. The full derivation is given in the Appendix by Eq. (A7). For NiO the nonlocal contribution to the matrix elements is found to be small, i.e. of a few percent.
Figure 3 shows the imaginary part of the dielectric function calculated within the LSDA. The resulting optical spectrum is not in agreement with experiment (see Fig. 4). In particular, the optical gap is considerably underestimated and the first structure has a much higher intensity compared to experiment. Conversely, for $`U`$ = 5 eV, our calculated imaginary part of the dielectric function within the LDA+U is in a better agreement with experiment, as shown in Fig. 4. The optical band gap and the oscillator strength of the first excitation peak are in excellent agreement with experiment. However, at higher photon energies the agreement with experiment is only qualitative, which is not expected owing to the mean-field approximation of this simple model. A much higher value of $`U`$, i.e. 8 eV, produces a much larger optical gap in contrast to experiment. In agreement with our conclusion that a much smaller value of $`U`$ is required to describe NiO, Dudarev and coworkers also found that $`U=6.2`$ eV reproduces the lattice parameter and the measured electron-energy loss spectra. It is too early to draw a definitive conclusion about the excited states of NiO, as our LDA+U model is a mean-field-like model in which excitonic effects are not included. To our knowledge, no calculations of excitonic effects have so far been attempted. Our results represent the first investigation of the low-lying excited states of NiO that considers all the subtleties of chemical bonding and strong electron-electron interaction.
The interband transitions are responsible for the first structure in the optical spectrum of NiO, located between 4.1 and 4.4 eV. We found that 40.2% of the contribution results from the transition from band 15 (second highest occupied band) to band 17 (lowest empty band), 36.2% from the transition 16$``$18, and 15.9% from the interband transition 16$``$17. To analyze the character of the initial and final states of the interband transitions, the band structure dispersion along some of the high-symmetry directions is shown in Fig. 5 together with the DOS of the states that give rise to the first optical peak. The arrows between the parallel bands indicate the interband transitions from the initial to the final state responsible for the first structure in the optical spectrum. Figure 6 shows the charge density plot of the initial and final states $`\mathrm{\Psi }_n^{𝐤,\sigma }`$ for bands 16 (highest occupied band) and 18 (second lowest empty band) at point $`𝐤=(\frac{127}{120},\frac{127}{360},\frac{\sqrt{3}}{90})\frac{\pi }{a}`$ located between the high-symmetry points K and U, where the optical matrix element value is among the largest. It also shows that the initial state is of mixed O $`2p`$ and Ni $`3d`$ character, whereas the final state mainly is of Ni $`e_g`$ character. The optical transitions then are between the O $`2p`$ and the Ni $`3d`$ states, resulting in excitations of the charge-transfer type. Our interpretation differs from that of Fujimori and Minami, who used a configuration interaction within the metal-ligand cluster and claimed that the $`dd`$ charge-transfer transitions are the origin of fundamental edge. The drawback of the cluster calculation is that in reality the oxygen $`2p`$ orbitals are delocalized and therefore not well described in a small cluster. On the other hand, an earlier, band-structure-based interpretation by Messick and coworkers assigned the peak to one electron interband transitions associated with Ni $`3d`$ to the Ni $`4s`$ states. This interpretation is not correct either because the Ni $`4s`$ are far above the top of the valence states (greater than 6 eV), and only quadrupolar interband transitions are permitted between the $`3d`$ and $`4s`$, which substantially reduces the peak intensity.
## V Conclusion
We have presented a new implementation of LDA+U model based on the PAW method, which is an all-electron method without any shape approximation to the potential or the charge density. We tested the method on NiO and obtained results that are in good agreement with previous LDA+U calculations and a recent GW model calculation. In particular, we obtained the correct antiferromagnetic insulating ground state of NiO.
We discussed the results in terms of the strength of the Hubbard interaction $`U`$. The optimum value of $`U`$ has been determined by comparison with the experimental dielectric function as well as with the ground state properties. We observed a large enhancement of the O $`2p`$ character at the top of the valence state, resulting in a more charge-transfer than Ni $`dd`$ LSDA -type band gap. The calculated antiferromagnetic moment is in good agreement with experiment.
Our calculated dielectric function for an intermediate value of $`U`$, namely 5 eV, is in good agreement with experiment. The low-lying, strong structure in the optical spectrum has been assigned to an interband transition from O $`2p`$ states at top of the valence band to the Ni $`e_g`$ states at the conduction-band bottom. Hence the origin of the first optical peak is due to a charge-transfer excitation.
Our calculation is supported by a recent LDA+U calculation by Dudarev and coworkers, who also argue that a much smaller value of $`U`$ than the one obtained from constrained LDA calculation is needed to describe the electron energy loss spectra and the equilibrium lattice parameter. It should be interesting to apply this method to other transition-metal oxides and check the applicability of LSDA+U for producing excitation energies.
## VI Acknowledgment
Part of this work was done during our (O.B. and M.A.) visit to the Ohio State University in the summer of 1998. We thank J. W. Wilkins and J. G. LePage for useful discussions. Supercomputer time was provided by CINES (project gem1101) on the IBM SP2 and by the Université Louis Pasteur on the SGI O2000 supercomputer.
## A Optical transition matrix elements in LDA+U
The interband transition matrix elements for a given Hamiltonian are obtained as follows:
$`M_{v,𝐤𝐪}^{c,𝐤}`$ $`=`$ $`\mathrm{\Psi }_{v,𝐤𝐪}|e^{i\mathrm{𝐪𝐫}}|\mathrm{\Psi }_{c,𝐤}={\displaystyle \frac{\mathrm{\Psi }_{v,𝐤𝐪}|(ϵ_{v,𝐤𝐪}ϵ_{c,𝐤})e^{i\mathrm{𝐪𝐫}}|\mathrm{\Psi }_{c,𝐤}}{ϵ_{v,𝐤𝐪}ϵ_{c,𝐤}}}`$ (A1)
$`=`$ $`{\displaystyle \frac{\mathrm{\Psi }_{v,𝐤𝐪}|He^{i\mathrm{𝐪𝐫}}e^{i\mathrm{𝐪𝐫}}H|\mathrm{\Psi }_{c,𝐤}}{ϵ_{v,𝐤𝐪}ϵ_{c,𝐤}}}={\displaystyle \frac{\mathrm{\Psi }_{v,𝐤𝐪}|[H,e^{i\mathrm{𝐪𝐫}}]_{}|\mathrm{\Psi }_{c,𝐤}}{ϵ_{v,𝐤𝐪}ϵ_{c,𝐤}}}.`$ (A2)
The commutator involving of the LDA contribution $`H^{\mathrm{LDA}}`$ to the Hamiltonian $`H=H^{\mathrm{LDA}}+H^U`$ is
$$[H^{\mathrm{LDA}},e^{i\mathrm{𝐪𝐫}}]_{}=\frac{1}{2}[^2,e^{i\mathrm{𝐪𝐫}}]_{}=\frac{1}{2}(2i𝐪+q^2)=\mathrm{𝐪𝐩}+O(q)^2.$$
(A3)
The quadratic and higher-order terms in $`𝐪`$ can be ignored in the long-wavelength limit appropriate for optical transitions. The commutator involving the non-LDA part $`H^U`$ is obtained as
$$[H^U,e^{i\mathrm{𝐪𝐫}}]_{}=i𝐪[H^U,𝐫]_{}+O(q)^2=i𝐪\underset{t,m,m^{}}{}\frac{V_{m,m^{},\sigma }^t}{I}[P_{m^{},m}^t,𝐫]_{}+O(q)^2$$
(A4)
Next we use the relation which holds for the special form of the projector operator presented in Eq. (12)
$$P_{m^{\prime \prime },m^{}}^t^{}|\varphi _{t,\mathrm{},m,\alpha }=\theta _{\mathrm{\Omega }_t}(𝐫)|\varphi _{t,\mathrm{},m^{\prime \prime },\alpha }\delta _{t,t^{}}\delta _{\mathrm{},2}\delta _{m,m^{}},$$
(A5)
i.e. the projection operator changes the magnetic quantum number of a specific $`d`$-like partial wave from $`m`$ to $`m^{\prime \prime }`$ and truncates it beyond the atomic sphere $`\mathrm{\Omega }_t`$. Hence we obtain
$$\varphi _\mathrm{\Lambda }|[P_{m^{},m}^t,𝐫]_{}|\varphi _\mathrm{\Lambda }^{}^t=P_{m^{},m}^t\varphi _\mathrm{\Lambda }|\theta _{\mathrm{\Omega }_t}(𝐫)𝐫|\varphi _\mathrm{\Lambda }^{}\varphi _\mathrm{\Lambda }|\theta _{\mathrm{\Omega }_t}(𝐫)𝐫|P_{m,m^{}}^t\varphi _\mathrm{\Lambda }^{}.$$
(A6)
Finally, we obtain the expression of the matrix elements for the dipole transition, with the PAW LDA+U formalism:
$`\underset{q0}{lim}M_{v,𝐤𝐪}^{c,𝐤}`$ $`=`$ $`{\displaystyle \frac{𝐪}{(ϵ_{v,𝐤𝐪}ϵ_{c,𝐤})}}\{\stackrel{~}{\mathrm{\Psi }}_{v,𝐤}|𝐩|\stackrel{~}{\mathrm{\Psi }}_{c,𝐤}+{\displaystyle \underset{\mathrm{\Lambda },\mathrm{\Lambda }^{}}{}}\delta _{t,t^{}}\stackrel{~}{\mathrm{\Psi }}_{v,𝐤}|\stackrel{~}{p}_\mathrm{\Lambda }[\varphi _\mathrm{\Lambda }|𝐩|\varphi _\mathrm{\Lambda }^{}\stackrel{~}{\varphi }_\mathrm{\Lambda }|𝐩|\stackrel{~}{\varphi }_\mathrm{\Lambda }^{}`$ (A7)
$`+`$ $`i{\displaystyle \underset{m,m^{}}{}}{\displaystyle \frac{V_{m,m^{}}}{I}}(P_{m,m^{}}^t\varphi _\mathrm{\Lambda }|\theta _{\mathrm{\Omega }_t}(𝐫)𝐫|\varphi _\mathrm{\Lambda }^{}\varphi _\mathrm{\Lambda }|\theta _{\mathrm{\Omega }_t}(𝐫)𝐫|P_{m,m^{}}^t\varphi _\mathrm{\Lambda }^{})]\stackrel{~}{p}_\mathrm{\Lambda }^{}|\stackrel{~}{\mathrm{\Psi }}_{c,𝐤}\}.`$ (A8)
The difference between wave functions $`|\stackrel{~}{\mathrm{\Psi }}_{v,𝐤}`$ and $`|\stackrel{~}{\mathrm{\Psi }}_{v,𝐤𝐪}`$ has been ignored because it only contributes to terms that are proportional to $`q^2`$, which are ignored in the long-wavelength limit.
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# Key aspects in the photometric monitoring of 𝛿 Sct stars
## 1 Introduction
In the last two decades a great observational effort has been undertaken to solve the light curves of $`\delta `$ Sct stars. These complicated variables, located on or just above the Zero–Age–Main–Sequence, constitute the lower limit of the classical instability strip. The presence of convective zones makes these stars very interesting from the evolutionary point of view, since various stages, from pre–main sequences to giant, can be observed.
Photometric monitoring is the most practiced approach to study their properties and several stars have been deeply investigated; for a review, see Poretti (2000). However, the spectroscopic approach tells us that many modes that are not photometrically detectable are actually excited (Mantegazza 2000).
The observed frequencies are between 60 and 405 $`\mu `$Hz and multiperiodicity is very common; rotation acts as an important factor in the increase of the number of excited modes. Only a few stars show a monoperiodic behaviour above the current limit of the detectable amplitude from ground, i.e. $``$1 mmag. Of course, the number of detected modes is strongly dependent on the accuracy, the number and the sampling of the measurements.
The observed modes are in the domain of pressure ($`p`$) modes; there is no observational evidence that gravity ($`g`$) modes are excited in $`\delta `$ Sct stars, even if some cases suggest they might be. At the moment, $`g`$-modes seem to be present only in the $`\gamma `$ Dor stars, which in turn do not show $`p`$–modes.
How the photometric results can be used to identify modes (i.e. to classify the oscillation in terms of quantum numbers $`n`$, $`\mathrm{}`$ and $`m`$) is discussed by Garrido (2000).
## 2 The state–of–art in the study of $`\delta `$ Sct stars
Several approaches have been used in the past to attack the problem of the detection of excited modes in their light curves. The debate on the nature of the variations, multiperiodic or irregular, occupied the years from 1974 to 1980. At that time the observations consisted in time–series collected over a few nights, a practice that had been borrowed from the study of high amplitude $`\delta `$ Sct stars. In the years that followed, a number of teams observed more carefully $`\delta `$ Sct stars and clearly evidenced their multiperiodicity, detecting stable frequencies and obtaining very satisfactory fits. It appeared also evident that nonradial modes are excited. The number of frequencies is so high that the aliases at 1 cd<sup>-1</sup> become a serious problem in the mode detection. For that, multisite campaigns were planned and the best results now available have been obtained from these ground–based observations.
The effort made by the observers these last years allowed us to detect a large number of frequencies in many stars: we will recall here the cases of FG Vir, XX Pyx, 4 CVn and BH Psc. By means of multisite campaigns we are able to detect terms with amplitude less than 1 mmag. The mode identification techniques are based on both the phase shifts and amplitude ratios of the light and colour curves and the synergic approach performed by considering spectroscopic curves. Their full exploitation should guarantee an important role for our researches in stellar physics.
Figure 1 shows the frequency content of three well-studied pulsators in which all the terms with a half-amplitude larger than 1 mmag are detected. It appears that their content is completely different: high frequencies only for XX Pyx, low frequencies only for 4 CVn, two groups of intermediate frequency values for FG Vir. The combination terms are also reported in the upper part of each panel; 4 CVn displays both $`f_if_j`$ and $`f_i+f_j`$ terms, located right and left of the bunch of independent frequencies. FG Vir displays only three combination modes in the same region of the independent frequencies. XX Pyx displays only a 2$`f`$ harmonic, at 76.2 cd<sup>-1</sup>.
The obvious conclusion is that $`\delta `$ Sct stars are very complicated pulsators and that a general recipe cannot be applied to predict the range of the excited modes or the importance of combination terms. However, a curious similarity in the frequency content was noted by Mantegazza & Poretti (1999) comparing 4 CVn to HD 2724: at least 7 frequencies have almost the same value.
## 3 The amplitude variations
The variations of the amplitudes have been detected when a $`\delta `$ Sct star has been re–observed in order to investigate some unclear facts: a comparison of the mode amplitudes in the two different seasons often evidences strong differences. Even if the data sampling can affect the reliability of these conclusions, some results seem to be very much in favour of a physical variability (Breger 2000). These variations can be seen as a further complication, even if we can argue that the disappearance or the damping of some terms can enhance or make discernible other modes, increasing the number of known frequencies.
Amplitude variations are observed in a large variety of stars, both multiperiodic (XX Pyx, Handler et al. 1997; 4 CVn, Breger 2000; 44 Tau, Poretti et al. 1992) and monoperiodic (28 And, Rodríguez et al. 1998; BF Phe, Poretti et al. 1996). The observations of mode growth in the light curve of a relatively simple pulsator as V663 Cas can clarify what happens in much more complicated cases: some modes can be damped and then re–excited (Poretti et al. 1996). Some observational facts favour this explanation even if the model of a beating between two close frequencies with similar, constant amplitude cannot be ruled out.
BH Psc is another example. It was observed by our group in 1991 (on two occasions), 1994 and 1995 at the European Southern Observatory, Chile. The analysis of the first two campaigns showed a very complex light variability resulting from the superimposition of more than 10 pulsation modes with frequencies between 5 and 12 cd<sup>-1</sup> and semi–amplitudes between 17 and 3 mmag (Mantegazza, Poretti, & Zerbi 1995). The fit left a high r.m.s. residual 2.3 times greater than the one measured between the two comparison stars. A second photometric campaign was carried out in October and November 1994; we hoped to reveal more terms and to check the stability of their amplitudes. The analysis of the new data allowed us to single out 13 frequencies (Mantegazza, Poretti, & Bossi 1996). They are concentrated in the region 5–11.5 cd<sup>-1</sup>: this distribution is slightly larger than the one observed for 4 CVn, but they can be considered very similar. However, more terms should be present, with an amplitude below 3 mmag, since about 15% of the variance could not be explained with the detected terms and the noise. Moreover, the standard deviation around the mean value was considerably higher in 1991 than in 1994 (26 mmag against 18 mmag); considerations about the amplitudes of the modes confirmed that the pulsation energy was lower in 1994 than in 1991. Moreover, comparing the 1991 September and 1991 October data it is quite evident that we can detect such variations on a baseline of 1 month or less (see Tab. 1, where the results of a preliminary analysis are shown).
## 4 Scientific topics to be matched by MONS
The precision in the measurements achievable by MONS will of course allow us to detect terms with amplitudes much smaller than 1 mmag. This in turn will generate a much more detailed knowledge of the frequency content of $`\delta `$ Sct variables. We could then obviously match one of the requests of the asteroseismology theory, i.e. increasing the number of detected terms. This is quite obviuos. However, the values of the amplitudes of these modes are also an important parameter that can help in the theoretical models. The observational evidence is that these amplitudes are variable. Can MONS then give us an important clue as regard this aspect ?
To reach the goal, observing twice the same field can be very useful. Doing so, we could:
1. gain in frequency resolution since, as shown in Fig. 1, there are many close terms. It should also be noted that multisite campaigns span more than 50 d and that a single 30-50 d MONS run will not have a sufficient frequency resolution to solve already known close peaks. Owing to the expected large number of peaks, interference between different terms can seriously affect the quality of the results: an unappropriate frequency resolution could be considered a weakness of the MONS results;
2. detect amplitude variations in a very clear way. The absence of alias terms combined with the appropriate frequency resolution will allow a clarification on the nature of the amplitude variations, i.e. beating or intrinsic variations.
## 5 A plausible solution
In the current configuration, can MONS satisfy these two requirements ?
Among the $`\delta `$ Sct stars, FG Vir is by far the most studied: not only an excellent observational basis is available, but also more and more thorough theoretical models are proposed on the basis of the observed frequencies (Breger et al. 1999). FG Vir ($`\alpha `$ 12$`{}_{}{}^{\mathrm{h}}14_{}^{\mathrm{m}}15^\mathrm{s}`$, $`\delta 5^{}43^{}00^{\prime \prime }`$, 2000.0) is located very close to $`\beta `$ Vir ($`\alpha `$ 11$`{}_{}{}^{\mathrm{h}}50_{}^{\mathrm{m}}42^\mathrm{s}`$, $`\delta +1^{}45^{}53^{\prime \prime }`$, 2000.0), a F9 V star candidate for searching solar–like oscillations. Then, FG Vir ($`V`$=6.1) can be measured not only as a primary target with the Main Camera, but also with the Forward–pointing star tracker (BST or SI-2) when $`\beta `$ Vir is measured in the Main Camera. From this lucky coincidence, we can obtain two runs on a well–studied $`\delta `$ Sct star, greatly increasing the number of detected frequencies and saving frequency resolution. Moreover, in both cases, it will be possible to measure II Vir, another $`\delta `$ Sct star close to FG Vir.
Is that the only result ? Certainly not. There would be one more bonus. Looking through the backward–pointing star tracker (FST or SI-1) we could see, in both cases, an interesting field, located at 176 from FG and $`\beta `$ Vir. Three intriguing objects are located in this field: namely, BH Psc ($`\alpha `$ 23$`{}_{}{}^{\mathrm{h}}59_{}^{\mathrm{m}}31^\mathrm{s}`$, $`\delta 2^{}50^{}37^{\prime \prime }`$, 2000.0), i.e. a $`\delta `$ Sct showing strong amplitude variations and HD 224639 and HD 224945, two $`\gamma `$ Dor variables. The requirements about frequency resolution and amplitude variations also apply to $`\gamma `$ Dor variables.
As a conclusion, it looks possible for MONS to perform an observational study of $`\delta `$ Sct stars in a very powerful way (increase in the number of detected terms, two–colour photometry, sufficiently long baseline allowing frequency resolution and amplitude monitoring) just considering some scheduling precautions and putting an adequate target in the Main Camera.
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# Speed of ion trap quantum information processors
## I Introduction
Experimental methods which allow both coherent control, and rapid and reliable measurement of the quantum state, have been available for some time for single ions held in ion traps. In an influential paper, Cirac and Zoller showed furthermore how laser manipulation of a string of ions in a linear ion trap can achieve general coherent evolution of the joint state of the ions, using currently available technologies, and with good scaling properties. Specifically, the coherent evolution can be driven so as to realize any unitary transformation on the joint state of the ions, including transitions from product states to entangled states. Soon after their theoretical work, the essential ingredients of the method were realized experimentally for a single trapped ion , and more recently two ions were driven from a product state to an entangled state with around 90% fidelity using closely related ideas . This is the only efficient creation of entanglement yet observed in any area of physics.
The combination of universal driven unitary evolution, an exponential scaling of the available Hilbert space with system size, and reliable measurement of the resulting state, are the essential ingredients for a future quantum computer. The scaling of system size in this argument is measured by the way the whole physical apparatus becomes larger, and operation of quantum logic gates slower, as more qubits are added to the system. In the case of the linear ion trap, the adding of further qubits (ions) into the trap is straightforward. The growth of the physical apparatus, and the slowing down of the processing, is dominated by the optical and electronic equipment needed to control the ions. This growth and slow-down is a polynomial, not exponential, function of the number of qubits, for small numbers of qubits. Technical problems will place a limit, as yet unknown, on the highest number of qubits which might be feasible, but below this limit the ion trap has the good scaling properties which allow the essential principles of quantum computing to be experimentally realised. Hence, laser-cooled ion trap experiments have a dual interest both for studying new avenues in atomic physics (e.g. interferometry with entangled particles) and for understanding quantum information processing in a specific realisable system.
In this work we examine, both theoretically and experimentally, the speed with which general processing operations can be driven in the ion trap system. We are concerned with the intrinsic limitations imposed by the physics of the system, such as the tightness of the trap and the presence of a rich energy level structure in the vibrational modes. The initial proposal of Cirac and Zoller showed how to drive quantum operations in the limit of small Rabi frequency (slow operations), and initial experiments have been in this limit. Our aim here is to clarify the trade-off between precision and speed of the quantum logic gates, and to derive the optimal way to operate the processor. This optimum is a compromise between speed-associated problems such as off-resonant excitation of unwanted transitions, and the basic decoherence rate due to environmental coupling which will limit the performance when the gates are too slow. Our main observation is to confirm the statement in , that the maximum gate rate obtained by the Cirac Zoller method scales as the geometric mean of the trap vibrational frequency and the recoil frequency associated with the multi-ion string. This is in contrast to statements made elsewhere that the maximum gate rate is roughly proportional to the vibrational frequency, though this faster rate may be available if further tricks are adopted, such as excitation in the node of a laser standing wave . We extend the discussion in by making a quantitative statement, equation (20), of the relation between speed and precision of the gate.
We also calculate the maximum gate rate for the method proposed by Monroe et al. , where a controlled-not operation is achieved without the need for an additional transition and laser frequency, by driving multiple Rabi cycles on the carrier, for well-chosen values of the Lamb-Dicke parameter.
We experimentally study the Cirac-Zoller method, using a single trapped Calcium ion cooled to the ground state of the motion in one dimension. The experiments are akin to those in where the methods were first demonstrated for one ion, except that we examine the regime as yet unexplored, where the gate rate is above the recoil frequency, so that the calculations carried out in the small Rabi frequency limit are no longer appropriate.
The paper is set out as follows. Section II presents the theory of the logic gates, and calculations of their fidelity as a function of the main parameters, primarily the Rabi frequency of the atom–light interaction, and the trap tightness. The calculations involve numerical solution of the Schrödinger equation for the system. We note the need to tune the laser light to the resonant frequency correctly, that is, taking into account the light shift (a.c. Stark shift) of the levels. We present results in the first instance considering only one mode of vibration of the ion string, and then in section II G we briefly discuss the influence of the other modes. Section III presents our experiments. We cool a single trapped calcium ion to the ground state of the motion in one dimension, and then observe Rabi flopping on the first blue motional sideband of the narrow 729 nm transition. By driving the sideband in the regime where the Rabi frequency is of the order of the trap vibrational frequency, we observe the expected trade-off between speed and precision of the operations, caused by off-resonant excitation of unwanted motional states.
## II Theory of switching rate
### A Preliminaries
Consider a string of two or more ions in a linear trap. The trap is strongly confining along $`x`$ and $`y`$ directions, and less strongly confining along the $`z`$ axis. We will assume motion in $`x`$ and $`y`$ directions is unexcited, and consider one normal mode of vibration of the ion string along $`z`$, writing its angular frequency $`\omega _z`$. The other normal modes along $`z`$ will be assumed to be unexcited throughout. (We will consider the element of approximation introduced by this assumption at the end.) For brevity, we will refer to this vibrational degree of freedom as “the normal mode” or “the vibration”. It has an evenly spaced ladder of energy levels $`E_n=(n+1/2)\mathrm{}\omega _z`$. Typically in an experiment one would choose the mode of interest to be the second or third, i.e. $`\omega _z=\sqrt{3}\omega _{z,\mathrm{cm}}`$ or $`\sqrt{29/5}\omega _{z,\mathrm{cm}}`$, where $`\omega _{z,\mathrm{cm}}`$ is the frequency of the centre of mass mode. For modes other than the centre of mass, the Lamb-Dicke parameter describing the ion–light coupling will vary from one ion to another , but this presents no problem as long as it is taken into account when choosing laser pulse intensities and/or durations. We will assume that the laser light is directed onto one ion at a time.
The Lamb-Dicke parameter is $`\eta =\eta _1/\sqrt{N}`$ where $`N`$ is the number of ions, and $`\eta _1=\sqrt{E_R/\mathrm{}\omega _z}`$. $`E_R=(r\mathrm{}k_z)^2/(2M)`$ is the recoil energy for a single ion initially at rest undergoing a $`\pi `$ pulse interaction with the laser field; $`k_z`$ is the $`z`$ component of the wavevector $`𝐤`$ of the laser field, $`M`$ is the mass of the ion, and $`r=1`$ for a single-photon transition, $`r=2`$ for a Raman transition assuming the geometry $`k_{1z}=k_{2z}`$ for the two Raman beams. Example recoil frequencies $`E_R/h`$ are given in table (I).
We will assume throughout that $`\eta <1`$. This is the regime in which the ion trap processor is used in practice, both because it facilitates the initial cooling to the ground state of motion, and because the processor then runs faster, as we will discuss.
The Cirac-Zoller method adopts the internal state of each ion as a qubit (we restrict attention to two internal states and thus one qubit per ion). Gates between qubits are obtained via excitation of the common vibrational mode, and measurement of the state of one or more qubits is by observing fluorescence. The method uses the fact that arbitrary single-qubit rotations, combined with any one 2-qubit gate such as controlled-not ($`{}_{}{}^{C}X`$) or controlled-rotation ($`{}_{}{}^{C}Z`$), between arbitrary pairs of qubits, form a universal set . That is, any unitary transformation can be decomposed into elements from the set. In the trapped ion system, a single-qubit rotation is a transition in the internal state of a single ion, driven by a laser pulse. A $`{}_{}{}^{C}X`$ between the internal states of any pair of ions $`A`$ and $`B`$ is achieved by a swap ($`S`$) operation between the internal state of $`A`$ and the vibration, followed by $`{}_{}{}^{C}X`$ between the vibration and the internal state of ion $`B`$ (with vibration as control, ion internal state as target) followed by $`S`$ again on ion $`A`$. Therefore the only operations that will concern us are single-qubit rotations, and $`S`$ and $`{}_{}{}^{C}X`$ between a single ion and the vibration. All of these are achieved by excitation of a chosen transition. The unwanted off-resonant excitation of other transitions are the main subject of this paper.
The strength of the ion-laser interaction is parametrized by the Rabi frequency, given by
$`\mathrm{\Omega }_{nm}`$ $`=`$ $`n\left|\mathrm{exp}\left(i\eta (\widehat{a}^{}+\widehat{a})\right)\right|m\mathrm{\Omega }_{\mathrm{free}}`$ (1)
$``$ $`C_{nm}\mathrm{\Omega }`$ (2)
where $`\mathrm{\Omega }_{\mathrm{free}}`$ is the Rabi frequency for a free ion, the states $`|n`$ are vibrational energy eigenstates, $`\mathrm{\Omega }=\mathrm{exp}(\eta ^2/2)\mathrm{\Omega }_{\mathrm{free}}`$ and the factor $`C_{nm}`$ is given by the following equation (3). Values of $`C_{nm}`$ are listed in table (II) for the low-lying vibrational levels.
$$C_{nm}=\sqrt{m!n!}(i\eta )^{|fm|}\underset{j=0}{\overset{\mathrm{min}(m,n)}{}}\frac{(1)^j\eta ^{2j}}{j!(j+|nm|)!(\mathrm{min}(m,n)j)!}.$$
(3)
The single-qubit rotations can be much faster than the two-qubit gates, because the energy separation $`\mathrm{}\omega _0`$ of the internal energy levels is much larger than that of the vibrational levels, and by choosing a laser beam direction $`k_z=0`$ (single photon transition) or $`k_{1z}=k_{2z}`$ (Raman transition), the Lamb-Dicke parameter can be made to vanish during single-qubit operations, which means these operations do not couple to the vibrational state ($`C_{nm}`$ becomes $`\delta _{nm}`$). Therefore, for $`\mathrm{\Delta }n=0`$, $`\mathrm{\Omega }`$ can be large compared to $`\omega _z`$ without causing off-resonant excitation of $`\mathrm{\Delta }n0`$ transitions. Therefore, the speed of the ion trap processor is limited by the $`S`$ and $`{}_{}{}^{C}X`$ gates.
The $`S`$ gate is achieved by a $`\pi `$ pulse on either the first red vibrational sideband, that is, at frequency $`\omega _0\omega _z`$, or the first blue vibrational sideband at frequency $`\omega _0+\omega _z`$. The duration of the $`S`$ gate is therefore
$$T_S=\frac{\pi }{\eta \mathrm{\Omega }}.$$
(4)
where the gate becomes exact in the limit $`\mathrm{\Omega }0`$. The choice of red or blue sideband can be dictated by experimental convenience. The two are equivalent in terms of their quantum computational effect, since the logical operations produced can be made identical simply by relabelling the states (i.e. changing which internal state of the ion is called 0, and which is called 1). We will treat the red sideband throughout our theoretical discussion, but the results will apply equally to blue sideband excitation. In fact, we used a blue sideband in the experiments described in section III.
We will consider two methods for the $`{}_{}{}^{C}X`$ gate. The first is that described in the original proposal of Cirac and Zoller, where $`{}_{}{}^{C}X`$ is obtained by single-bit rotations (Hadamard gates) combined with $`{}_{}{}^{C}Z`$, and $`{}_{}{}^{C}Z`$ is obtained from a $`2\pi `$ pulse on the first red sideband (or blue, depending on the relative positions of the levels) of an auxiliary transition in the ion, i.e. laser frequency $`\omega _{\mathrm{aux}}\omega _z`$, gate duration
$$T_{C1}=\frac{2\pi }{\eta \mathrm{\Omega }}.$$
(5)
The other method is that of Monroe et al. , who proposed using a $`2m\pi `$ pulse on the carrier (frequency $`\omega _0`$) where $`m`$ is an integer and $`\eta ^2=1/(2m)`$. This is especially useful if an auxiliary transition is not available. When the vibrational state is $`|n=0`$, this drives an integer number of Rabi oscillations of the internal state, while if the vibrational state is $`|n=1`$, this drives a half-integer number of Rabi oscillations, since $`C_{11}=1\eta ^2=(2m1)/(2m)`$. The method requires $`\eta `$ to be fixed in the experiment to one of the special values $`1/\sqrt{2m}`$, which is easily done in practice by adjusting the trap confinement and/or laser beam direction. The duration of a $`{}_{}{}^{C}X`$ gate by the Monroe method is
$$T_{C2}=\frac{2m\pi }{\mathrm{\Omega }}=\frac{\pi }{\eta ^2\mathrm{\Omega }}.$$
(6)
Since the ion trap is operated with $`\eta <1`$, it might be thought that $`T_{\mathrm{C2}}`$ is necessarily greater than $`T_{\mathrm{C1}}`$. In fact, this is not necessarily the case, since these gate times are limited by the maximum allowable Rabi frequency $`\mathrm{\Omega }`$, and this can depend on the type of gate.
Sorensen and Molmer have proposed a general method to implement gates such as controlled-not using bichromatic laser fields, achieving good fidelity even when the vibrational degrees of freedom are not in their ground state. The speed limitations of this method have recently been re-examined , so we will not discuss it here, other than to say the method appears to be useful.
### B Solution of time-dependent Schrödinger equation
Our aim is to find the maximum switching speed of the processor. It is seen from equations (4), (5) and (6) that this is determined by the maximum allowable Rabi frequency during the gates which operate on the vibrational state. The Rabi frequency cannot be arbitrarily large, because in the limit $`\mathrm{\Omega }\omega _z`$, the evolution would be independent of the vibrational state.
In order to find the maximum Rabi frequency, we need to solve the time-dependent Schrödinger equation for the system, without making the approximation of small $`\mathrm{\Omega }`$. We assume a two-level ion. This means we will not explicitly examine the type of $`{}_{}{}^{C}X`$ gate which uses an auxiliary level (equation (5)), but in any case this operation is closely related to the $`S`$ operation which we will examine, so results for the rate of $`S`$ will apply to this type of $`{}_{}{}^{C}X`$ apart from the factor 2. We will study $`S`$ and the Monroe $`{}_{}{}^{C}X`$ gate.
The Hamiltonian for the ion that is being illuminated by the laser during a given gate is
$$H=H_0+H_I$$
(7)
where $`H_0`$ is diagonal, with diagonal elements given by $`\mathrm{}(0,\omega _z,2\omega _z,\mathrm{},\omega _0,\omega _0+\omega _z,\omega _0+2\omega _z,\mathrm{})`$, and the interaction Hamiltonian
$$H_I=\mathrm{}\frac{\mathrm{\Omega }}{2}\left(\begin{array}{cc}\mathrm{𝟎}& \widehat{C}e^{i\omega _Lt}\\ \widehat{C}^{}e^{i\omega _Lt}& \mathrm{𝟎}\end{array}\right)$$
(8)
where $`\widehat{C}`$ is the matrix having elements $`C_{nm}`$. Note that the only approximation so far is to ignore the other vibrational modes.
We first adopt a frame rotating with the laser frequency: $`|\stackrel{~}{\psi }(t)=U|\psi (t),`$ where $`U`$ is a diagonal unitary matrix, with diagonal elements $`\mathrm{exp}(i\{1,1,1,\mathrm{},1,1,1,\mathrm{}\}\omega _Lt/2)`$. The states $`|\stackrel{~}{\psi }(t)`$ satisfy the Schrödinger equation $`i\mathrm{}(d/dt)|\stackrel{~}{\psi }=\stackrel{~}{H}|\stackrel{~}{\psi }`$, where $`\stackrel{~}{H}i\mathrm{}(dU/dt)U^{}+UHU^{}`$ is now a time-independent Hamiltonian. We can therefore write the solution to the Schrödinger equation
$$|\stackrel{~}{\psi }(t)=e^{i\stackrel{~}{H}t/\mathrm{}}|\stackrel{~}{\psi }(0)V^{}e^{iV\stackrel{~}{H}V^{}t/\mathrm{}}V|\stackrel{~}{\psi }(0)$$
(9)
where $`V`$ is the matrix of eigenvectors of $`\stackrel{~}{H}`$, so that $`V\stackrel{~}{H}V^{}`$ is diagonal.
The quantity of interest, from the point of view of quantum information processing in the trap, is the final state expressed as a superposition of computational basis states. If at $`t=0`$ the computational basis states are $`|\stackrel{~}{u}(0)=|u(0)`$, then at other times they are $`\mathrm{exp}(i\stackrel{~}{H}_0t/\mathrm{})|\stackrel{~}{u}(0)`$ since then the coefficients $`\stackrel{~}{u}(t)|\stackrel{~}{\psi }(t)`$ do not evolve in the absence of gate operations. Therefore we would like to calculate
$$u(0)\left|e^{i\left(\stackrel{~}{H}\stackrel{~}{H}_0\right)t/\mathrm{}}\right|\psi (0).$$
(10)
Define $`P\mathrm{exp}(i(\stackrel{~}{H}\stackrel{~}{H}_0)t/\mathrm{})`$, and let $`G`$ be the precise unitary operator for the intended gate. Then the degree to which the laser pulse produces the intended gate is given by the overlap between the final state and one which would be obtained from $`G`$. To indicate the degree of imperfection of the laser pulse operation, we therefore calculate
$$f_{\mathrm{min}}=\underset{\psi }{\mathrm{min}}\left|\psi \left|G^{}P\right|\psi \right|^2$$
(11)
and define the imprecision to be $`ϵ=(1f_{\mathrm{min}})^{1/2}`$.
Let us comment on whether or not $`ϵ0`$ represents imperfection in the system. We are assuming no imperfection in the sense of unknown dynamics (e.g. laser intensity and frequency noise). Therefore as long as we can calculate numerically the effect of the laser pulse operation, we have accurate knowledge of the expected state of the system: the fact that the laser pulse gate differs from any particular “ideal” gate is not a source of any imprecision at all. However, the intention is to use the ion trap as a quantum processor, to do quantum calculations which we lack the computing power to simulate classically. Putting together many laser pulses, we are therefore assuming we cannot predict the effect on the whole quantum process of having $`ϵ0`$ in all the operations. Therefore $`ϵ`$ must be regarded as imprecision in the device, which must be minimized.
The question of driving complicated evolution in a predictable manner is central to other techniques such as NMR spectroscopy, and there exist methods to combine pulses causing different types of driven rotation in order to undo the effect of certain terms in the Hamiltonian. Equivalent methods can almost certainly be found for the ion trap system, but in any case operations close to ideal ones will remain the best starting point for any such method.
### C Performance without correction for light shift
In this section, we will examine the operation of the processor in the regime $`\eta <1`$, in the case that the light shifts (a.c. Stark shifts) are not taken into account when choosing the laser pulse frequency, phase, and duration.
First consider the $`S`$ gate: a $`\pi `$ pulse on the first red sideband. We set $`\omega _L=\omega _0\omega _z`$ in (II B), and examine the eigenvalues of $`\stackrel{~}{H}`$. The separations between the eigenvalues enable one to deduce the Rabi flopping frequencies and the light shifts in the system. The Rabi flopping frequency on the $`|n=0|n=1`$ transition is $`\eta \mathrm{\Omega }`$, as expected, which confirms that the gate time is $`T_S=\pi /(\eta \mathrm{\Omega })`$ as in equation (4).
The frequency of the transition is shifted by
$$\mathrm{\Delta }\omega =\frac{1}{2}\left(1+\frac{\eta ^2}{2}\right)\left(\frac{\mathrm{\Omega }}{\omega _z}\right)^2\omega _z+\frac{1}{8}\left(\frac{\mathrm{\Omega }}{\omega _z}\right)^4\omega _z+\mathrm{}.$$
(12)
This is primarily the light shift caused by the presence of $`\mathrm{\Delta }n=0`$ transitions, which are off-resonant by $`\omega _z`$.
As was noted by Wineland et al. , the primary error in the gate operation is caused by the light shift $`\mathrm{\Delta }\omega `$, which can be significant compared to the Rabi flopping frequency $`\eta \mathrm{\Omega }`$. To understand its effect, we now model the transition of interest, that is $`|g,n=1|e,n=0`$, as a two-level system, driven by a $`\pi `$ pulse off-resonant by $`\mathrm{\Delta }\omega `$. The only term in the propagator of first order in $`\mathrm{\Delta }\omega `$ is $`i(\mathrm{\Delta }\omega /\alpha )\mathrm{sin}(\alpha t/2)`$ where $`\alpha ^2=\eta ^2\mathrm{\Omega }^2+\mathrm{\Delta }\omega ^2`$. Putting $`\eta \mathrm{\Omega }t=\pi `$, we obtain
$$ϵ\left|\frac{i\mathrm{\Delta }\omega }{\alpha }\right|\frac{\mathrm{\Omega }}{2\eta \omega _z}$$
(13)
Therefore, to attain a given gate precision $`ϵ`$ the Rabi frequency must satisfy $`\mathrm{\Omega }2ϵ\eta \omega _z`$, which gives
$`{\displaystyle \frac{1}{T_S}}={\displaystyle \frac{\eta \mathrm{\Omega }}{\pi }}`$ $``$ $`{\displaystyle \frac{2ϵ}{\pi }}\eta ^2\omega _z`$ (14)
$`=`$ $`4ϵ{\displaystyle \frac{E_R}{Nh}}.`$ (15)
The result thus has the simple form that the processor speed for swap operations is limited by the $`N`$-ion recoil frequency.
Next, consider the $`{}_{}{}^{C}X`$gate using the special Lamb-Dicke parameter method of Monroe et al., that is, a $`2m\pi `$ pulse on the carrier with $`2m=1/\eta ^2`$. Substituting $`\omega _L=\omega _0`$ into (II B), and examining the eigenvalues of $`\stackrel{~}{H}`$, we find the transition frequencies are shifted by
$$\mathrm{\Delta }\omega \frac{(\eta \mathrm{\Omega })^2}{\omega _z}.$$
(16)
Since $`\mathrm{\Delta }\omega \mathrm{\Omega }`$, the light shifts for this gate are much less significant. Modelling each driven transition as a two-level system, using the same method described above for the $`S`$ gate, we find that the effect of $`\mathrm{\Delta }\omega `$ is $`ϵ\eta ^2\mathrm{\Omega }/\omega _z`$. This is small compared to the effect of off-resonant excitation of $`\mathrm{\Delta }n=\pm 1`$ transitions. These are driven in the regime where the detuning is large compared to their Rabi frequency, leading to a propagator with terms of order $`\sqrt{2}\eta \mathrm{\Omega }/\omega _z`$, so we find
$$ϵ\frac{\sqrt{2}\eta \mathrm{\Omega }}{\omega _z}.$$
(17)
Setting $`\mathrm{\Omega }ϵ\omega _z/\sqrt{2}\eta `$, we find
$$\frac{1}{T_{C2}}\sqrt{2}ϵ\sqrt{\frac{E_R}{Nh}\frac{\omega _z}{2\pi }},\mathrm{\Omega }\omega _z.$$
(18)
It might appear from this that the gate speed can increase without limit by using a tight trap, but this is a false impression because the expression is only valid in the regime $`\mathrm{\Omega }\omega _z`$, which means it is only valid for a gate rate small compared to $`\eta ^2\omega _z/\pi E_R/(Nh)`$.
### D Gates with correction for light shift: Red sideband
The light shift is readily corrected for by tuning the laser frequency to accurate resonance with the light-shifted transition frequency . We will now describe the correct laser pulse frequency, phase, and duration, and derive an expression for the imprecision caused by the primary remaining unwanted effect, which is off-resonant excitation. We find that the recoil frequency still strongly influences the processor speed.
We have studied this problem by numerically evaluating expressions (10) and (11). Only vibrational levels with $`n<4`$ were included, thus yielding an $`8\times 8`$ propagator matrix. When further vibrational levels were included in the calculation, the results were not significantly affected.
We find that the best implementation of a swap gate is obtained when the laser frequency is tuned to the light-shifted frequency of the first red sideband: $`\omega _L=\omega _0\omega _z+\mathrm{\Delta }\omega `$ where $`\mathrm{\Delta }\omega `$ is given by equation (12). The gate time is still as in equation (4). With no further adjustments, this produces a propagator close to $`U_\varphi S`$ where $`S`$ is a perfect swap operator, and $`U_\varphi `$ is diagonal with diagonal elements $`\mathrm{exp}(i\{1,1,1,\mathrm{},1,1,1,\mathrm{}\}\mathrm{\Delta }\omega T_S/2)`$. Therefore, to produce the intended gate, the system must be further evolved by an application of $`U_\varphi ^{}`$. This is equivalent to applying subsequent laser pulses with a correspondingly adjusted phase. Therefore the swap gate is completed by changing the phase of the oscillator used in the experimental apparatus to keep the laser in step with the atom’s internal resonance (at frequency $`\omega _0`$), by an amount $`\mathrm{\Delta }\varphi =\mathrm{\Delta }\omega T_S`$.
This adjustment permits the swap gate rate to exceed the recoil frequency. The term in the system propagator $`P`$ which now gives the dominant contribution to the imprecision $`ϵ`$ is off-resonant excitation of $`\mathrm{\Delta }n=0`$ transitions. We expect the amplitude for this process to be of order $`\mathrm{\Omega }/\sqrt{\mathrm{\Omega }^2+\omega _z^2}`$. Our numerical calculations confirm this, giving
$$ϵ\frac{\mathrm{\Omega }}{\sqrt{2(\mathrm{\Omega }^2+\omega _z^2)}}\frac{\mathrm{\Omega }}{\sqrt{2}\omega _z}$$
(19)
for the remaining imprecision in the adjusted swap operation. Therefore the gate rate is now limited by
$$\frac{1}{T_S}2\sqrt{2}ϵ\sqrt{\frac{E_R}{Nh}\frac{\omega _z}{2\pi }},$$
(20)
which is $`2\sqrt{2}ϵ`$ times the geometric mean of the $`N`$-ion recoil frequency and the vibrational frequency of the chosen normal mode.
If an auxilliary transition is available, the above argument will hold for the $`{}_{}{}^{C}Z`$ gate implemented by a $`2\pi `$ pulse to the first vibrational sideband of the auxilliary transition, the only difference is that the gate rate at given $`ϵ`$ is half that for the swap gate ($`T_{C1}=2T_S`$).
### E Carrier
The Monroe method $`{}_{}{}^{C}X`$ gate is executed by a laser pulse at the carrier frequency. The speed limit given by equation (18) cannot be exceeded, since it is caused by off-resonant excitation that cannot be avoided. However, we are interested in the behaviour outside the region of validity of (18), that is, when $`\mathrm{\Omega }\omega _z`$.
When $`\eta `$ is small but $`\mathrm{\Omega }`$ is not, the precision of the operation need not be limited by off-resonant excitation of $`\mathrm{\Delta }n=\pm 1`$ transitions, since their rate is low. The source of imprecision instead comes from the fact that the dynamics of the system can no longer be pictured as separate Rabi flopping on the $`n=00`$ transition and the $`n=11`$ transition. Rather, the interaction Hamiltonian drives the set of vibrational levels $`n=0,1,2,\mathrm{}\mathrm{\Omega }/\omega _z`$ as a single entity. The result is that the evolution from an initial state $`|n=0`$ is almost the same as the evolution from an initial state $`|n=1`$. Therefore the Rabi flopping must be driven for a longer time before a $`{}_{}{}^{C}X`$ gate is obtained, and the gate speed ceases to increase with $`\mathrm{\Omega }`$. The rate obtained at $`\mathrm{\Omega }\omega _z`$ was found to be the maximum. Recalling the expression (6) for $`T_{C2}`$, it is seen that this maximum is approximately the recoil frequency.
In the region $`\mathrm{\Omega }\omega _z`$, the gate can be optimized simply by adjusting the pulse duration. The correct pulse length is not precisely $`\pi /(\eta ^2\mathrm{\Omega })`$, but differs from this by an amount of order $`\pi /(\eta \mathrm{\Omega })`$. We have not found any simple expression for this adjustment, we surmise that this is because it is a complicated function of all the light shifts in the multi-level system. To run a processor in practice, the correction to the gate time can either be calculated numerically, or measured experimentally. By this adjustment, we were able to extend the region of validity of expression (18) up to $`\mathrm{\Omega }0.5\omega _z`$. Operating at this limit, the gate rate is always equal to the $`N`$-ion recoil frequency, independent of the trap tightness, but the gate becomes more precise as the trap gets tighter ($`ϵ\eta /\sqrt{2}`$, expression (17)).
One way to use the carrier excitation method to achieve faster $`{}_{}{}^{C}X`$ gates, is to use higher vibrational levels such as $`|n=2`$ (in conjunction with $`|n=0`$). However, this would mean that the $`S`$ gate would have to drive the second red sideband, making it slower (the expected limitation at small $`\eta `$ being off-resonant excitation of $`\mathrm{\Delta }n=0`$ transitions). Therefore only small gains, if any, are available by this route.
### F Sensitivity to laser intensity fluctuations
Although our main purpose is to study limitations imposed by the unavoidable properties of the system, we comment here on the sensitivity to one source of technical noise, namely laser intensity fluctuations, since these partially limit our experiments described in section III. The laser drives Rabi flopping at frequency $`\eta \mathrm{\Omega }`$ for the Cirac-Zoller swap gate, and at frequency $`\mathrm{\Omega }`$ for the Monroe gate. If, owing to intensity fluctuations, $`\mathrm{\Omega }`$ is imprecise by $`\mathrm{\Delta }\mathrm{\Omega }`$ then the action of either gate will be imprecise by $`ϵT_S\eta \mathrm{\Delta }\mathrm{\Omega }/\pi `$ and $`ϵT_{C2}\mathrm{\Delta }\mathrm{\Omega }/\pi `$ respectively. We therefore find $`\mathrm{\Delta }\mathrm{\Omega }/\mathrm{\Omega }ϵ`$ for the swap gate (also for the closely related $`{}_{}{}^{C}Z`$ gate), whereas $`\mathrm{\Delta }\mathrm{\Omega }/\mathrm{\Omega }\eta ^2ϵ`$ for the Monroe gate.
### G Allowance for all the normal modes
The only approximation we made was to neglect all but one of the normal modes of oscillation of the ion string. When we relax this assumption, the main features of the more general problem can be understood in terms of the Rabi frequencies and the normal mode frequencies. To keep the notation clear, we will describe a case of three normal modes, which we will label $`x`$, $`y`$ and $`z`$, but the discussion readily generalises to more normal modes. In our experiments described in section III there is a just a single ion, so the modes of interest are indeed vibration along the principle axes of the trap, but the modes here labelled $`x`$, $`y`$ and $`z`$ could equally refer to different normal modes of the $`z`$-oscillation of a three-ion string, in the case where the two other directions of oscillation are frozen out owing to the linear geometry (that is, they have much higher frequency, and are prepared in the ground state). Equation (2) is replaced by
$$\mathrm{\Omega }(n_x,m_x)(n_y,m_y)(n_z,m_z)=C_{n_xm_x}(\eta _x)C_{n_ym_y}(\eta _y)C_{n_zm_z}(\eta _z)\mathrm{\Omega }$$
(21)
in an obvious notation, where we now redefine $`\mathrm{\Omega }\mathrm{exp}((\eta _x^2+\eta _y^2+\eta _z^2)/2)\mathrm{\Omega }_{\mathrm{free}}`$. There are two regimes to consider. In the case that all the modes are prepared very close to the ground state, the main effect is that the light shift is no longer that given by equation (12), but one given by a sum of light shifts related to all the possible transitions. This is easily compensated by a suitable adjustment to the laser frequency. The limit on precision therefore remains that due to off-resonant transitions. For excitation on the red sideband ($`n_z=0,m_z=1`$), in addition to carrier transitions excited off-resonant by $`\omega _z`$ with Rabi frequency $`\mathrm{\Omega }`$, there are also sideband transitions excited off-resonant by $`\omega _x\omega _z`$ and $`\omega _y\omega _z`$, with Rabi frequencies $`\eta _x\mathrm{\Omega }`$ and $`\eta _y\mathrm{\Omega }`$ respectively. Since the latter contribute at a higher order in $`\eta `$ compared to the carrier term, they are negligible unless two modes are close in frequency. For a linear ion string the separation of the lowest mode frequencies is large enough that the additional terms make little contribution. Therefore our previous discussion, including equations (19) and (20), remains valid. In the experiments described in section III, on the other hand, $`\omega _y\omega _z0.025\omega _z`$ so unwanted excitation of the $`y`$ sideband has to be taken into account.
For the carrier excitation, in addition to the $`z`$ sidebands excited off-resonant by $`\omega _z`$ with Rabi frequency $`\eta _z\mathrm{\Omega }`$, there are also sideband transitions excited off-resonant by $`\omega _x`$ and $`\omega _y`$ with Rabi frequencies $`(1\eta _z^2)\eta _x\mathrm{\Omega }`$ and $`(1\eta _z^2)\eta _y\mathrm{\Omega }`$ respectively. These contribute at the same order in $`\eta `$ as the $`z`$ sidebands, and result in a change to equation (18) by a numerical factor of order unity.
Our experiments were carried out with only one mode cooled to the ground state. The other two modes had roughly thermal distributions with mean vibrational quantum number of order 10. The result of this is to blur all the Rabi frequencies and light shifts, since each time a given $`z`$ transition is excited, the Rabi frequency depends on the $`x`$ and $`y`$ vibrational quantum numbers:
$$\mathrm{\Omega }(n_x,n_x)(n_y,n_y)(n_z,m_z)\left(1n_x\eta _x^2\right)\left(1n_y\eta _y^2\right)C_{n_zm_z}\mathrm{\Omega }.$$
(22)
We are able to fit our experimental results by calculating the evolution for each value of $`n_x,n_y`$, and then averaging over a thermal population of the $`x`$ and $`y`$ vibrations.
### H Summary
In conclusion, we find that the swap gate and the Cirac-Zoller type of controlled-phase gate can be made faster, at a given level of precision, by making the trap tighter, but the speed increase is only in proportion to $`\omega _z^{1/2}`$. This is in agreement with earlier work , but here we have added a precise quantitative statement of the trade-off between speed and precision. Some previous studies have adopted $`T_S\pi /\omega _z`$ for the purpose of making rough estimates, but we find the error thus introduced is significant; for example, it leads to an overestimate of the switching rate by two orders of magnitude in . By exciting the ion in the node of laser standing wave, it is possible to avoid (or greatly reduce) the off-resonant carrier excitation , therefore $`T_S\pi /\omega _z`$ may be available, but this involves a significant extra complication of the experiment which may render it unfeasible in practice .
The “magic Lamb-Dicke parameter” method of Monroe et al can not be speeded up indefinitely, it has a natural speed limit given by the recoil frequency.
There are advantages in having the ions spaced by a micron or more, in order to allow laser addressing of one ion at a time, and/or resolving the fluorescence from different ions with appropriate imaging optics . If we set a limit $`s`$ for the smallest permissible seperation of the closest ions in the string during processing, we obtain .
$$\omega _{z,\mathrm{cm}}^2<\frac{8e^2}{4\pi ϵ_0Ms^3N^{1.71}}$$
(23)
Putting this in equation (20), and adopting the breathing mode $`\omega _z=\sqrt{3}\omega _{z,\mathrm{cm}}`$, we obtain for the fastest swap gate rate, at given precision and ion spacing,
$$\frac{1}{T_S}2.5\left(\frac{e^2}{4\pi ϵ_0}\right)^{1/4}ϵ\left(\frac{E_R}{h}\right)^{1/2}\frac{1}{M^{1/4}s^{3/4}N^{0.93}}.$$
(24)
If we approximate the $`N`$ dependence as $`N^1`$, and set $`s=10\lambda `$ where $`\lambda `$ is the laser wavelength, then we arrive at a “gate time per ion” which depends only on the choice of ion and transition, at given $`ϵ`$. This quantity, $`T_S/N`$, is given in table III for some example transitions, for the case of 99% gate fidelity (i.e. $`ϵ=0.1`$).
## III Experiments
In this section, we describe our experimental investigations, carried out in conditions where the Rabi frequency $`\mathrm{\Omega }_{\mathrm{sideband}}\eta \mathrm{\Omega }_{\mathrm{carrier}}`$ of the sideband excitation is significantly above the recoil frequency. We thus are able to investigate the regime where the speed of execution of the gates determines the fidelity of operation.
We use a single trapped <sup>40</sup>Ca<sup>+</sup> ion to demonstrate the principle of quantum information processing. The electronic $`|S_{1/2}`$, $`m=1/2`$ ground state and the metastable $`|D_{5/2}`$, $`m=5/2`$ state (1s lifetime) are used to implement one qubit. We apply a 4 Gauss bias magnetic field to lift the degeneracy of sublevels in the ground and excited state manifolds. The qubit can be coherently manipulated by laser light at 729 nm (see table I). Two basic operations are demonstrated. First there is the single qubit rotation which only affects the individual internal electronic state, and secondly we perform the swap operation which entangles the electronic state of the ion and its vibrational state. We did not investigate the Monroe et al. $`{}_{}{}^{C}X`$ gate because our experiments involved small Lamb Dicke parameters of order $`0.045`$. In this regime the Monroe method is $`500`$ times more sensitive than the Cirac-Zoller method to laser intensity noise and thermal populations in spectator modes, and we found it to be unworkable in these experiments.
We store the ion in a conventional spherical-quadrupole Paul trap with trap frequencies $`\omega _{x,y,z}=2\pi (4.0,1.925,1.850)`$MHz. To acquire each experimental data point, at given values of the parameters, we run a series of 100 cycles. A single experimental cycle is made up of 5 stages, which are (i) Doppler cooling, (ii) sideband cooling, (iii) coherent driven evolution, (iv) observation of fluorescence, (v) deshelving. We then record the fraction $`P_D`$ of the 100 cycles in which fluorescence was observed in the penultimate stage.
In more detail, the 5 stages of a cycle are as follows. Doppler cooling is performed on the $`|S_{1/2}`$ to $`|P_{1/2}`$ 397 nm electric dipole transition. The electronic ground state is then prepared in a pure $`|S_{1/2},m=1/2`$ state by optical pumping. The ion is cooled to the vibrational ground state of the $`\omega _z`$ mode by applying sideband cooling on the $`|S_{1/2},m=1/2|D_{5/2},m=5/2`$ electric quadrupole transition at 729 nm. A full description of the trap and the cooling procedure is given in Ref. . The wavevector $`𝐤`$ of the 729 nm radiation is nearly perpendicular to the $`y`$ direction, and has an angle of 40 and 50 to the $`z`$ and $`x`$ directions. The corresponding values of $`\eta _{x,y,z}^{729nm}`$ are (0.04, 0.01, 0.045). In the qubit operation step, the $`|S_{1/2},m=1/2|D_{5/2},m=5/2`$ transition is excited with a laser pulse of well controlled frequency, intensity, and timing. Then, by monitoring fluorescence at 397 nm under laser excitation, we detect whether a transition to the non-fluorescing state D<sub>5/2</sub> occurred. The scheme allows one to discriminate between the internal states of the ion with an efficiency close to 100% . In our experiment,the discrimination effiency is approximately 99.8$`\%`$, limited by the lifetime of the metastable state .
Finally, the ion is again repumped (deshelved) from the D<sub>5/2</sub> level to the electronic ground state via the P<sub>3/2</sub> level. The fraction $`P_D`$ of cycles in which fluorescence is observed indicates the population of the $`D_{5/2}`$ level after the coherent qubit operation step.
### A Carrier excitation
As discussed in section II A, the single qubit rotation can be driven fast and without loss in contrast, even in the range where $`\mathrm{\Omega }`$ exceeds the trap frequencies $`\omega _{x,y,z}`$, as long as the Lamb-Dicke parameters are sufficiently small. To reach this limit in the case of a single-photon transition under travelling wave excitation, the wavevector $`𝐤`$ of the exciting light field should be perpendicular to the chosen vibration direction, and to the direction having next smallest vibrational frequency. The remaining vibrational frequency must be large compared to $`\mathrm{\Omega }`$. In our case the laser is nearly perpendicular to one of the weak confinement directions, but not the other, so there is some unwanted excitation of the motion. A limit to execution speed is set by the available laser power to drive the qubit transition. However, Rabi frequencies as high as a few MHz can be achieved even in the case of dipole forbidden quadrupole transitions.
Rabi oscillations on the $`\mathrm{\Delta }n=0`$ carrier transition at a frequency $`\mathrm{\Omega }_{\mathrm{carrier}}=2\pi \mathrm{\hspace{0.33em}1090}`$ kHz are shown in Fig. 1a. For this, 100mW of light are focussed into a waist of 30$`\mu `$m at the position of the ion. We observe a contrast of more than 95%, the contrast being mainly limited by the thermally distributed axial vibration mode ($`\omega _x=2\pi \mathrm{\hspace{0.17em}4.0}`$MHz), which causes a spread in Rabi frequency owing to the $`n`$-dependence as indicated in equation (22). We find that our results are consistent with an average over a thermal phonon distribution in the x-oscillator having $`n=12(2)`$.
### B Sideband excitation
The second building block is the sideband excitation, involving a $`\mathrm{\Delta }n0`$ transition, here performed on the $`z`$ vibration. For convenience, we adopted the first blue rather than red sideband. The equivalence of blue and red sidebands for these studies was mentioned in section II A. The laser frequency was red-detuned from the blue sideband resonance (measured at low laser power) to compensate for the a.c. Stark effect. The accuracy of the $`\mathrm{\Delta }n=+1`$ operation is therefore limited by off-resonant excitations on the $`\mathrm{\Delta }n=0`$ transition, as discussed in section II D, see also Fig. 2. We indeed observe off-resonant carrier excitation, visible as fast oscillations with a frequency near $`\omega _z`$ on top of the $`\mathrm{\Omega }_{sideband}=2\pi \mathrm{\hspace{0.17em}48}`$ kHz Rabi oscillation, see Fig. 3. Note that this frequency is far beyond the recoil frequency, and thus the contrast shrinks to 75% for a $`\pi `$ pulse.
If the blue sideband $`\omega _z`$ is excited at a Rabi frequency of $`\mathrm{\Omega }_{\mathrm{sideband}}=2\pi \mathrm{\hspace{0.17em}7.4}`$kHz (Fig. 1b), the off-resonant carrier excitation is no longer detectable. A contrast of 95% for a $`\pi `$ pulse and 85% for a $`2\pi `$ pulse is measured, while equation (19) predicts $`ϵ=0.063`$ and hence a fidelity $`99.6`$% for the $`\pi `$ pulse. Under low-power excitation, we do not observe any light shift.
The experimental data can be fitted by numerically solving the Schrödinger equation in the truncated basis of the 4-level system consisting of $`|S_{1/2},n=0`$, $`|D_{5/2},n=0`$, $`|S_{1/2},n=1`$, and $`|D_{5/2},n=1`$ states (see Fig. 2), and averaging over a thermal distribution of population in the other vibrational modes. Three parameters in the calculation were independently measured: the Rabi frequency on the carrier $`\mathrm{\Omega }=2\pi 1090`$kHz, the Lamb Dicke parameter $`\eta _z`$, and the trap frequency $`\omega _z=2\pi \mathrm{\hspace{0.17em}1850}`$ kHz. The detuning of the laser field $`\delta `$ is varied for optimum contrast and we find excellent agreement between the data and a numerical simulation for a value of $`\delta _{\mathrm{th}}/2\pi `$= -375 kHz (Fig. 3). The discrepancy of 125 kHz between $`\delta _{\mathrm{th}}`$ and the experimentally determined value $`\delta _{\mathrm{exp}}/2\pi =250`$ kHz is probably caused by light shifts due to the $`D_{5/2}P_{3/2},S_{1/2}P_{1/2}`$, and $`S_{1/2}P_{3/2}`$ transitions. The sum of the calculated light shifts from the dipole transitions is approximately 70 kHz.
A third contribution to the overall light shift is expected from the other radial ($`y`$) mode at 1.925 MHz ($`\delta \omega /2\pi `$=75 kHz). Each level of the y-vibrational ladder is light-shifted, resulting in a change in the $`(n_y,n_y)(1_z,0_z)`$ z-sideband resonance frequency by approximately $`(\eta _y\mathrm{\Omega }_{carrier})^2(n_y+1)/(2(\omega _y\omega _z))=(n_y+1)(2\pi )\mathrm{\hspace{0.33em}0.8}`$ kHz. Assuming a mean vibrational quantum number of 25, a 20 kHz spread in shifts is expected. To model this situation we solved the Schrödinger equation for different $`y`$-oscillator occupations and averaged the results over a thermal distribution (see Fig. 4). The loss of contrast at $`\mathrm{\Omega }t=\pi `$ is dominated by this thermal averaging, rather than laser intensity and magnetic field fluctuations.
To compare the experimental findings with an optimum $`\pi `$ pulse, for a given pulse length, we investigated numerically the parameter space and found that an optimum $`\pi `$ pulse (92% contrast of Rabi $`|S_{1/2}|D_{5/2}`$ oscillation) could be achieved for a detuning of $`\delta _{\mathrm{th}}/2\pi =355`$ kHz. We conclude that the detuning chosen in the experiment was off by 20 kHz. Note, that the electron shelving detection method only measures the internal state occupation, not the vibrational one. Hence, the measured contrast does not directly yield the fidelity. Equation (19) predicts $`ϵ=0.41`$ for this case, and hence a fidelity of 83%. The 20 kHz deviation from the ideal detuning leads to a 5% loss in contrast and the thermally distributed $`y`$-vibration mode with $`n=25`$ further reduces the contrast to 75%. The corresponding fidelity was near 64%. For the details of the $`\pi `$ pulse dynamics, see Fig.4.
To conclude, we have presented a theoretical discussion of two types of quantum gate which couple the internal ion state with its motion, and an experimental study of one of these. Our calculations lead to quantitative statements of the precision of the operations, taking into account the complete Hamiltonian including all the vibrational states and the off-resonant coupling terms. These statements are given in equations (17), (18) for the Monroe et al “magic Lamb Dicke parameter” gate, and in equations (19), (20) for the Cirac-Zoller swap gate. If we set a limit of a fixed number of laser wavelengths for the spacing of the ions in a linear trap, then we arrive at a gate time per ion, for the swap gate, which depends only on the precision and the choice of ion and transition. This time is given in equation (24) and table III.
In our experimental studies we have demonstrated a contrast of 75% for a $`1/T_S=100`$kHz swap operation, and 95% for a $`1/T_S=14`$kHz swap operation on a single trapped ion. These gate rates are respectively 21 and 3 times the relevant recoil frequency. In both cases the contrast could be significantly improved by cooling a further vibrational mode.
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# Untitled Document
Impurity Scattering in Carbon Nanotubes with Superconducting Pair Correlations
Kikuo Harigaya
Electrotechnical Laboratory, Umezono 1-1-4, Tsukuba 305-8568, Japan
Abstract. Effects of the superconducting pair potential on the impurity scattering processes in metallic carbon nanotubes are studied theoretically. The backward scattering of electrons vanishes in the normal state. In the presence of the superconducting pair correlations, the backward scatterings of electron- and hole-like quasiparticles vanish, too. The impurity gives rise to backward scatterings of holes for incident electrons, and it also induces backward scatterings of electrons for incident holes. Negative and positive currents induced by such the scatterings between electrons and holes cancel each other. Therefore, the nonmagnetic impurity does not hinder the supercurrent in the regions where the superconducting proximity effects occur, and the carbon nanotube is a good conductor for Cooper pairs. Relations with experiments are discussed.
INTRODUCTION
Recent investigations (1,2) show that the superconducting proximity effect occurs when the carbon nanotubes contact with conventional superconducting metals and wires. The superconducting energy gap appears in the tunneling density of states below the critical temperature $`T_\mathrm{c}`$. On the other hand, the recent theories discuss the nature of the exceptionally ballistic conduction (3) and the absence of backward scattering (4) in metallic carbon nanotubes with impurity potentials at the normal states.
In this paper, we study the effects of the superconducting pair potential on the impurity scattering processes in metallic carbon nanotubes, using the continuum $`𝒌𝒑`$ model for the electronic states. We find the absence of backward scatterings of electron- and hole-like quasiparticles in the presence of superconducting proximity effects, and the nonmagnetic impurity does not hinder the supercurrent in the regions where the superconducting proximity effects occur. Therefore, the carbon nanotube is a good conductor for Cooper pairs as well as in the normal state. This finding is interesting in view of the recent experimental progress of the superconducting proximity effects of carbon nanotubes (1,2).
IMPURITY SCATTERING
IN NORMAL NANOTUBES
We will study the metallic carbon nanotubes with the superconducting pair potential. The model is as follows:
$$H=H_{\mathrm{tube}}+H_{\mathrm{pair}},$$
(1)
$`H_{\mathrm{tube}}`$ is the electronic states of the carbon nanotubes, and the model based on the $`𝒌𝒑`$ approximation (4,5) represents electronic systems on the continuum medium. The second term $`H_{\mathrm{pair}}`$ is the pair potential term owing to the proximity effect.
The hamiltonian of a graphite plane by the $`𝒌𝒑`$ approximation (4,5) in the secondly quantized representation has the following form:
$$H_{\mathrm{tube}}=\underset{𝒌,\sigma }{}\mathrm{\Psi }_{𝒌,\sigma }^{}E_𝒌\mathrm{\Psi }_{𝒌,\sigma },$$
(2)
where $`E_𝒌`$ is an energy matrix:
$$E_𝒌=\left(\begin{array}{cccc}0& \gamma \left(k_xik_y\right)& 0& 0\\ \gamma \left(k_x+ik_y\right)& 0& 0& 0\\ 0& 0& 0& \gamma \left(k_x+ik_y\right)\\ 0& 0& \gamma \left(k_xik_y\right)& 0\end{array}\right),$$
(3)
$`𝒌=(k_x,k_y)`$, and $`\mathrm{\Psi }_{𝒌,\sigma }`$ is an annihilation operator with four components: $`\mathrm{\Psi }_{𝒌,\sigma }^{}=(\psi _{𝒌,\sigma }^{(1)},\psi _{𝒌,\sigma }^{(2)},\psi _{𝒌,\sigma }^{(3)},\psi _{𝒌,\sigma }^{(4)})`$. Here, the fist and second elements indicate an electron at the A and B sublattice points around the Fermi point $`K`$ of the graphite, respectively. The third and fourth elements are an electron at the A and B sublattices around the Fermi point $`K^{}`$. The quantity $`\gamma `$ is defined as $`\gamma (\sqrt{3}/2)a\gamma _0`$, where $`a`$ is the bond length of the graphite plane and $`\gamma _0`$ ($``$ 2.7 eV) is the resonance integral between neighboring carbon atoms. When the above matrix is diagonalized, we obtain the dispersion relation $`E_\pm =\pm \gamma \sqrt{k_x^2+\kappa _{\nu \varphi }^2(n)}`$, where $`k_x`$ is parallel with the axis of the nanotube, $`\kappa _{\nu \varphi }(n)=(2\pi /L)(n+\varphi \nu /3)`$, $`L`$ is the circumference length of the nanotube, $`n`$ ($`=0`$, $`\pm 1`$, $`\pm 2`$, …) is the index of bands, $`\varphi `$ is the magnetic flux in units of the flux quantum, and $`\nu `$ ($`=0`$, 1, or 2) specifies the boundary condition in the $`y`$-direction. The metallic and semiconducting nanotubes are characterized by $`\nu =0`$ and $`\nu =1`$ (or 2), respectively. Hereafter, we consider the case $`\varphi =0`$ and the metallic nanotubes $`\nu =0`$.
The second term in Eq. (1) is the pair potential:
$$H_{\mathrm{pair}}=\mathrm{\Delta }\underset{𝒌}{}(\psi _{𝒌,}^{\left(1\right)}\psi _{𝒌,}^{\left(1\right)}+\psi _{𝒌,}^{\left(2\right)}\psi _{𝒌,}^{\left(2\right)}+\psi _{𝒌,}^{\left(3\right)}\psi _{𝒌,}^{\left(3\right)}+\psi _{𝒌,}^{\left(4\right)}\psi _{𝒌,}^{\left(4\right)}+\mathrm{h}.\mathrm{c}.)$$
(4)
where $`\mathrm{\Delta }`$ is the strength of the superconducting pair correlation of an $`s`$-wave pairing. We assume that the spatial extent of the regions where the proximity effect occurs is as long as the superconducting coherence length.
Now, we consider the impurity scattering in the normal metallic nanotubes. We take into account of the single impurity potential located at the point $`𝒓_0`$:
$$H_{\mathrm{imp}}=I\underset{𝒌,𝒑,\sigma }{}\mathrm{e}^{i\left(𝒌𝒑\right)𝒓_0}\mathrm{\Psi }_{𝒌,\sigma }^{}\mathrm{\Psi }_{𝒑,\sigma },$$
(5)
where $`I`$ is the impurity strength.
The scattering $`t`$-matrix at the $`K`$ point is
$$t_K=I\left[1I\frac{2}{N_s}\underset{𝒌}{}G_K(𝒌,\omega )\right]^1,$$
(6)
where $`G_K`$ is a propagator of a $`\pi `$-electron around the Fermi point $`K`$. The discussion about the $`t`$-matrix at the $`K^{}`$ point is qualitatively the same, so we only look at the $`t`$-matrix at the $`K`$ point. The sum for $`𝒌=(k,0)`$, which takes account of the band index $`n=0`$ only, is replaced with an integral:
$$\frac{2}{N_s}\underset{𝒌}{}G_K(𝒌,\omega )=\rho 𝑑ϵ\frac{1}{\omega ^2ϵ^2}\left(\begin{array}{cc}\omega & ϵ\\ ϵ& \omega \end{array}\right)\rho \pi i\mathrm{sgn}\omega \left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),$$
(7)
where $`\rho =a/2\pi L\gamma _0`$ is the density of states at the Fermi energy. Therefore, we obtain
$$t_K=\frac{I}{1+I\rho \pi i\mathrm{sgn}\omega }\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right).$$
(8)
The transformation into the energy-diagonal representation where the branches with $`E=\pm \gamma |k|`$ are diagonal has the same form of $`t_K`$.
The scattering matrix $`t_K`$ in the energy-diagonal representation is diagonal, and the off-diagonal matrix elements vanish. This means that only the scattering processes from $`k`$ to $`k`$ and from $`k`$ to $`k`$ are effective. The scatterings from $`k`$ to $`k`$ and from $`k`$ to $`k`$ are cancelled. Such the absence of the backward scattering has been discussed recently (4).
IMPURITY SCATTERING WITH SUPERCONDUCTING PAIR POTENTIAL
We consider the single impurity scattering when the superconducting pair potential is present. In the Nambu representation, the scattering $`t`$-matrix at the $`K`$ point is
$$\stackrel{~}{t}_K=\stackrel{~}{I}\left[1\frac{2}{N_s}\underset{𝒌}{}\stackrel{~}{G}_K(𝒌,\omega )\stackrel{~}{I}\right]^1,$$
(9)
where $`\stackrel{~}{G}_K`$ is the Nambu representation of $`G_K`$ and
$$\stackrel{~}{I}=I\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& 1\end{array}\right).$$
(10)
The sign of the scattering potential for holes is reversed from that for electrons, so the minus sign appears at the third and fourth diagonal matrix elements.
The sum over $`𝒌`$ is performed as in the previous section, and we obtain the scattering $`t`$-matrix (with the same form in the energy-diagonal representation):
$$\stackrel{~}{t}_K=\frac{I}{1+\left(I\rho \pi \right)^2}\left(\begin{array}{cccc}1+\alpha \omega & 0& \alpha \mathrm{\Delta }& 0\\ 0& 1+\alpha \omega & 0& \alpha \mathrm{\Delta }\\ \alpha \mathrm{\Delta }& 0& 1+\alpha \omega & 0\\ 0& \alpha \mathrm{\Delta }& 0& 1+\alpha \omega \end{array}\right)$$
(11)
where $`\alpha =I\rho \pi i/\sqrt{\omega ^2\mathrm{\Delta }^2}`$.
Hence, we find that the off-diagonal matrix elements become zero in the diagonal $`2\times 2`$ submatrix. This implies that the backward scatterings of electron-line and hole-like quasiparticles vanish in the presence of the proximity effects, too. Off-diagonal $`2\times 2`$ submatrix has the diagonal matrix elements whose magnitudes are proportional to $`\mathrm{\Delta }`$. The finite correlation gives rise to backward scatterings of the hole of the wavenumber $`k`$ when the electron with $`k`$ is incident. The back scatterings of the electrons with the wavenumber $`k`$ occur for the incident holes with $`k`$, too. Negative and positive currents induced by such the two scattering processes cancel each other. Therefore, the nonmagnetic impurity does not hinder the supercurrent in the regions where the superconducting proximity effects occur. This effect is interesting in view of the recent experimental progress of the superconducting proximity effects (1,2).
SUMMARY
We have investigated the effects of the superconducting pair potential on the impurity scattering processes in metallic carbon nanotubes. The backward scattering of electrons vanishes in the normal state. In the presence of the superconducting pair correlations, the backward scatterings of electron- and hole-like quasiparticles vanish, too. The impurity gives rise to backward scatterings of holes for incident electrons, and it also induces backward scatterings of electrons for incident holes. Negative and positive currents induced by such the scatterings between electrons and holes cancel each other. Therefore, the carbon nanotube is a good conductor for the Cooper pairs coming from the proximity effects.
REFERENCES
1. A. Y. Kasumov et al, Science 284, 1508 (1999).
2. A. F. Morpurgo, J. Kong, C. M. Marcus, and H. Dai, Science 286, 263 (1999).
3. C. T. White and T. N. Todorov, Nature 393, 240 (1998).
4. T. Ando and T. Nakanishi, J. Phys. Soc. Jpn. 67, 1704 (1998).
5. H. Ajiki and T. Ando, J. Phys. Soc. Jpn. 62, 1255 (1993).
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# A Cluster of Black Holes at the Galactic Center
## 1 Introduction
The measurement of proper motions and radial velocities of stars within the central parsec of the Galaxy has led to the conclusion that a black hole of mass $`3.0\pm 0.3\times 10^6M_{}`$ is present in the center (Eckart & Genzel 1997; Genzel et al. 2000). There is also increasing evidence that massive black holes are found in the centers of other galaxies (Richstone et al. 1998).
The central region of the Galaxy is also peculiar because the relaxation time among stars can be shorter than the age of the Galaxy, owing to the high density. The process of relaxation leads to a stellar cusp, which has a density profile $`\rho r^{7/4}`$ when all the stars have the same mass (Bahcall & Wolf 1976). Several interesting physical processes take place among the stars in this cusp: stars can come close enough to physically collide with each other, and they can also come sufficiently close to the black hole to be tidally disrupted or swallowed (e.g., Frank & Rees 1976; Lightman & Shapiro 1977; Quinlan, Hernquist, & Sigurdsson 1995; Sigurdsson & Rees 1997).
One of the consequences of the relaxation is that the most massive objects will sink to the center of the stellar cusp. Among an old stellar population, the most massive objects should be black holes formed in the final core collapse of massive stars. We assume in this paper that most massive stars with $`M30M_{}`$ produce black holes, most with a mass of $`7M_{}`$ (Bailyn et al. 1998). The high mass of these black holes implies that their dynamical friction time to move to the center of the Galaxy is shorter than a Hubble time over a much larger volume than the one where ordinary stars have a short relaxation time. We will find in §2 that this should lead to the formation of a cluster of stellar black holes around the central supermassive black hole (hereafter “Sgr A\*”), and that other stars are ejected from the region occupied by this cluster. In §3 we discuss the rate at which the black holes in this cluster are captured by Sgr A\*, and we find that most of the black holes should still be present in the cluster. Several observable consequences of the presence of this black hole cluster are discussed in §4.
## 2 Cluster Formation
The deprojected light profile in the inner kpc of the Galaxy scales as $`r^{1.8}`$, while the predicted profile around a massive black hole $`(r1\mathrm{pc})`$ scales as $`r^{7/4}`$ (Bahcall & Wolf 1976). For simplicity, we therefore adopt a density profile $`\rho (r)`$
$$\rho (r)r^{7/4}.$$
(1)
From the model fit of Genzel et al. (2000) to the velocity dispersion data, we find that the total mass within $`r_0=1.8\mathrm{pc}`$ is $`2M_{\mathrm{cbh}}`$ (see their Fig. 17), where
$$M_{\mathrm{cbh}}=3\times 10^6M_{}$$
(2)
is the mass of Sgr A\*. Hence, the total distributed mass inside 1.8 pc is $`M_{\mathrm{cbh}}`$, and the density profile is
$$\rho _{}(r)=\frac{5}{16\pi }\frac{M_{\mathrm{cbh}}}{r_0^3}\left(\frac{r}{r_0}\right)^{7/4},r_01.8\mathrm{pc}.$$
(3)
We assume that this density profile is entirely composed of stars, brown dwarfs, and stellar remnants. To calculate the mass fraction of black holes, $`\eta _{\mathrm{bh}}`$, we use the following initial mass function. For the range $`M>1M_{}`$, we adopt a Salpeter law $`dN/dmm^\alpha `$ with $`\alpha =2.35`$. For $`0.7M_{}<m<1M_{}`$, we adopt $`\alpha =2`$ from Zoccali et al. (2000). For $`0.05M_{}<m<0.7M_{}`$, we adopt $`\alpha =1.65`$ by correcting the result of Zoccali et al. (2000) for binaries according to the adjustment of Gould, Bahcall, & Flynn (1997) and by extending the power law beyond the last observed point at $`0.15M_{}`$. We cut off the mass function at $`m0.05M_{}`$ in accordance with the prelimary indications from microlensing (Han & Gould 1976). We assume that all progenitors with masses 1–8 $`M_{}`$ have become $`0.6M_{}`$ white dwarfs, those with masses 8–30 $`M_{}`$ have become $`1.4M_{}`$ neutron stars, and those with masses 30–100 $`M_{}`$ have become $`7M_{}`$ black holes. We then find
$$\eta _{\mathrm{bh}}=1.6\%,m=0.23M_{},$$
(4)
where $`m`$ is the mean mass of the population.
Once this population is formed, the black holes will sink toward the center on the dynamical friction timescale (Binney & Tremaine 1987)
$$t_{\mathrm{df}}^1=\mathrm{ln}\mathrm{\Lambda }\frac{4\pi G\rho Gm_{\mathrm{bh}}}{v^3}_0^v𝑑u\mathrm{\hspace{0.17em}4}\pi u^2f(u),$$
(5)
where $`m_{\mathrm{bh}}=7M_{}`$ is the mass of the black hole, $`v`$ is its velocity, $`f(u)`$ is the velocity distribution of the ambient stars, and the integral gives the fraction of ambient stars with speeds below $`v`$. For a Gaussian velocity distribution with dispersion $`\sigma `$, and a typical black hole speed $`v^2=3\sigma ^2`$, the value of the integral is $`0.54`$. For the Keplerian part of the potential ($`r<r_0`$), $`\mathrm{ln}\mathrm{\Lambda }=\mathrm{ln}(M_{\mathrm{cbh}}/m_{\mathrm{bh}})=13`$. For $`r>r_0`$, $`\mathrm{ln}\mathrm{\Lambda }`$ rises slightly but we ignore this in the interest of simplicity. We evaluate $`t_{\mathrm{df},0}`$, the fiducial dynamical friction time at $`r_0`$,
$$t_{\mathrm{df},0}=1.4\mathrm{Gyr}$$
(6)
and note its scaling in the two regimes
$$t_{\mathrm{df}}=t_{\mathrm{df},0}\left(\frac{r}{r_0}\right)^{1/4}(r<r_0),t_{\mathrm{df}}=t_{\mathrm{df},0}\left(\frac{r}{r_0}\right)^{17/8}(r>r_0).$$
(7)
Hence, after a time $`t`$, all the black holes that were originally within a radius $`r`$ will collect in a cluster near the center, where $`r`$ is given by,
$$\frac{r}{r_0}=\left(\frac{t}{4t_{\mathrm{df},0}}\right)^4(t<4t_{\mathrm{df},0}),\frac{r}{r_0}=\left[\frac{17}{7}\left(\frac{t}{t_{\mathrm{df},0}}4\right)+1\right]^{8/17}(t>4t_{\mathrm{df},0}).$$
(8)
This implies an infall rate of black holes
$$\frac{dN_{\mathrm{bh}}}{dt}=\frac{5\eta _{\mathrm{bh}}M_{\mathrm{cbh}}}{4t_{\mathrm{df},0}m_{\mathrm{bh}}}\left(\frac{t}{4t_{\mathrm{df},0}}\right)^4(t<4t_{\mathrm{df},0}),$$
$$\frac{dN_{\mathrm{bh}}}{dt}=\frac{10\eta _{\mathrm{bh}}M_{\mathrm{cbh}}}{7t_{\mathrm{df},0}m_{\mathrm{bh}}}\left[\frac{17}{7}\left(\frac{t}{t_{\mathrm{df},0}}4\right)+1\right]^{7/17}(t>4t_{\mathrm{df},0}).$$
(9)
If we assume that the bulge formed at a time $`t_{\mathrm{bulge}}10`$Gyr, then from equation (8), all the black holes within a radius $`r_{\mathrm{df}}=5\mathrm{pc}`$ will have migrated to the center by now. The cumulative total and current rate of precipitation are therefore
$$N_{\mathrm{bh}}2.4\times 10^4,\frac{dN_{\mathrm{bh}}}{dt}2.9\mathrm{Myr}^1.$$
(10)
In other words, provided that our assumption of the fraction of massive stars that form black holes in their final core collapse is correct, we must conclude that a large number of stellar black holes, with a total mass of $`5\%`$ of the Sgr A\* mass, have migrated to the center and, unless they have subsequently been captured by Sgr A\*, they should have formed a cluster of black holes in the center of the stellar cusp.
As the black holes precipitate, they start dominating the total density in some central region at some point, and then the low-mass stars are expelled from this region over a relaxation time. Assuming that most of the energy is lost from the cluster by direct capture of black holes near the center, the black holes should also relax to a density profile proportional to $`r^{7/4}`$, for which the outward energy flow is constant. As this energy flow is transmitted to the low-mass stars outside the black hole cluster, the cluster will need to expand and push out the low-mass stars. To derive the relative density of the black hole profile compared to the empirically normalized stellar profile, we invoke the steady state energy-flow condition between two species of stars A and B of mass $`m_A`$ and $`m_B`$, which dominate the total density $`\rho _A`$ and $`\rho _B`$ at radii $`r_A`$ and $`r_B`$, respectively. The total energy at radius $`r`$ is proportional to $`\rho \sigma ^2r^3`$, and the relaxation time is proportional to $`\sigma ^3/(\rho m)`$. Therefore, the constant energy flow condition yields
$$\frac{\rho _A\sigma _A^2r_A^3}{\sigma ^3/(\rho _Am_A)}=\frac{\rho _B\sigma _B^2r_B^3}{\sigma ^3/(\rho _Bm_B)},$$
(11)
Making use of the Kepler-potential relation $`\sigma ^2r`$, this implies,
$$\frac{\rho _A}{\rho _B}=\left(\frac{r_A}{r_B}\right)^{7/4}\left(\frac{m_A}{m_B}\right)^{1/2}.$$
(12)
Thus, the mass density of black holes in the central region is below that implied by extrapolating equation (3), by a factor $`(m/m_{\mathrm{bh}})^{1/2}`$, that is, $`\rho _{\mathrm{bh}}(r)=0.18\rho _{}(r)`$. Hence, if all of the black holes precipitated over the lifetime of the galaxy from a radius of $`5\mathrm{p}\mathrm{c}`$ remain in the cluster at present, the cluster should extend over a radius $`r_{\mathrm{bh}}`$,
$$r_{\mathrm{bh}}=\left(\frac{\eta _{\mathrm{bh}}^2m_{\mathrm{bh}}}{m}\right)^{2/5}\mathrm{\hspace{0.17em}5}\mathrm{pc}=0.7\mathrm{pc}.$$
(13)
The timescale to achieve this expansion is the relaxation time in the expanded (lower density) cluster, which is $`3`$ times longer than the dynamical friction timescale because of the lower density.
## 3 Rate of Capture of the Black Holes
In the previous section we found that about 24,000 stellar black holes should have migrated to the center of the stellar cusp around Sgr A\*. We now address the question of the rate at which these black holes will be removed from the cluster by coalescing with Sgr A\*. If this rate is low, most of the black holes should be in the cluster at present. If the rate is high enough, then many fewer black holes will be present, and a balance between the rate at which black holes are precipitating in the cluster by dynamical friction and the rate at which they are being captured should be established.
The dominant process by which black holes will be eliminated is by a random walk into a highly eccentric orbit as their orbits change over the relaxation timescale, from which they can be captured by Sgr A\*. This process was first studied by Frank & Rees (1976). In the case of stars, tidal disruption can eliminate them from the cluster once they come close enough to Sgr A\*; obviously, orbiting black holes will be eliminated only when they are swallowed by Sgr A\*, possibly after having lost orbital energy by emitting gravitational waves.
Before describing in more detail the mechanism by which black holes are captured, we need to discuss the process of orbital diffusion by which black holes will migrate into the eccentric orbits from which they can be captured.
### 3.1 Orbital Diffusion
A black hole can be captured by Sgr A\* from an orbit of any semimajor axis, provided that its peribothron $`q=a(1e)`$ is sufficiently small. This will lead to a distribution of black holes in phase space that is strongly depleted at eccentricities very close to unity, and diffusion of black holes will take place toward orbits of decreasing peribothron. In order to investigate quantitatively this black hole migration, we first evaluate the diffusion tensor in velocity space. We sketch the derivation here and leave the details and the justifications for the various approximations to Appendix A.
The diffusion equation is given by
$$_v𝐣+\frac{f}{t}=0,j_k\underset{l}{}\kappa _{kl}\frac{f(𝐯,𝐫,t)}{v_l}$$
(14)
where $`f`$ is the phase-space density, $`𝐣`$ is the “velocity current density”, and $`\kappa _{kl}`$ is the diffusion tensor. By symmetry, $`\kappa _{kl}=`$ diag($`\kappa _{},\kappa _{},\kappa _r`$), where $`\kappa _{}`$ and $`\kappa _r`$ are the components of $`\kappa `$ perpendicular and parallel to the radial direction. In general, the diffusion tensor depends on the spatial position and the velocity. A useful physical interpretation of the diffusion tensor components is that, over a small interval of time $`\delta t`$, the rms change in the velocity of a black hole in a direction perpendicular to its initial velocity is equal to $`(2\kappa _{}\delta t)^{1/2}`$. The total rms change in any direction is therefore $`(2tr(\kappa )\delta t)^{1/2}`$ (where tr means the trace), and the relaxation time is of order $`v_{\mathrm{esc}}^2/[2tr(\kappa )]`$.
We assume that the unperturbed (by black-hole capture) phase-space density, $`f_0`$, is a function only of the energy, and hence find, for a $`\rho r^\alpha `$ density profile in a Kepler potential,
$$f_0(𝐮,𝐫)=g\left(\frac{u^2}{v_{\mathrm{esc}}^2}\right)h(r),g(x)(1x)^{\alpha 3/2}\mathrm{\Theta }(1x),$$
(15)
$$h(r)\frac{(3\alpha )\alpha !}{8\pi ^2(\alpha 3/2)!(1/2)!}\frac{N_{\mathrm{bh}}}{(2GM_{\mathrm{cbh}}r_{\mathrm{bh}})^{3/2}}\left(\frac{r}{r_{\mathrm{bh}}}\right)^{3/2\alpha },$$
(16)
where $`v_{\mathrm{esc}}(r)`$ is the local escape velocity, $`N_{\mathrm{bh}}`$ is the total number of black holes within a radius $`r_{\mathrm{bh}}`$, $`M_{\mathrm{cbh}}`$ is the central mass, and $`\mathrm{\Theta }`$ is a step function. Thus, for $`\alpha =7/4`$, $`f_0r^{1/4}(1v^2/v_{\mathrm{esc}}^2)^{1/4}`$. In Appendix A, we show that for $`\alpha =3/2`$, the velocity dependences of the parallel and perpendicular components of the diffusion tensor are
$$\kappa _{}(v,r)=\left(1\frac{1}{5}\frac{v^2}{v_{\mathrm{esc}}^2}\right)\kappa _0(r),\kappa _r(v,r)=\left(1\frac{3}{5}\frac{v^2}{v_{\mathrm{esc}}^2}\right)\kappa _0(r),$$
(17)
where $`\kappa _0`$ is the (isotropic) diffusion coefficient at $`v=0`$. We also justify in the Appendix using this velocity dependence as an adequate approximation for the case $`\alpha =7/4`$, for which we find,
$$\kappa _0(r)=\frac{7}{6\sqrt{\pi }}\frac{(3/4)!}{(1/4)!}\frac{N_{\mathrm{bh}}(Gm_{\mathrm{bh}})^2v_{\mathrm{esc}}^2}{(2GM_{\mathrm{cbh}}r_{\mathrm{bh}})^{3/2}}\left(\frac{r}{r_{\mathrm{bh}}}\right)^{1/4}\mathrm{ln}\mathrm{\Lambda }_1,$$
(18)
The quantity $`\mathrm{ln}\mathrm{\Lambda }_1`$ can be slightly smaller than $`\mathrm{ln}\mathrm{\Lambda }=\mathrm{ln}(M_{\mathrm{cbh}}/m_{\mathrm{bh}})`$, as discussed in Appendix A, depending on the scale over which the diffusion takes place, because scatterings with large velocity changes do not contribute to the diffusion rate over a small range of velocities.
### 3.2 Three Regimes of Capture
Before identifying the condition for capture from highly eccentric orbits, it will be useful to compute the core radius, $`r_c`$, where the energy-loss time to gravitational radiation on a circular orbit is equal to the relaxation time. We define the relaxation time for circular orbits as $`t_{r,0}=v_c^2/[2tr(\kappa )]`$, where $`v_c=(GM_{\mathrm{cbh}}/r)^{1/2}`$ is the circular velocity. The energy loss by gravitational radiation for circular orbits is
$$\frac{d\mathrm{ln}E}{dt}=\frac{64}{5}\frac{G^3m_{\mathrm{bh}}M_{\mathrm{cbh}}^2}{c^5r_c^4}.$$
(19)
Equating this to $`t_{r,0}^1`$, we find,
$$r_c=1.57\left(\frac{GM_{\mathrm{cbh}}}{c^2}\right)^{2/3}r_{\mathrm{bh}}^{1/3}\left(\frac{M_{\mathrm{cbh}}}{N_{\mathrm{bh}}m_{\mathrm{bh}}\mathrm{ln}\mathrm{\Lambda }}\right)^{4/15}$$
$$=8.6\mathrm{AU}\left(\frac{M_{\mathrm{cbh}}}{3\times 10^6M_{}}\right)^{2/3}\left(\frac{r_{\mathrm{bh}}}{0.7\mathrm{pc}}\right)^{1/3}\left(\frac{M_{\mathrm{cbh}}}{1.38N_{\mathrm{bh}}m_{\mathrm{bh}}\mathrm{ln}\mathrm{\Lambda }}\right)^{4/15}.$$
(20)
This radius is extremely small compared to the cluster radius $`r_{bh}`$. If black holes were captured by diffusing to orbits with $`ar_c`$, and then losing energy by gravitational waves at low eccentricity, the capture rate would therefore also be extremely small. In reality, as we shall see now, black holes will be captured from a large range of radii, mostly from highly eccentric orbits. We must therefore define a relaxation time scale for these high eccentricity orbits.
To do this, we first calculate the orbit-averaged rate of change of the peribothron $`qa(1e)`$ in a highly eccentric orbit. From angular momentum conservation, we have $`v_{}^2=(1+e)qGM_{\mathrm{bh}}/r^2qv_{\mathrm{esc}}^2/r`$. Hence,
$$P\frac{dq}{dt}=_0^P𝑑t\frac{2\kappa _{}r}{v_e^2}=2^{5/4}\mathrm{ln}\mathrm{\Lambda }_1\frac{53}{33}\left(\frac{m_{\mathrm{bh}}}{M_{\mathrm{cbh}}}\right)^2N_{\mathrm{bh}}\left(\frac{a}{r_{\mathrm{bh}}}\right)^{9/4}r_{\mathrm{bh}},$$
(21)
We then define the eccentric relaxation time $`t_{r,1}(a)=(d\mathrm{ln}q/dt)^1`$. We now use $`t_{r,1}`$ to demonstrate that there are three regimes of capture: 1) for $`a<a_{\mathrm{trans}}`$, capture is dominated by gravitational radiation, 2) for $`a_{\mathrm{trans}}<a<a_{\mathrm{crit}}`$, it is dominated by direct capture, and 3) for $`a>a_{\mathrm{crit}}`$, capture falls off very rapidly and can be ignored.
A particle orbiting around a Schwarzschild black hole with semi-major axis $`a`$ much greater than the Schwarzschild radius will be directly captured by the black hole if its peribothron $`q`$, computed by extrapolating the Newtonian orbit, is less than 4 Schwarzschild radii (e.g., Misner, Thorne, & Wheeler 1973). When $`q=8GM_{\mathrm{cbh}}/c^2`$, the particle is actually brought to 2 Schwarzschild radii by relativistic effects, where the maximum of the effective radial potential is located; the particle has then overcome the angular momentum barrier and can directly fall into the black hole. Therefore, the phase space density should drop to zero below the minimum peribothron
$$q_{\mathrm{min}}=\frac{8GM_{\mathrm{cbh}}}{c^2}0.24\mathrm{AU}.$$
(22)
The black hole can also be captured at larger peribothron if the timescale to change the eccentricity by diffusion is longer than the timescale for losing its orbital energy by gravitational waves. For $`1e1`$, the gravitational radiation decay rate is,
$$\frac{d\mathrm{ln}E}{dt}=\frac{170}{3}(2q)^{7/2}\frac{G^3m_{\mathrm{bh}}M_{\mathrm{cbh}}^2}{c^5a^{1/2}}.$$
(23)
We determine $`q_{\mathrm{min}}`$, the minimum peribothron that avoids capture, by setting $`d\mathrm{ln}E/dt=t_{r,1}^1`$,
$$q_{\mathrm{min}}(a)=0.35\left(\frac{8GM_{\mathrm{cbh}}}{c^2}\right)\left(\frac{M_{\mathrm{cbh}}}{1.73N_{\mathrm{bh}}m_{\mathrm{bh}}\mathrm{ln}\mathrm{\Lambda }_1}\right)^{2/5}\left(\frac{r_{\mathrm{bh}}}{a}\right)^{1/2},$$
(24)
where the factor $`1.73`$ is the ratio $`M_{\mathrm{cbh}}/(N_{\mathrm{bh}}m_{\mathrm{bh}}\mathrm{ln}\mathrm{\Lambda }_1)`$ for the values we use in this paper, and for $`\mathrm{ln}\mathrm{\Lambda }_1=9.7`$ (see Appendix A). Hence, the transition from capture by gravitational radiation to direct capture occurs at $`a_{\mathrm{trans}}=0.35^2r_{\mathrm{bh}}`$ = 17,000 AU.
Above some critical semimajor axis $`a_{\mathrm{crit}}`$, the diffusion of $`q`$ over a single period exceeds $`q_{\mathrm{min}}`$ as given by equation (22) for direct capture. Hence, during each period, $`Pa^{3/2}`$, the black holes are captured with a probability that decreases with semimajor axis as $`(q_{\mathrm{min}}/a)a^1`$, so that the capture rate falls off $`a^{5/2}`$. Since the mass within radius $`r`$ increases as $`r^{3\alpha }`$, captures from orbits with $`aa_{\mathrm{crit}}`$ produce a negligible loss of black holes. We evaluate $`a_{\mathrm{crit}}`$ by setting $`t_{r,1}=P`$, and find
$$a_{\mathrm{crit}}=0.41\mathrm{pc}\left(\frac{N_{\mathrm{bh}}}{2.4\times 10^4}\right)^{4/9}\left(\frac{M_{\mathrm{cbh}}}{3\times 10^6M_{}}\right)^{4/3}\left(\frac{m_{\mathrm{bh}}}{7M_{}}\right)^{8/9}\left(\frac{r_{\mathrm{bh}}}{0.7\mathrm{pc}}\right)^{5/9}.$$
(25)
### 3.3 Capture Rate From The Loss Cylinder
While $`\kappa `$ is a function of both $`𝐯`$ and $`r`$, we will solve the diffusion equation at fixed $`r`$, and temporarily assume that $`\kappa `$ is independent of $`𝐯`$. We will introduce variation in $`\kappa `$ only when we evaluate the loss rate. This is a very good approximation, as we discuss in Appendix A. We focus first on the case of $`a_{\mathrm{trans}}<a<a_{\mathrm{crit}}`$, for which capture is typically direct as described by equation (22). Since $`q_{\mathrm{min}}`$ is independent of $`a`$, the captured orbits at fixed $`r`$ are characterized by a cylinder in velocity space,
$$\frac{v_{}}{v_{\mathrm{esc}}}\sqrt{\frac{q_{\mathrm{min}}}{r}}=2\frac{v_{\mathrm{esc}}}{c},\frac{|v_r|}{v_{\mathrm{esc}}}1,$$
(26)
where $`v_r`$ and $`𝐯_{}`$ are the radial and perpendicular components of the velocity. Note that this structure in velocity space is definitely a “loss cylinder” and not a “loss cone” as it is often described in the literature. Making use of yet another good approximation described in Appendix A, we solve the steady-state diffusion equation (14), at fixed $`v_r`$:
$$f(𝐯_{};v_r)=f_0(v_r)\left[1\frac{\mathrm{ln}(v_{}^2/(v_{\mathrm{esc}}^2v_r^2))}{\mathrm{ln}(q_{\mathrm{min}}/r)}\right].$$
(27)
Hence, the capture rate per unit volume, $`dC/dV`$, is given by
$$\frac{dC}{dV}=_{\mathrm{cyl}\mathrm{bnd}}𝐣𝑑𝐀=8\pi _0^{v_{\mathrm{esc}}}𝑑v_r\frac{\kappa _{}(v_r)f_0(v_r)}{\mathrm{ln}(r/q_{\mathrm{min}})},$$
(28)
where $`d𝐀`$ is the area element on the boundary of the capture cylinder (“cyl-bnd”). We evaluate this using equations (15) and (17),
$$\frac{dC}{dV}=\frac{35}{16\pi ^{3/2}}\frac{(3/4)!}{(1/4)!}\frac{\mathrm{ln}\mathrm{\Lambda }_1}{\mathrm{ln}(r/q_{\mathrm{min}})}\frac{N_{\mathrm{bh}}^2(Gm_{\mathrm{bh}})^2}{(2GM_{\mathrm{cbh}}r_{\mathrm{bh}})^{3/2}}\left(\frac{r}{r_{\mathrm{bh}}}\right)^2r_{\mathrm{bh}}^3$$
$$=\frac{6\kappa _0}{v_{\mathrm{esc}}^2}\frac{n(r)}{\mathrm{ln}(r/q_{\mathrm{min}})},$$
(29)
where $`n(r)=5N_{\mathrm{bh}}/(16\pi r_{\mathrm{bh}}^3)(r/r_{\mathrm{bh}})^{7/4}`$ is the number density of black holes at radius $`r`$, and $`6\kappa _0/v_{\mathrm{esc}}^2`$ is of order the relaxation time. Integrating over volume, the total capture rate within some maximum radius $`r_{\mathrm{max}}`$ (determined below) is
$$C=k\frac{\mathrm{ln}\mathrm{\Lambda }_1}{\mathrm{ln}(r_{\mathrm{max}}/q_{\mathrm{min}})}N_{\mathrm{bh}}^2\left(\frac{m_{\mathrm{bh}}}{M_{\mathrm{cbh}}}\right)^2\frac{r_{\mathrm{max}}}{r_{\mathrm{bh}}}\frac{2\pi }{P_{\mathrm{bh}}}$$
(30)
where $`k=(2\pi )^{1/2}(35/8)(3/4)!/(1/4)!1.76`$, $`P_{\mathrm{bh}}`$ is the orbital period at $`a=r_{\mathrm{bh}}`$, and where we have made the evaluation treating $`\mathrm{ln}(r/q_{\mathrm{min}})\mathrm{ln}(r_{\mathrm{max}}/q_{\mathrm{min}})`$ and $`\mathrm{ln}\mathrm{\Lambda }_1`$ as constants. Notice that the capture rate increases as $`N_{\mathrm{bh}}^{4/5}`$ as black holes are added to the cluster, because $`r_{\mathrm{bh}}N_{\mathrm{bh}}^{4/5}`$, and $`P_{\mathrm{bh}}N_{\mathrm{bh}}^{6/5}`$.
As described already by Frank & Rees (1976), the result we have reached in equations (29) and (30) shows that the capture rate of the black holes is essentially proportional to the relaxation time, and the fact that $`q_{\mathrm{min}}/r1`$ increases the time required for the black hole to find a capture orbit only logarithmically. The reason is that, over a relaxation time, an orbiting black hole will totally change its eccentricity not only due to some large deflection in a close encounter, but also due to many small deflections that will change the eccentricity by very small amounts, allowing the black hole to effectively sweep over all possible eccentricities and find the very narrow range of eccentricity where it can be captured. However, this is no longer true for $`a>a_{\mathrm{crit}}`$: when $`q`$ is brought below $`q_{\mathrm{min}}`$ by the random deflections, the black hole will most likely miss being captured unless it happens to be at peribothron.
The maximum radius $`r_{\mathrm{max}}`$ of integration of the capture rate in equation (30) must therefore be of order $`a_{\mathrm{crit}}`$. For a highly eccentric orbit, the time-averaged radius is $`r=(3/2)a`$. We therefore adopt $`r_{\mathrm{max}}=(3/2)a_{\mathrm{crit}}`$. Since capture is dominated by black holes near $`r_{\mathrm{max}}`$, we evaluate $`\mathrm{ln}(r/q_{\mathrm{min}})`$ there and find,
$$\mathrm{ln}(r_{\mathrm{max}}/q_{\mathrm{min}})=\mathrm{ln}\frac{3c^2a_{\mathrm{crit}}}{16GM_{\mathrm{cbh}}}13.2.$$
(31)
At the same time, the term $`\mathrm{ln}\mathrm{\Lambda }_1`$ can be approximated as (see Appendix A)
$$\mathrm{ln}\mathrm{\Lambda }_1\mathrm{ln}\mathrm{\Lambda }\frac{1}{4}\mathrm{ln}\frac{c^2r_{\mathrm{max}}}{8GM_{\mathrm{cbh}}}9.7$$
(32)
Hence, the ratio of logarithms in equation (30) is $`\mathrm{ln}\mathrm{\Lambda }_1/\mathrm{ln}(r_{\mathrm{max}}/q_{\mathrm{min}})0.73`$. We are finally able to evaluate the capture rate explicitly,
$$C\frac{N_{\mathrm{bh}}}{30\mathrm{Gyr}}.$$
(33)
Since this timescale is much longer than a Hubble time, most of the 24,000 black holes that have entered the cluster are still in it and have not been captured. Hence, the actual radius of the black hole cluster is close to our initial estimate given by equation (13).
Note that we have everywhere used the direct capture formula (22) to calculate $`q_{\mathrm{min}}`$ rather than the gravitational radiation formula (24), which applies at $`r<a_{\mathrm{trans}}17,000`$AU. Recall, however, that $`q_{\mathrm{min}}`$ only enters into the logarithm term. The capture rate from the inner part of the cluster is therefore higher than we have assumed, but not dramatically. Notice that since $`q_{\mathrm{min}}`$ is a function of $`a`$ in equation (24), the loss structure in velocity space is not a cylinder as in the case of direct capture \[see eq. (26)\]. Rather, this structure is fatter near $`v_r0`$ and narrower near $`v_r\pm v_{\mathrm{esc}}`$.
## 4 Observable Consequences
The large number of black holes that move to the center by dynamical friction have several observable effects which we now discuss.
The most important effect is that the stars that formerly resided in the region that is now occupied by the cluster of black holes are ejected into orbits at larger radius. Therefore, the old stellar population of mass $`m`$ should have a very large core radius, given by $`r_{cs}=r_{\mathrm{bh}}(m_{\mathrm{bh}}/m)^{2/7}12\mathrm{pc}`$. This core radius ($`40^{\prime \prime }`$) is much larger than the value expected from stellar collisions alone, which would only produce a core at the radius $`0.03\mathrm{pc}`$ where the orbital velocities are comparable to the escape velocities from the stellar surfaces. The most straightforward test of our model is therefore to measure the distribution of low-mass stars and find out if this very large core indeed exists.
In fact, the bright ($`K<15`$) stars in the inner galaxy exhibit a power-law profile ($`\alpha 7/4`$) all the way to $`r0.1\mathrm{pc}`$ where a core sets in (Genzel et al. 2000; Schmitt 1995). However, this is not necessarily inconsistent with a core in the old population: most of the observed stars could well be young. Indeed, Genzel et al. (2000) have found that the stars they observed very near to the Galactic center are young supergiants, which must have formed in a recent ($`10`$Myr) burst of star formation. If such bursts of star formation have occurred repeatedly, then we would expect that the inner region $`r<r_{\mathrm{bh}}0.7\mathrm{pc}`$ should contain the remains of the bursts from the last $`t_r(r_{\mathrm{bh}})2`$Gyr, but not earlier. Blum, Sellgren, & DePoy (1996a) found that the luminosity function in the inner $`2^{}`$ (5 pc) extends to several magnitudes brighter than the one in Baade’s Window and inferred that young stars must be present in significant numbers. Blum, Sellgren, & DePoy (1996b) obtained spectra of 19 of these brighter stars. Of those they could age-date, most were younger than 200 Myr.
The intermediate-age ($`2`$Gyr) and older stars ejected from the central cluster could be expected to be found up to a few $`r_0`$ from Sgr A\* on “Oort Cloud” like orbits. That is, a star from this population would keep getting jolted by the black holes to a more and more eccentric orbit until diffusive encounters near apocenter drove it into an orbit with a pericenter just beyond $`r_{\mathrm{bh}}`$. Narayanan, Gould, & Depoy (1996) list 16 presumably intermediate-age stars $`(K_0<5)`$ at projected radii of 8 to 20 pc. If these come primarily from the central pc, they should be preferentially on radial orbits. Of course, there has been recent star formation outside the central pc as well. For example, Blum et al. (1996b) estimate the ages IRS 24 and IRS 23 at $`100`$Myr and $`200`$Myr respectively, and these both lie at a projected distance of 1.7 pc. Hence, not all the Narayanan et al. (1996) stars are necessarily ejected.
Another interesting consequence of the ejection of low-mass stars from the center is that the capture of ordinary stars by supermassive black holes could be much rarer than commonly believed. These captures should lead to bright optical flares (e.g., Rees 1988; Ulmer 1999), which could be found in supernova searches. A cluster of black holes should have formed around all the supermassive black holes in galactic centers by the dynamical friction process described here, and these clusters would reduce by a large factor the rate at which stars can come close enough to the black hole to be tidally disrupted. Of course, since there is a central density cusp containing at least young stars around Sgr A\*, there will be some tidal captures, but the absence of old stars should reduce the number of expected flares. In addition, for galaxies such as ellipticals that are poor in neutral gas, one would not expect continuous star formation near the central black holes. Consequently, after the black hole cluster had ejected all the old stars, no young stars would form to replace them. For these galaxies, flares from stellar captures could be extremely rare.
We mention also microlensing of a background bulge star as another potentially observable effect, although requiring a large improvement in sensitivity and resolution of infrared imaging over our current capabilities. If a bulge star at a distance $`2\mathrm{kpc}`$ behind the center could be observed, the angular Einstein radius of Sgr A\* would be $`b=0.^{\prime \prime }8`$, corresponding to a linear size of $`r_E=0.03\mathrm{pc}`$. The two images of the star could therefore be comfortably resolved. The star would typically take several hundred years to complete the “microlensing event” by Sgr A\*. These two images would then have a microlensing optical depth to be lensed by one of the cluster black holes, $`\tau (N_{\mathrm{bh}}m_{\mathrm{bh}}/M_{BH})(r_E/r_{bh})^{5/4}A10^3A`$, where $`A`$ is the magnification of the image. Imaging down to $`K=21`$, one should on average find about one background star at $`2\mathrm{kpc}`$ within an Einstein radius of Sgr A\*, and about 100 similar stars near the core radius of $`1\mathrm{pc}`$ (for which the Einstein radius of Sgr A\* would be only $`0.^{\prime \prime }02`$). Of course, the identification of the two images of a star lensed by Sgr A\* would be of enormous interest by itself; despite the large number of orbiting black holes in the cluster, the expected rate of the “planet-like” events is still low, and several lensed stars would need to be identified to have a good chance of detecting the short events over a period of $`10`$ years.
Work by AG was supported in part by grant AST 97-27520 from the NSF. We are happy to acknowledge discussions with Martin J. Rees.
## Appendix A Diffusion Tensor in a Kepler Potential
Here, we calculate the diffusion tensor $`\kappa `$ for the general case of a $`\rho r^\alpha `$ density profile in a Kepler potential. assuming that the velocity distribution is isotropic. Recall that the phase space distribution is given by equations (15) and (16). We begin our evaluation at $`𝐯=0`$ where the diffusion tensor, $`\kappa _0`$, is isotropic. A single encounter with another black hole at speed $`u`$ changes the velocity by $`\mathrm{\Delta }v=2Gm_{\mathrm{bh}}/(bu)`$. Hence the time averaged growth of $`v^2`$ is
$$\frac{dv^2}{dt}=2\pi b𝑑bd^3u\left(\frac{2Gm_{\mathrm{bh}}}{bu}\right)^2uf_0(u,r)$$
$$=\frac{2(3\alpha )\alpha !}{(\alpha 1/2)!(1/2)!}\frac{N_{\mathrm{bh}}(Gm_{\mathrm{bh}})^2v_{\mathrm{esc}}^2}{(2GM_{\mathrm{cbh}}r_{\mathrm{bh}})^{3/2}}\left(\frac{r}{r_{\mathrm{bh}}}\right)^{3/2\alpha }\mathrm{ln}\mathrm{\Lambda }_1,$$
(A1)
where $`\mathrm{\Lambda }_1`$ is the ratio of maximum to minimum impact parameters, which we evaluate below.
If the diffusion equation is written (for isotropic $`\kappa _0`$) as $`\kappa _0_v^2f=f/t`$, then
$$\kappa _0=\frac{1}{6}\frac{dv^2}{dt},$$
(A2)
which combined with equation (A1) gives $`\kappa _0`$.
Next, we evalute $`\kappa _{}(v)`$ and $`\kappa _r(v)`$ for the special case of $`\alpha =3/2`$ for which the phase-space density (eqs. and ) is independent of both radius and velocity and consequently the problem is tractable analytically. Up to an overall normalization $`\gamma `$, the transverse velocity diffusion is given by
$$\frac{dv_{}^2}{dt}=\frac{\gamma }{2}_{v_{\mathrm{esc}}}^{v_{\mathrm{esc}}}𝑑u_r_{\theta =0}^{\theta _{\mathrm{max}}(u_r)}\left(\frac{1+\mathrm{cos}^2\theta }{2}\right)\frac{(u_r+v)^2d\mathrm{tan}^2\theta }{|u_r+v|\mathrm{sec}\theta },$$
(A3)
where $`\mathrm{cos}\theta _{\mathrm{max}}(u_r)|u_r+v|/(v_{\mathrm{esc}}^2+v^2+2u_rv)^{1/2}`$. For the total velocity diffusion $`dv_{}^2/dt`$, we find a similar expression, but with $`(1+\mathrm{cos}^2\theta )/21`$. We evaluate these expressions and find
$$\frac{dv^2}{dt}=\gamma v_{\mathrm{esc}}^2\left(1\frac{1}{3}\frac{v^2}{v_{\mathrm{esc}}^2}\right),\frac{dv_{}^2}{dt}=\frac{2}{3}\gamma v_{\mathrm{esc}}^2\left(1\frac{1}{5}\frac{v^2}{v_{\mathrm{esc}}^2}\right).$$
(A4)
We therefore conclude that $`\kappa _{}(v)/\kappa _0=10.2(v/v_{\mathrm{esc}})^2`$ and $`\kappa _r(v)/\kappa _0=10.6(v/v_{\mathrm{esc}})^2`$. These results apply to $`\alpha =3/2`$, but we adopt them for $`\alpha =7/4`$ as well. If the true coefficient for $`\alpha =7/4`$ is $`0.2(1+ϵ)`$ rather than 0.2, then this introduces an error into the capture formula (30) of only $`(1/6)ϵ`$. Since $`ϵ`$ is likely to be small, this correction is negligible.
We now justify two other approximations made in the capture calculation in § 3. First, in equation (27), we effectively considered the diffusion coefficient as fixed on slices of the velocity sphere perpendicular to the cylinder and passing through $`v_r`$. In fact, it is fixed on spherical shells of constant speed. This “plane parallel” approximation is justified by two related considerations. First, $`\kappa `$ varies very slowing, only by a factor $`1.25`$ over the entire velocity sphere. Second, most of the depleted region of velocity space is relatively near the cylinder, so that the error in $`\kappa `$ made by the plane-parallel approximation is extremely small.
Next, the boundary condition for the solution to the diffusion equation (27) sets the phase-space density $`f`$ at the edge of the velocity sphere $`(u=v_{\mathrm{esc}})`$ equal to the unperturbed density $`f_0`$ at $`u=v_r`$. Strictly speaking, if the boundary condition is set at $`(u=v_{\mathrm{esc}})`$, then one should use $`f_0(v_{\mathrm{esc}})`$. However, for the $`\alpha =7/4`$ profile, this is zero. (Our approximation would be exact for $`\alpha =3/2`$). The basis for our approximation is that the density returns essentially to $`f_0`$ long before the edge of the velocity sphere, at which point $`f_0`$ is not much different from $`f_0(v_r)`$.
Finally, we evaluate $`\mathrm{\Lambda }_1=b_{\mathrm{max}}/b_{\mathrm{min}}`$, the ratio of the maximum to minimum impact parameters. This can also be written $`\mathrm{\Lambda }_1=ϵ_{\mathrm{max}}/ϵ_{\mathrm{min}}`$, where $`(ϵ_{\mathrm{min}}v_{\mathrm{esc}},ϵ_{\mathrm{max}}v_{\mathrm{esc}})`$ is the range of impulses relevant to the problem at hand. For general relaxation $`ϵ_{\mathrm{max}}=1`$ and $`ϵ_{\mathrm{min}}2Gm_{\mathrm{bh}}/(rv_{\mathrm{esc}}^2)=m_{\mathrm{bh}}/M_{\mathrm{cbh}}`$. However, while the harder scatters all contribute to the overall relaxation, they do not contribute to diffusion into the loss cone because the black hole simply “jumps over” the capture cylinder. Unfortunately, it is not trivial to identify exactly what the largest allowed jumps should be: while jumps larger than $`(q_{\mathrm{min}}/r)^{1/2}v_{\mathrm{esc}}`$ do not directly lead to capture, they do help maintain the overall velocity profile given by equation (27). Note that the profile differs significantly from the background for many $`e`$-foldings. If the jumps larger than $`(q_{\mathrm{min}}/r)^{1/2}`$ did not contribute at all, then $`\mathrm{\Lambda }_1=(q_{\mathrm{min}}/r)^{1/2}\mathrm{\Lambda }`$. In view of the distribution’s slow approach to the background level near the capture cylinder, we adopt $`\mathrm{\Lambda }_1=\mathrm{\Lambda }(q_{\mathrm{min}}/r)^{1/4}=9.7`$.
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# NON-GAUSSIAN SIGNATURE INDUCED BY THE SZ EFFECT OF GALAXY CLUSTERS
## 1 Introduction
One of the major goals of cosmology is to understand the origin of the initial density perturbations. Do they come from inflation$`^\mathrm{?}`$ or from topological defects$`^\mathrm{?}`$? One way to answer this question is to make a statistical analysis of the CMB anisotropies. In fact, inflation predicts a Gaussian distribution of the primary anisotropies whereas the topological defects predict a non-Gaussian distribution. However, not only topological defects produce non-Gaussian signals. There are other astrophysical sources of non-gaussianity such as gravitational lensing (cf. F. Bernardeau’s contribution in this proceedings) and the secondary anisotropies such as the Sunyaev-Zel’dovich effect$`^\mathrm{?}`$ of galaxy clusters.
In this context, we have developed a method for the statistical analysis based on the wavelet decomposition of a signal. The method has been tested on an arbitrary set of Gaussian and non-Gaussian maps and then applied to CMB simulated maps to investigate the effect of the SZ contribution on the statistical signature of the anisotropies.
## 2 Detection method
Our method is based on a wavelet analysis of an image. That is, rather than analysing the signal in the direct space, we analyse the associated coefficients in the wavelet decomposition. The wavelet analysis can be viewed as a convolution of the signal by band pass filters: a scaling function similar to a low pas filter and a wavelet function similar to a high pass filter. The low pass filter will give an image with a lower resolution and the high pass will give the associated details. We have used the multi-scale wavelet analysis$`^\mathrm{?}`$ to decompose an input signal into a series of successive low resolution images and their associated detail images. At each level of the decomposition, the reference image has a resolution reduced by a factor of two and only this reference low pass image is decomposed at each level. This is the principle of the dyadic decomposition. We have chosen an anti-symmetric wavelet$`^\mathrm{?}`$ function similar to a first derivative operator because we expect that the non-Gaussian signatures arise from sharp gradients in the signal. Such a function is particularly sensitive to gradients and will detect them easily. The high pass filter applied to both directions of the image gives at each level three detail images corresponding to vertical, horizontal and diagonal details. The features of interest, particularly the statistical signatures, are studied at each scale (or level) for each type of details. We have qualified and tested our detection method on a set of Gaussian and non-Gaussian maps (100 of each) having the same power spectrum. This condition allows us to attribute the differences that are detected between the Gaussian and non-Gaussian maps to nothing but their statistical nature.
Our detection method$`^\mathrm{?}`$ is based on the three following major steps. First, we perform the wavelet decomposition of both 100 Gaussian maps and 100 non-Gaussian maps. We thus obtain at each decomposition level the wavelet coefficients associated with the vertical, horizontal and diagonal details. We also define the multi-scale gradient as the sum of the squared coefficients associated with the vertical and horizontal details. The non-Gaussian signature can be detected through the third or fourth order moment of the distribution, respectively the skewness and the kurtosis. In the following, we focus only on the excess of kurtosis because our signals are very moderately skewed. The second step of the method is to compute for each Gaussian and non-Gaussian map and at each level the excess of kurtosis of the wavelet coefficients associated with the corresponding details. We check that contrary to the Gaussian maps, the excess of kurtosis of the non-Gaussian signal is always centred around a non-zero value. We also note that the dispersion around the mean excess of kurtosis increases at high decomposition levels even for the Gaussian maps, due to the lack of coefficients at these levels. Therefore, in the following, we restrict our study to the first three levels.
The last step of the analysis is to quantify the detectability of the non-Gaussian signature. This is done by comparing the probability distribution functions (PDF) of the different processes. We compute the median excess of kurtosis of the 100 non-Gaussian maps at each decomposition scale and estimate the probability that it belongs to the PDF of the Gaussian maps. A low probability favours of a non-Gaussian signal whereas a high probability indicates a Gaussian process. We can also perform a more global comparison of the PDFs through the Kolmogorov-Smirnov test$`^\mathrm{?}`$. It gives, with a very good accuracy, the probability for a distribution to be different from a Gaussian.
## 3 Application to CMB
We apply our method to simulated CMB data including primary and secondary anisotropies. Our goal is to estimate the statistical non-Gaussian contamination induced by the secondary anisotropies. Therefore, we use an inflation model that generates Gaussian distributed primary anisotropies to which we add the simulated$`^\mathrm{?}`$ contributions due to the Sunyaev-Zel’dovich (thermal and kinetic) effect of galaxy clusters.
We have simulated 100 statistical realisations of the resulting CMB maps and performed the multi-scale decomposition. Following our proposed method, we compute the excess of kurtosis at different scales for the coefficients associated with the diagonal, vertical and horizontal details (Table 1).
At the first three decomposition scales, the excess of kurtosis is very large due to the SZ contribution. We also note, that the computations using the diagonal details are more sensitive to non-gaussianity and thus more powerful in detecting it. In fact, the galaxy clusters exhibit very peaked profiles or even point-like behaviour. The diagonal details are very sensitive to symmetric profiles. We find that the SZ effect a major source of non-gaussianity among the secondary effects.
## 4 Effects of the instrumental configurations
We apply our statistical discriminators to test for non-gaussianity within the context of the representative instrumental configuration of the future Planck Surveyor satellite for CMB observations.
We use the same astrophysical contributions as those described above (primary and SZ). The difference lies in the fact that the maps are convolved with a 6 arcminute Gaussian beam. We also take into account the expected Gaussian noise of Planck ($`(\delta T/T)_{rms}\mathrm{2.\hspace{0.17em}10}^6`$ per 1.5 arcminute pixel). The convolution by a 6 arcminute beam suppresses the power at the corresponding scale (Scale I) and affects the second decomposition scale whereas the third is not significantly altered. At the third decomposition scale, we find for the multi-scale gradient $`k=0.62_{0.60}^{+1.43}`$. Whereas we find for the horizontal and vertical details, and for the diagonal details respectively, $`k=0.07_{0.08}^{+0.11}`$ and $`0.16\pm 0.10`$. In order to quantify the detectability of non-gaussianity in the Planck-like configuration, we generate Gaussian distributed maps with same power spectrum as the studied signal. We plot (Fig. 1) the PDF of the Gaussian (dashed line) and non-Gaussian (solid line) processes. We derive the probability that the median excess of kurtosis measured on the “real sky” belongs to the Gaussian process. Using the multi-scale gradient, we find that the probability of detecting non-gaussianity is 71.9% at the second decomposition scale. There is no significant detection elsewhere. Whereas using the coefficients of the diagonal details, the probability of detecting a non-Gaussian signature at the third scale is 94.5%. We apply the K-S test to the distribution of the excess of kurtosis for the diagonal details and find a probability of 96.6% of detecting non-gaussianity. Since the comparison of the two distributions using the K-S test is very sensitive to departures from gaussianity. It thus gives better results on the detection of the non-Gaussian signature.
## 5 Conclusions
In the present study, we investigate the statistical signature induced by the SZ effect of galaxy clusters when the primary anisotropies result from an inflationary scenario and are thus Gaussian distributed. We use discriminators based on the statistical properties of the coefficients in a four level wavelet decomposition. In our study, we find that the SZ effect of clusters generates a very large non-Gaussian signature that dominates by far all other secondary anisotropies. We apply our statistical tests to a Planck-like configuration in order to estimate the capabilities of the satellite to detect the non-Gaussian signature of the SZ effect. In this case, we detect unambiguously the non-Gaussian signature at the third decomposition scale ($`12`$ arcminutes), the first ($`3`$ arcminutes) and second ($`6`$ arcminutes) scales being affected by the beam convolution.
## References
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# Error-correcting codes and image restoration with multiple stages of dynamics
## I INTRODUCTION
The corruption of signals by noise is a common problem encountered in information processing. To retrieve signals from messages corrupted during the transmission through noisy channels, various error-correcting codes have been proposed . In particular, the error-correction mechanism of a class of parity-checking codes can be considered as the search for thermodynamically stable states of a Hamiltonian constructed in terms of the message bits . These codes have been demonstrated to saturate the Shannon information bound in the limit that each encoded bit checks the parity of an infinitely large number of message bits . While in practice, each encoded bit can only check the parity of a finite number of message bits, these codes still maintain a very low bit error probability.
The need to retrieve signals from corrupted messages is also inherent in image restoration . Although parity-checking bits may not be explicitly introduced for the task, prior knowledge about the images plays a similar role. For example, the smoothness of real-world images provides a mechanism for checking the pixel values in comparison with those of their neighbors. A corresponding Hamiltonian, consisting of a ferromagnetic bias to reflect the smoothening tendency, can be constructed in terms of the image pixels. Modern techniques of image restoration based on Markov random fields correspond to the search for thermodynamically stable states of the Hamiltonian system, using methods such as simulated annealing .
In a recent paper, we have shown that the problems of error-correcting codes and image restoration can be formulated in a unified framework . In both tasks, the choice of the so-called hyperparameters is an important factor in determining their performances. Hyperparameters refer to the coefficients of the various interactions appearing in the Hamiltonian of the tasks. In error correction, they determine the statistical significance given to the parity-checking terms and the received bits. Similarly in image restoration, they determine the statistical weights given to the prior knowledge and the received data. It was shown, by the use of inequalities, that the optimal choice of the hyperparameters correspond to the Maximum Posterior Marginal (MPM) method, where there is a match between the source and model priors. The choice of these values correspond to the Nishimori point in the space of hyperparameters . It is equivalent to a thermodynamic process at finite temperature, and the task performance is better than the Maximum A Posteriori probability (MAP) method, where the values of the hyperparameters are taken to infinity, equivalent to a zero temperature process. Furthermore, from the analytic solution of the infinite-range model and the Monte Carlo simulation of finite-dimensional models, it was shown that an inappropriate choice of the hyperparameters can lead to a rapid degradation of the tasks.
In fact, hyperparameter estimation has been the subject of many previous studies , a recently popular one using the “evidence framework” . However, if the prior models the source poorly, no hyperparameters can be reliable . Even if they can be estimated accurately through steady-state statistical measurements, they may fluctuate when interfered by bursty noise sources in communication channels. Hence it is important to devise decoding or restoration procedures which are robust against the uncertainties in hyperparameter estimation.
In this paper we propose the technique of selective freezing as a method to increase the tolerance to uncertainties in hyperparameter estimation. The technique has been studied for pattern reconstruction in neural networks, where it led to an improvement in the retrieval precision, a widening of the basin of attraction, and a boost in the storage capacity . The idea is best illustrated for Ising bits or pixels with binary states $`\pm 1`$, though it can be easily generalized to other cases. In a finite temperature thermodynamic process, the Ising variables keep moving under thermal agitation. Some of them have smaller thermal fluctuations than the others, implying that they are more certain to stay in one state than the other. This stability implies that they have a higher probability to stay in the correct state for error-correction or image restoration tasks, even when the hyperparameters are not optimally tuned. It may thus be interesting to separate the thermodynamic process into two stages. In the first stage we select those relatively stable bits or pixels whose time-averaged states have a magnitude exceeding a certain threshold. In the second stage we subsequently fix (or freeze) them in the most probable thermodynamic states (for Ising variables this corresponds to the sign of the time-averaged state). Thus these selectively frozen bits or pixels are able to provide a more robust assistance to the less stable bits or pixels in their search for the most probable states. The selective freezing procedure reduces to the usual finite-temperature decoding or restoration process if all bits or pixels are frozen (since nothing happens in the second stage), or no bits or pixels are frozen (since the second stage is merely a continuation of the equilibration process of the first stage).
The two-stage thermodynamic process can be studied analytically in the mean-field model, which provides a qualitative guide to the behavior of more realistic cases of lower dimensions. However, it is necessary to give a remark about the theoretical approach. That is, as far as we have tried, the analytical solution has been inaccessible by the more conventional replica method. Rather, we have to use the cavity method to obtain the equations for the order parameters. In particular, the cavity method leads to the appearance of a term called the trans-susceptibility, which correctly describes the effects of the thermodynamics of the first stage on that of the second.
The paper is organized as follows. In Section II we briefly review the formulation of error-correcting codes and image restoration in a unified framework. In Sections III and IV, we consider the mean-field model for error-correcting codes and image restoration respectively. We derive the equations for the order parameters of the two-stage thermodynamics using the cavity method, and present numerical results illustrating the robustness of selective freezing against uncertainties in hyperparameter estimation. We further demonstrate that even when the noise model changes without the receiver/restoration agent realizing the change (i.e. it makes a wrong estimation of the prior), the task performance is still robust. For the more realistic cases of lower dimensions, simulation results illustrate the relevance of the infinite-range model in providing qualitative guidance. The conclusion is given in Section V.
## II FORMULATION
Consider an information source which generates data represented by a set of Ising spins $`\{\xi _i\}`$, where $`\xi _i=\pm 1`$ and $`i=1,\mathrm{},N`$. The data is generated according to the source prior $`P_s(\{\xi _i\})`$. For error-correcting codes transmitting unbiased messages, all sequences are equally probable and $`P_s(\{\xi \})=2^N`$. For images with smooth structures, the prior consists of ferromagnetic Boltzmann factors, which increase the tendencies of the neighboring spins to stay at the same spin states, that is,
$$P_s(\{\xi \})=\frac{1}{Z(\beta _s)}\mathrm{exp}\left(\frac{\beta _s}{z}\underset{ij}{}\xi _i\xi _j\right).$$
(1)
Here $`ij`$ represents pairs of neighboring spins, $`z`$ is the valency of each site, and the partition function $`Z(\beta _s)`$ is given by
$$Z(\beta _s)=\mathrm{Tr}_\xi \mathrm{exp}\left(\frac{\beta _s}{z}\underset{ij}{}\xi _i\xi _j\right).$$
(2)
The data is coded by constructing the codewords, which are the products of $`p`$ spins $`J_{i_1\mathrm{}i_p}^0=\xi _{i_1}\mathrm{}\xi _{i_p}`$ for appropriately chosen sets of of indices $`\{i_1,\mathrm{},i_p\}`$, the choice of which determines the type of code. Each spin may appear in a number of $`p`$-spin codewords; the number of times of appearance is called the valency $`z_p`$. The Sourlas code is equivalent to the infinite-range model in which all possible codewords of $`p`$ spins are chosen from $`N`$ spins. On the other hand, the Kabashima-Saad code consists of combinations in which each spin appears in a random pre-selection of $`z_p`$ codewords. For conventional image restoration, codewords with only $`p=1`$ are transmitted, corresponding to the pixels in the image; the inclusion of terms with $`p>1`$, and their positive effects on restoring the original image, have also been discussed in . For simplicity, we restrict ourselves to the case of a single non-vanishing value of $`p`$ with $`p2`$, and $`p=1`$.
When the signal is transmitted through a noisy channel, the output consists of the sets $`\{J_{i_1\mathrm{}i_p}\}`$ and $`\{\tau _i\}`$, which are the corrupted versions of $`\{J_{i_1\mathrm{}i_p}^0\}`$ and $`\{\xi _i\}`$ respectively. In the binary symmetric channel, the outputs $`J_{i_1\mathrm{}i_p}`$ are equal to $`J_{i_1\mathrm{}i_p}^0`$ with probabilities $`p_J`$ and $`1p_J`$ respectively, and $`\tau _i`$ equal to $`\xi _i`$ with probabilities $`p_\tau `$ and $`1p_\tau `$ respectively. Thus
$$P_{\mathrm{out}}(\{J\},\{\tau \}|\{\xi \})\mathrm{exp}\left(\beta _JJ_{i_1\mathrm{}i_p}\xi _{i_1}\mathrm{}\xi _{i_p}+\beta _\tau \tau _i\xi _i\right),$$
(3)
where
$$\beta _J=\frac{1}{2}\mathrm{ln}\frac{1p_J}{p_J}\mathrm{and}\beta _\tau =\frac{1}{2}\mathrm{ln}\frac{1p_\tau }{p_\tau }.$$
(4)
The first summation in the exponent of Eq. (3) extends over an appropriate set of the indices $`(i_1,\mathrm{},i_p)`$.
The Gaussian channel is defined by, for a given sequence $`\{\xi _i\}`$,
$$P_{\mathrm{out}}(\{J\},\{\tau \}|\{\xi \})\mathrm{exp}\left(\frac{1}{2J^2}(J_{i_1\mathrm{}i_p}J_0\xi _{i_1}\mathrm{}\xi _{i_p})^2\frac{1}{2\tau ^2}(\tau _ia\xi _i)^2\right).$$
(5)
$`J_0`$ and $`a`$ are the strengths of the signals to be fed into the channel, and $`J^2`$ and $`\tau ^2`$ are the variances of the noise. We note that by letting $`\beta _J`$ and $`\beta _\tau `$ to be $`J_0/J^2`$ and $`a/\tau ^2`$ respectively, the input-dependent terms of Eq. (5) reduce to those of Eq. (3), which therefore can be regarded as the noise model for both binary symmetric and Gaussian channels.
According to Bayesian statistics, the posterior probability that the source sequence is $`\{\sigma \}`$, given the outputs $`\{J\}`$ and $`\{\tau \}`$, takes the form
$$P(\{\sigma \}|\{J\},\{\tau \})P_{\mathrm{out}}(\{J\},\{\tau \}|\{\sigma \})P_s(\{\sigma \}).$$
(6)
Using Eq. (3) and (1), we have
$$P(\{\sigma \}|\{J\},\{\tau \})\mathrm{exp}\left(\beta _JJ_{i_1\mathrm{}i_p}\sigma _{i_1}\mathrm{}\sigma _{i_p}+\beta _\tau \tau _i\sigma _i+\frac{\beta _s}{z}\underset{ij}{}\sigma _i\sigma _j\right).$$
(7)
It often happens that the receiver at the end of the noisy channel does not have precise information on $`\beta _J`$, $`\beta _\tau `$ or $`\beta _s`$. One then has to estimate these parameters. If the receiver estimates $`\beta `$, $`h`$ and $`\beta _m`$ for $`\beta _J`$, $`\beta _\tau `$ and $`\beta _s`$ respectively, then the mean of the posterior distribution of $`\sigma _i`$ is equal to the thermal average
$$\sigma _i=\frac{\mathrm{Tr}\sigma _ie^{H\{\sigma \}}}{\mathrm{Tr}e^{H\{\sigma \}}},$$
(8)
where the Hamiltonian is given by
$$H\{\sigma \}=\beta J_{i_1\mathrm{}i_p}\sigma _{i_1}\mathrm{}\sigma _{i_p}h\tau _i\sigma _i\frac{\beta _m}{z}\underset{ij}{}\sigma _i\sigma _j.$$
(9)
One then regards $`\mathrm{sgn}\sigma _i`$ as the $`i`$th bit of the decoded/restored information.
To reduce the sensitivity of the decoding/restoration process to the uncertainties in parameter estimation, we propose a two-stage process of selective freezing instead of the one-stage thermodynamic process implied by Eq. (8). In the first stage the spins evolve thermodynamically as prescribed in Eq. (8), and the thermal averages $`\sigma _i`$ of the spins are monitored. We may relate $`\sigma _i`$ to an effective field $`H_i`$ by $`\sigma _i=\mathrm{tanh}H_i`$. Spins with larger magnitudes of $`\sigma _i`$ correspond to larger magnitudes of $`H_i`$. They are more likely to agree with the correct message or image bit, and are less likely to change signs even when the hyperparameters vary. Their relative stability can be used to assist the less stable spins to boost their robustness against hyperparameter uncertainties. Hence we select those spins with $`|\sigma _i|`$ exceeding a given threshold $`\theta `$, and freeze them in the second stage of the thermodynamics. The average of the spin $`\stackrel{~}{\sigma }_i`$ in the second stage is then given by
$$\stackrel{~}{\sigma }_i=\frac{\mathrm{Tr}\stackrel{~}{\sigma }_i\underset{j}{}\left[\mathrm{\Theta }\left(\sigma _j^2\theta ^2\right)\delta _{\stackrel{~}{\sigma }_j,\mathrm{sgn}\sigma _j}+\mathrm{\Theta }\left(\theta ^2\sigma _j^2\right)\right]e^{\stackrel{~}{H}\{\stackrel{~}{\sigma }\}}}{\mathrm{Tr}_j\left[\mathrm{\Theta }\left(\sigma _j^2\theta ^2\right)\delta _{\stackrel{~}{\sigma }_j,\mathrm{sgn}\sigma _j}+\mathrm{\Theta }\left(\theta ^2\sigma _j^2\right)\right]e^{\stackrel{~}{H}\{\stackrel{~}{\sigma }\}}},$$
(10)
where $`\mathrm{\Theta }`$ is the step function, $`\stackrel{~}{H}\{\stackrel{~}{\sigma }\}`$ is the Hamiltonian for the second stage, and has the same form as Eq. (9) in the first stage. To increase the flexibility in the process, the parameters $`\beta `$, $`h`$ and $`\beta _m`$ can be replaced by $`\stackrel{~}{\beta }`$, $`\stackrel{~}{h}`$ and $`\stackrel{~}{\beta }_m`$ respectively in the second stage. One then regards $`\mathrm{sgn}\stackrel{~}{\sigma }_i`$ as the $`i`$th spin of the decoding/restoration process.
The most important quantity in selective freezing is the overlap of the decoded/restored bit $`\mathrm{sgn}\stackrel{~}{\sigma }_i`$ and the original bit $`\xi _i`$ averaged over the output probability and the spin distribution. This is given by
$$M_{\mathrm{sf}}=\underset{\xi }{}𝑑J𝑑\tau P_s(\{\xi \})P_{\mathrm{out}}(\{J\},\{\tau \}|\{\xi \})\xi _i\mathrm{sgn}\stackrel{~}{\sigma }_i.$$
(11)
Following Appendix A of , we can prove the following inequality
$$M_{\mathrm{sf}}M(\beta =\beta _J,h=\beta _\tau ,\beta _m=\beta _s),$$
(12)
where the right hand side is the overlap of the single-stage dynamics when the model parameters $`\beta `$, $`h`$ and $`\beta _m`$ match the source parameters $`\beta _J`$, $`\beta _\tau `$ and $`\beta _s`$ respectively. Hence selective freezing cannot outperform the single-stage process if the hyperparameters can be estimated precisely. However, we remark that the purpose of selective freezing is rather to provide a relatively stable performance when the hyperparameters cannot be estimated precisely. This cannot be revealed from the inequality, but will be confirmed by the analytic and simulation results in Sections III and IV.
## III THE INFINITE-RANGE MODEL FOR ERROR-CORRECTING CODES
Let us now suppose that the output of the transmission channel consists of only the set of $`p`$-spin interactions $`\{J_{i_1\mathrm{}i_p}\}`$. The Hamiltonian (9) then becomes
$$H\{\sigma \}=\beta \underset{i_1<\mathrm{}<i_p}{}J_{i_1\mathrm{}i_p}\sigma _{i_1}\mathrm{}\sigma _{i_p},$$
(13)
where we have set $`\beta _m=0`$ for the case that all messages are equally probable.
Analytical solutions for the overlap are in general unavailable. We therefore consider the infinite-range model in which the exchange interactions are present for all possible pairs of sites in the Hamiltonian of Eq. (13).
To consider the transition between error-free and errored regimes, we are interested in the noise model in which $`J_{i_1\mathrm{}i_p}`$ is Gaussian with mean $`p!j_0\xi _{i_1}\mathrm{}\xi _{i_p}/N^{p1}`$ and variance $`p!J^2/2N^{p1}`$. Since all messages are equally probable, we can apply a gauge transformation $`\sigma _i\sigma _i\xi _i`$ and $`J_{i_1\mathrm{}i_p}J_{i_1\mathrm{}i_p}\xi _{i_1}\mathrm{}\xi _{i_p}`$ to (13), and arrive at an equivalent $`p`$-spin model with a ferromagnetic bias, where
$$P(J_{i_1\mathrm{}i_p})=\left(\frac{N^{p1}}{\pi J^2p!}\right)^{1/2}\mathrm{exp}\left[\frac{N^{p1}}{J^2p!}\left(J_{i_1\mathrm{}i_p}\frac{p!}{N^{p1}}j_0\right)^2\right].$$
(14)
The Nishimori point for this model is located at $`\beta =2j_0/J^2`$.
The infinite-range model is exactly solvable using mean-field theoretical techniques for disordered systems such as the replica or cavity method . Here we use the cavity method because of its more transparent physical interpretation, and some obstacles encountered in the use of the replica method.
The cavity method uses a self-consistency argument to consider what happens when a spin is added or removed from the system. The central quantity in this method is the cavity field, which is the local field of a spin when it is added to the system, assuming that the exchange couplings act only one-way from the system to the new spin (but not from the spin back to the system). Since the exchange couplings feeding the new spin have no correlations with the system, the cavity field becomes a Gaussian variable in the limit of large valency.
### A Average spin in the first stage
We start with the so-called “clustering property” for mean-field systems ,
$$\sigma _{i_1}\mathrm{}\sigma _{i_p}=\sigma _{i_1}\mathrm{}\sigma _{i_p},$$
(15)
where $``$ represents thermodynamic averages. As shown in Appendix A, the clustering property enables us to express the thermal averages of a spin in terms of the cavity field, say, for spin 1,
$$\sigma _1=\mathrm{tanh}\beta h_1;h_1=\underset{1<j_2<\mathrm{}<j_p}{}J_{1j_2\mathrm{}j_p}\sigma _{j_2}^{\backslash 1}\mathrm{}\sigma _{j_p}^{\backslash 1},$$
(16)
where the superscript $`\backslash 1`$ denotes the thermal averages for a Hamiltonian in which $`\sigma _1`$ and the associated exchange interactions are absent, but otherwise identical to Eq. (13). Thus $`h_1`$ is the cavity field obeying a Gaussian distribution, whose mean and variance are $`pj_0m^{p1}`$ and $`pJ^2q^{p1}/2`$ respectively, where $`m`$ and $`q`$ are the magnetization and Edwards-Anderson order parameter respectively, given by
$$m\frac{1}{N}\underset{i}{}\sigma _i\mathrm{and}q\frac{1}{N}\underset{i}{}\sigma _i^2.$$
(17)
It is convenient to write
$$\beta h_i=\widehat{m}+\sqrt{\widehat{q}}u_i,$$
(18)
where
$$\widehat{m}=p\beta j_0m^{p1}\mathrm{and}\widehat{q}=\frac{p}{2}\beta ^2J^2q^{p1},$$
(19)
and $`u_i`$ is a Gaussian variable with mean 0 and variance 1.
### B Order parameters in the first stage
Applying self-consistently the cavity argument to all terms in Eq. (17), we can obtain self-consistent equations for $`m`$ and $`q`$:
$`m`$ $`=`$ $`{\displaystyle Du\mathrm{tanh}G},`$ (20)
$`q`$ $`=`$ $`{\displaystyle Du\mathrm{tanh}^2G},`$ (21)
where $`Dudue^{u^2/2}/\sqrt{2\pi }`$ is the Gaussian measure, $`G=\widehat{m}+\sqrt{\widehat{q}}u`$. The overlap for the one-stage decoding process is given by
$$M\frac{1}{N}\underset{i}{}\mathrm{sgn}\sigma _i=\mathrm{erf}\frac{\widehat{m}}{\sqrt{2\widehat{q}}}.$$
(22)
Now we consider selective freezing. If we introduce a freezing threshold $`\theta `$ so that all spins with $`\sigma _i^2>\theta ^2`$ are frozen, then the freezing fraction $`f`$ is given by
$$f\frac{1}{N}\underset{i}{}\mathrm{\Theta }\left(\sigma _i^2\theta ^2\right)=1\frac{1}{2}\mathrm{erf}\frac{u_+}{\sqrt{2}}+\frac{1}{2}\mathrm{erf}\frac{u_{}}{\sqrt{2}},$$
(23)
where $`u_\pm =(\pm u_0\widehat{m})/\sqrt{\widehat{q}}`$ with $`\mathrm{tanh}u_0=\theta `$.
### C Average spin in the second stage
Assuming that the spin $`\stackrel{~}{\sigma }_1`$ is dynamic in the second stage, we can write
$$H\{\stackrel{~}{\sigma }\}H\{\stackrel{~}{\sigma }\}^{\backslash 1}\stackrel{~}{\beta }\underset{1<j_1\mathrm{}<j_{p1}}{}\stackrel{~}{\sigma }_1J_{1j_1\mathrm{}j_{p1}}\underset{s=1}{\overset{p1}{}}\left[\stackrel{~}{\sigma }_{j_s}\mathrm{\Theta }\left(\theta ^2\sigma _{j_s}^2\right)+\mathrm{sgn}\sigma _{j_s}\mathrm{\Theta }\left(\sigma _{j_s}^2\theta ^2\right)\right],$$
(24)
where $`H\{\stackrel{~}{\sigma }\}^{\backslash 1}`$ is the Hamiltonian when spin 1 is completely removed from the system in both stages of the thermodynamic process. Removing spin 1 may cause the thermal averages of other spins to adjust slightly in the first stage. Hence some dynamic spins (with $`\sigma _k^2<\theta ^2`$) may become frozen ones (with $`\sigma _k^2>\theta ^2`$) and vice versa, so that strictly speaking, further terms should be considered in Eq. (24) to account for these secondary effects. For example, if spin $`k`$ is induced to switch from dynamic to frozen (or vice versa) on removal of spin 1, then the Taylor expansion of $`H\{\stackrel{~}{\sigma }\}`$ implies that an extra term
$`\stackrel{~}{\beta }\left(\mathrm{sgn}\sigma _k^{\backslash 1}\stackrel{~}{\sigma }_k\right)\left[\delta (\sigma _k^{\backslash 1}\theta )\delta (\sigma _k^{\backslash 1}+\theta )\right](\sigma _k\sigma _k^{\backslash 1})`$ (25)
$`{\displaystyle \underset{1<j_1\mathrm{}<j_{p1}k}{}}J_{kj_1\mathrm{}j_{p1}}{\displaystyle \underset{s=1}{\overset{p1}{}}}\left\{\stackrel{~}{\sigma }_{j_s}\mathrm{\Theta }\left[\theta ^2(\sigma _{j_s}^{\backslash 1})^2\right]+\mathrm{sgn}\sigma _{j_s}^{\backslash 1}\mathrm{\Theta }\left[(\sigma _{j_s}^{\backslash 1})^2\theta ^2\right]\right\}`$ (26)
should be incorporated in Eq. (24). Here, we have neglected these terms for clarity. Nevertheless, justification a posteriori can be provided for their deletion.
Using a cavity argument similar to Appendix A, we can show that
$$\stackrel{~}{\sigma }_1=\mathrm{tanh}\stackrel{~}{\beta }\left\{\underset{1<j_1\mathrm{}<j_{p1}}{}J_{1j_1\mathrm{}j_{p1}}\underset{s=1}{\overset{p1}{}}\left[\stackrel{~}{\sigma }_{j_s}^{\backslash 1}\mathrm{\Theta }\left(\theta ^2\sigma _{j_s}^2\right)+\mathrm{sgn}\sigma _{j_s}^{\backslash 1}\mathrm{\Theta }\left(\sigma _{j_s}^2\theta ^2\right)\right]\right\}.$$
(27)
However, the effective field on the right hand side of Eq. (27) is still not a cavity field because $`\sigma _{j_s}`$, which are used in the step functions to decide whether the spin $`j_s`$ is dynamic or frozen in the second stage, is different from $`\sigma _{j_s}^{\backslash 1}`$. Hence it may have correlations with spin 1. Taylor expansion of $`\sigma _{j_s}`$ about $`\sigma _{j_s}^{\backslash 1}`$ yields
$`\stackrel{~}{\sigma }_1=\mathrm{tanh}\stackrel{~}{\beta }\{\stackrel{~}{h}_1+{\displaystyle \underset{1jj_1\mathrm{}<j_{p2}}{}}J_{1jj_1\mathrm{}j_{p2}}`$ (28)
$`{\displaystyle \underset{s=1}{\overset{p2}{}}}\left[\stackrel{~}{\sigma }_{j_s}^{\backslash 1}\mathrm{\Theta }\left[\theta ^2(\sigma _{j_s}^{\backslash 1})^2\right]+\mathrm{sgn}\sigma _{j_s}^{\backslash 1}\mathrm{\Theta }\left[(\sigma _{j_s}^{\backslash 1})^2\theta ^2\right]\right]`$ (29)
$`[\mathrm{sgn}\sigma _j^{\backslash 1}\stackrel{~}{\sigma }_j^{\backslash 1}][\delta (\sigma _j^{\backslash 1}\theta )\delta (\sigma _j^{\backslash 1}+\theta )](\sigma _j\sigma _j^{\backslash 1})\},`$ (30)
where $`\stackrel{~}{h}_1`$ is the generic cavity field which is now completely uncorrelated with spin 1. It is given by
$$\stackrel{~}{h}_1=\underset{1<j_1\mathrm{}<j_{p1}}{}J_{1j_1\mathrm{}j_{p1}}\underset{s=1}{\overset{p1}{}}\left\{\stackrel{~}{\sigma }_{j_s}^{\backslash 1}\mathrm{\Theta }\left[\theta ^2(\sigma _{j_s}^{\backslash 1})^2\right]+\mathrm{sgn}\sigma _{j_s}^{\backslash 1}\mathrm{\Theta }\left[(\sigma _{j_s}^{\backslash 1})^2\theta ^2\right]\right\}.$$
(31)
To evaluate the difference $`\sigma _j\sigma _j^{\backslash 1}`$ appearing in Eq. (30), we have to apply the cavity method a second time, by comparing the changes when both spins 1 and $`j`$ are removed. This is done in Appendix B and the result is
$$\sigma _j\sigma _j^{\backslash 1}=\left(\beta \mathrm{sech}^2\beta h_j^{\backslash 1}\right)\left(h_{j1}\mathrm{tanh}\beta h_1^{\backslash j}\right),$$
(32)
where
$$h_{1j}=h_{j1}=\underset{1jk_1\mathrm{}<k_{p2}}{}J_{1jk_1\mathrm{}k_{p2}}\sigma _{k_1}^{\backslash 1j}\mathrm{}\sigma _{k_{p2}}^{\backslash 1j}.$$
(33)
When Eqs. (32-33) are substituted into Eq. (30), the significant contribution comes from the terms which pair up $`J_{1jj_1\mathrm{}j_{p2}}`$ and $`J_{1jk_1\mathrm{}k_{p2}}`$. The various terms appearing in the summation over $`jj_1<\mathrm{}<j_{p2}`$ involve thermal averages in the absence of spins 1 or $`j`$. We assume that the effects of removing a spin is negligible (which can be shown to be equivalent to the replica symmetric approximation in the replica method ). Then replacing the components of the terms by their mean values, and counting that $`N^{p2}/(p2)!`$ terms appearing in the summation over $`j_1<\mathrm{}<j_{p2}`$, we arrive at
$`\stackrel{~}{\sigma }_1=\mathrm{tanh}\stackrel{~}{\beta }\{\stackrel{~}{h}_1+{\displaystyle \frac{p}{2}}(p1)J^2{\displaystyle \frac{1}{N}}{\displaystyle \underset{j}{}}[\delta (\sigma _j\theta )\delta (\sigma _j+\theta )][\mathrm{sgn}\sigma _j\stackrel{~}{\sigma }_j]`$ (34)
$`\left(\beta \mathrm{sech}^2\beta h_j\right)\left(r^{p2}\mathrm{tanh}\beta h_1\right)\},`$ (35)
where $`r`$ is the order parameter describing the spin correlations of the two thermodynamic stages:
$$r\frac{1}{N}\underset{i}{}\sigma _i\left\{\stackrel{~}{\sigma }_i\mathrm{\Theta }\left[\theta ^2\sigma _i^2\right]+\mathrm{sgn}\sigma _i\mathrm{\Theta }\left[\sigma _i^2\theta ^2\right]\right\}.$$
(36)
Eq. (35) can be simplified by introducing the trans-susceptibility $`\chi _{tr}`$, which describes the response of a spin in the second stage to variations of the cavity field in the first stage, namely
$$\chi _{tr}\frac{1}{N}\underset{i}{}\frac{\stackrel{~}{\sigma }_i}{h_i}.$$
(37)
Since $`\stackrel{~}{\sigma }_i`$ equals $`\mathrm{sgn}h_i`$ for $`\mathrm{tanh}^2\beta h_i>\theta ^2`$, and $`\mathrm{tanh}\beta \stackrel{~}{h}_i`$ otherwise, we get
$$\chi _{tr}=\frac{1}{N}\underset{i}{}\left[\delta (\sigma _i\theta )\delta (\sigma _i+\theta )\right]\left[\mathrm{sgn}\sigma _i\stackrel{~}{\sigma }_i\right]\beta \mathrm{sech}^2\beta h_i.$$
(38)
Eq. (35) can thus be simplified to
$$\stackrel{~}{\sigma }_1=\mathrm{tanh}\stackrel{~}{\beta }\left\{\stackrel{~}{h}_1+\frac{p}{2}(p1)J^2r^{p2}\chi _{tr}\mathrm{tanh}\beta h_1\right\}.$$
(39)
### D Order parameters in the second stage
The cavity field $`\stackrel{~}{h}_1`$ in the second stage is a Gaussian variable. Its mean and variance are $`pj_0\stackrel{~}{m}^{p1}`$ and $`pJ^2\stackrel{~}{q}^{p1}/2`$ respectively, where $`\stackrel{~}{m}`$ and $`\stackrel{~}{q}`$ are the magnetization and Edwards-Anderson order parameter respectively, given by
$`\stackrel{~}{m}`$ $``$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{i}{}}\left[\mathrm{\Theta }(\theta ^2\sigma _i^2)\stackrel{~}{\sigma }_i+\mathrm{\Theta }(\sigma _i^2\theta ^2)\mathrm{sgn}\sigma _i\right],`$ (40)
$`\stackrel{~}{q}`$ $``$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{i}{}}\left[\mathrm{\Theta }(\theta ^2\sigma _i^2)\stackrel{~}{\sigma }_i^2+\mathrm{\Theta }(\sigma _i^2\theta ^2)\right].`$ (41)
Furthermore, the covariance between $`h_1`$ and $`\stackrel{~}{h}_1`$ is $`pJ^2r^{p1}/2`$, where $`r`$ is given in Eq. (36).
Algebraic manipulations can be simplified if we write, for $`i=1`$,
$`\beta h_i`$ $`=`$ $`\widehat{m}+\sqrt{\widehat{q}}u_i,`$ (42)
$`\stackrel{~}{\beta }\stackrel{~}{h}_i`$ $`=`$ $`\widehat{\stackrel{~}{m}}+\sqrt{\widehat{\stackrel{~}{q}}}(\eta u_i+\sqrt{1\eta ^2}v_i),`$ (43)
where $`u_i`$ and $`v_i`$ are independent Gaussian variables with mean 0 and variance 1, $`\widehat{m}`$, $`\widehat{q}`$ are given in Eq. (19), and
$`\widehat{\stackrel{~}{m}}=p\stackrel{~}{\beta }j_0\stackrel{~}{m}^{p1},\mathrm{and}\widehat{\stackrel{~}{q}}={\displaystyle \frac{p}{2}}\stackrel{~}{\beta }^2J^2\stackrel{~}{q}^{p1},`$ (44)
$`\widehat{r}={\displaystyle \frac{p}{2}}\beta \stackrel{~}{\beta }J^2r^{p1},\mathrm{and}\eta ={\displaystyle \frac{\widehat{r}}{\sqrt{\widehat{q}\widehat{\stackrel{~}{q}}}}}.`$ (45)
Applying self-consistently the same cavity argument to all terms in Eqs. (40), (41), (36) and (38) and performing the Gaussian average over $`u_i`$ and $`v_i`$, we arrive at the following self-consistent equations for $`\stackrel{~}{m}`$, $`\stackrel{~}{q}`$, $`r`$ and $`\chi _{tr}`$:
$`\stackrel{~}{m}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{erf}{\displaystyle \frac{u_+}{\sqrt{2}}}{\displaystyle \frac{1}{2}}\mathrm{erf}{\displaystyle \frac{u_{}}{\sqrt{2}}}+{\displaystyle _u_{}^{u_+}}Du{\displaystyle Dv\mathrm{tanh}L},`$ (46)
$`\stackrel{~}{q}`$ $`=`$ $`1{\displaystyle \frac{1}{2}}\mathrm{erf}{\displaystyle \frac{u_+}{\sqrt{2}}}+{\displaystyle \frac{1}{2}}\mathrm{erf}{\displaystyle \frac{u_{}}{\sqrt{2}}}+{\displaystyle _u_{}^{u_+}}Du{\displaystyle Dv\mathrm{tanh}^2L},`$ (47)
$`r`$ $`=`$ $`\left({\displaystyle _{\mathrm{}}^u_{}}+{\displaystyle _{u_+}^{\mathrm{}}}\right)Du\left|\mathrm{tanh}G\right|+{\displaystyle _u_{}^{u_+}}Du{\displaystyle Dv\mathrm{tanh}G\mathrm{tanh}L},`$ (48)
$`\chi _{tr}`$ $`=`$ $`{\displaystyle \frac{\mathrm{exp}(u_+^2/2)}{J\sqrt{\pi pq^{p1}}}}{\displaystyle Dv(1\mathrm{tanh}L_v^{(+)})}+{\displaystyle \frac{\mathrm{exp}(u_{}^2/2)}{J\sqrt{\pi pq^{p1}}}}{\displaystyle Dv(1+\mathrm{tanh}L_v^{()})},`$ (49)
where
$`L`$ $`=`$ $`\widehat{\stackrel{~}{m}}+\sqrt{\widehat{\stackrel{~}{q}}}(\eta u+\sqrt{1\eta ^2}v)+{\displaystyle \frac{p}{2}}(p1)\stackrel{~}{\beta }J^2r^{p2}\chi _{tr}\mathrm{tanh}G,`$ (50)
$`L_v^{(\pm )}`$ $`=`$ $`\widehat{\stackrel{~}{m}}+\sqrt{\widehat{\stackrel{~}{q}}}(\eta u_\pm +\sqrt{1\eta ^2}v)\pm {\displaystyle \frac{p}{2}}(p1)\stackrel{~}{\beta }J^2r^{p2}\chi _{tr}\theta .`$ (51)
Eqs. (20-21), (46-49) for the order parameters $`m`$, $`q`$, $`\stackrel{~}{m}`$, $`\stackrel{~}{q}`$, $`r`$ and $`\chi _{tr}`$ form a close set of equations. The performance of selective freezing is measured by
$$M_{\mathrm{sf}}\frac{1}{N}\underset{i}{}\left[\mathrm{\Theta }(\theta ^2\sigma _i^2)\mathrm{sgn}\stackrel{~}{\sigma }_i+\mathrm{\Theta }(\sigma _i^2\theta ^2)\mathrm{sgn}\sigma _i\right].$$
(52)
From the above parameters, $`M_{\mathrm{sf}}`$ can be derived as:
$$M_{\mathrm{sf}}=\frac{1}{2}\mathrm{erf}\frac{u_+}{\sqrt{2}}\frac{1}{2}\mathrm{erf}\frac{u_{}}{\sqrt{2}}+_u_{}^{u_+}Du\mathrm{erf}\frac{L_u}{\sqrt{2\stackrel{~}{q}(1\eta ^2)}},$$
(53)
where $`L_u=\widehat{\stackrel{~}{m}}+\sqrt{\widehat{\stackrel{~}{q}}}\eta u+[p(p1)/2]\stackrel{~}{\beta }J^2r^{p2}\chi _{tr}\mathrm{tanh}G`$.
We have also tried to derive the above equations using the replica method. However, in the nearest results that we could find, terms involving the trans-susceptbility are absent, which we believe to be unphysical. Therefore the replica approach to the order parameter equations remain an open question.
We show an example of the case $`p=2`$ and $`j_0=J=1`$ in Fig. 1, where the overlap $`M_{\mathrm{sf}}`$ is plotted as a function of the decoding temperature $`T(=\beta ^1=\stackrel{~}{\beta }^1)`$ for various given values of freezing fraction $`f`$. When $`f=0`$ (no spins frozen) and $`f=1`$ (all spins frozen), the dynamics is equivalent to one with single stage, and the overlap reaches its maximum at the Nishimori point $`T=J^2/2j_0`$ as expected. We observe that the tolerance against variations in $`T`$ is enhanced by selective freezing for certain values of $`f`$.
It is therefore interesting to consider the appropriate values of $`f`$ for the best overlap at a given decoding temperature. Figs. 2(a-f) shows that at high temperatures such as in Figs. 2(a-c), there is a single maximum and its position is fairly independent of temperature, lying around $`f=0.9`$ in the present case. At intermediate temperatures such as in Figs. 2(d-e), there appear two maxima and as temperature changes, there is a discontinuous jump in the maximum position. Fig. 2(f) shows that when the temperature is lower than the Nishimori point ($`T_N=0.5`$), the overlap cannot be improved by selective freezing.
Figure 3 compares the overlap of the one-stage dynamics with that of the best of selective freezing. It shows that when the decoding temperature is mis-determined to be higher than its optimal value at the Nishimori point, selective freezing can provide a fairly robust performance. Furthermore, the choice of the freezing fraction for such robust performance appears to be quite independent of the temperature. The solid line in Fig. 4 locates the position for the best overlap and, as observed from Figs. 2(a-f), lies in the vicinity of $`f0.9`$ for a large range of temperature. The unshaded region in the same figure also indicates that selective freezing leads to an improvement in the overlap over a wide range of the parameter space.
We have also studied the dependence of the overlap on varying the freezing threshold $`\theta `$ rather than the freezing fraction $`f`$. However, Fig. 5 shows that the optimal value of $`\theta `$ has a much larger dependence on the temperature. This is due to the sensitive dependence of the thermal averages of the spins on temperature. At high temperatures, most spins are thermally agitated, and the freezing threshold has to be set to a very low value in order to freeze a given fraction of spins. On the other hand, at low temperatures, most spins are relatively stable, and the freezing threshold has to be set to a very high value in order to keep a given fraction of spins dynamic in the second stage. We conclude that the freezing fraction is a better controlling parameter for the decoding performance.
The advantages of selective freezing are confirmed by Monte Carlo simulations shown in Fig. 6. For one-stage dynamics, the overlap is maximum at the Nishimori point ($`T_N=0.5`$) as expected. However, it deterriorates rather rapidly when the decoding temperature increases. In contrast, selective freezing maintains a more steady performance, especially when $`f=0.9`$.
## IV THE MEAN-FIELD MODEL FOR IMAGE RESTORATION
In conventional image restoration problems, a given degraded image consists of the set of pixels $`\{\tau _i\}`$, but not the set of exchange interactions $`\{J_{i_1,\mathrm{},i_p}\}`$. On the other hand, effective restoration requires the introduction of a model prior distribution of the pixels for smooth images. In this case the Hamiltonian corresponds to that of a random field Ising model,
$$H\{\sigma \}=h\underset{i}{}\tau _i\sigma _i\frac{\beta _m}{z}\underset{ij}{}\sigma _i\sigma _j.$$
(54)
In mean-field systems, each pixel $`i`$ has an extensive valency. The pixels $`\tau _i`$ are the degraded versions of the source pixels $`\xi _i`$, corrupted by noise which, for convenience, is assumed to be Gaussian with mean $`a\xi _i`$ and variance $`\tau ^2`$, i.e.
$$P(\tau _i|\xi _i)=\frac{\mathrm{exp}\left[\frac{1}{2\tau ^2}(\tau _ia\xi _i)^2\right]}{\sqrt{2\pi \tau ^2}}.$$
(55)
In turn, the source pixels satisfy the prior distribution in Eq. (1). Applying the cavity argument for mean-field systems, the prior distribution becomes factorizable,
$$P(\xi _i)=\frac{\mathrm{exp}(\beta _sm_0\xi _i)}{2\mathrm{cosh}\beta _sm_0},$$
(56)
where $`m_0=\mathrm{tanh}\beta _sm_0`$. The order parameter in the first stage is given by
$$m\frac{1}{N}\underset{i}{}\sigma _i=\frac{1}{2\mathrm{cosh}\beta _sm_0}\underset{\xi =\pm 1}{}\mathrm{exp}(\beta _sm_0\xi )Dx\mathrm{tanh}U,$$
(57)
where $`U=\beta _mm+ha\xi +h\tau x`$. The overlap for the one-stage restoration process is given by
$$M\frac{1}{N}\underset{i}{}\xi _i\mathrm{sgn}\sigma _i=\frac{1}{2\mathrm{cosh}\beta _sm_0}\underset{\xi =\pm 1}{}\mathrm{exp}(\beta _sm_0\xi )\xi \mathrm{erf}\frac{\beta _mm+ha\xi }{\sqrt{2}h\tau }.$$
(58)
Next we consider selective freezing in the second stage with a freezing threshold $`\theta `$. The freezing fraction is given by
$$f\frac{1}{N}\underset{i}{}\mathrm{\Theta }\left(\sigma _i^2\theta ^2\right)=\frac{1}{2\mathrm{cosh}\beta _sm_0}\underset{\xi =\pm 1}{}\mathrm{exp}(\beta _sm_0\xi )\left[1\frac{1}{2}\mathrm{erf}\frac{u_+(\xi )}{\sqrt{2}}+\frac{1}{2}\mathrm{erf}\frac{u_{}(\xi )}{\sqrt{2}}\right],$$
(59)
where $`u_\pm (\xi )=(\pm u_0\beta _mmha\xi )/h\tau `$ with $`\mathrm{tanh}u_0=\theta `$. The order parameter of the second stage is given by
$`\stackrel{~}{m}`$ $``$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{i}{}}\left[\mathrm{\Theta }(\theta ^2\sigma _i^2)\stackrel{~}{\sigma }_i+\mathrm{\Theta }(\sigma _i^2\theta ^2)\mathrm{sgn}\sigma _i\right]`$ (60)
$`=`$ $`{\displaystyle \frac{1}{2\mathrm{cosh}\beta _sm_0}}{\displaystyle \underset{\xi =\pm 1}{}}\mathrm{exp}(\beta _sm_0\xi )\left[{\displaystyle \frac{1}{2}}\mathrm{erf}{\displaystyle \frac{u_+(\xi )}{\sqrt{2}}}{\displaystyle \frac{1}{2}}\mathrm{erf}{\displaystyle \frac{u_{}(\xi )}{\sqrt{2}}}+{\displaystyle _{u_{}(\xi )}^{u_+(\xi )}}Dx\mathrm{tanh}L\right],`$ (61)
where $`L=\beta _m\stackrel{~}{m}+ha\xi +h\tau x`$. The overlap for selective freezing is given by
$`M_{\mathrm{sf}}`$ $``$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{i}{}}\xi _i\left[\mathrm{\Theta }(\theta ^2\sigma _i^2)\mathrm{sgn}\stackrel{~}{\sigma }_i+\mathrm{\Theta }(\sigma _i^2\theta ^2)\mathrm{sgn}\sigma _i\right]`$ (62)
$`=`$ $`{\displaystyle \frac{1}{2\mathrm{cosh}\beta _sm_0}}{\displaystyle \underset{\xi =\pm 1}{}}\mathrm{exp}(\beta _sm_0\xi )\xi \mathrm{erf}{\displaystyle \frac{g(\beta _m\stackrel{~}{m})+ha\xi }{\sqrt{2}h\tau }},`$ (63)
where
$$g(\beta _m\stackrel{~}{m})=\{\begin{array}{ccccc}\beta _mmu_0& & & & \beta _m\stackrel{~}{m}<\beta _mmu_0,\\ \beta _m\stackrel{~}{m}& & & & \beta _mmu_0<\beta _m\stackrel{~}{m}<\beta _mm+u_0,\\ \beta _mm+u_0& & & & \beta _m\stackrel{~}{m}>\beta _mm+u_0.\end{array}$$
(64)
We note that since the spin-glass interaction is absent in this case, there are no trans-susceptibility effects. This is unlike the case of error-correcting codes, in which $`\chi _{tr}`$ is nonzero when $`J`$ is nonzero.
The three cases of the function $`g(\beta _m\stackrel{~}{m})`$ in Eq. (64) correspond to three situations. When $`\beta _m\stackrel{~}{m}<\beta _mmu_0`$, all the dynamic spins in the second stage have negative thermodynamic averages and therefore take the value $`1`$ in the two-stage restoration process. This is equivalent to a one-stage restoration process in which all spins with thermodynamic averages above the threshold $`+\theta `$ are frozen to $`+1`$, and to $`1`$ otherwise. Similarly, when $`\beta _m\stackrel{~}{m}>\beta _mm+u_0`$, all the dynamic spins in the second stage have positive thermodynamic averages. Only when $`\beta _mmu_0<\beta _m\stackrel{~}{m}<\beta _mm+u_0`$, do we have the dynamic spins frozen to partly $`+1`$ and partly $`1`$.
We can consider the condition for the optimal performance $`M_{\mathrm{sf}}`$ of selective freezing. For a given distribution of data and noise, $`g(\beta _m\stackrel{~}{m})`$ is the only adjustable parameter in Eq. (63), playing the same role as the adjustable parameter $`\beta _mm`$ for one-stage dynamics in Eq. (58). In the space of $`h`$ and $`\beta _m`$, the performance is optimal along the line $`h/\beta _\tau =\beta _mm/\beta _sm_0`$ for one-stage dynamics ($`\beta _\tau =a/\tau ^2`$ for Gaussian noise). Analogously, there exists a line of optimal performance defined by $`h/\beta _\tau =g(\beta _m\stackrel{~}{m})/\beta _sm_0`$ for selective freezing.
An example of the lines of optimal performance is shown in Fig. 7. It is interesting to note the kinks for certain freezing fractions. They correspond to transitions of cases in which the dynamic spins are partially or completely frozen to $`\pm 1`$.
A comparison of Eqs. (58) and (63) shows that selective freezing performs as well as one-stage dynamics, but cannot outperform it. Nevertheless, selective freezing provides a rather stable performance when the hyperparameters cannot be estimated precisely. In image restoration, the usual practice is to choose a fixed ratio of $`\beta _m/h`$. Fig. 8 confirms this stability along the line of operation with $`\beta _m/h`$ set to the optimal ratio $`\beta _s/\beta _\tau `$. Note especially that the lines with $`f=0.7`$ and $`0.9`$ attain a nearly optimal value of $`M_{\mathrm{sf}}`$ over a wide range of parameters. The kink at $`f=0.9`$ is, again, due to the appearance of the $`1`$ frozen dynamic spins (to the right of the kink).
The stable performance of selective freezing can be partly explained by the proximity of the lines of optimal performance with the line of operation which, as discussed in , is an important factor in hyperparameter estimation. This is illustrated by the optimal lines for small values of $`f`$ near the Nishimori point $`(T_m,h)=(1.05^1,1)`$ in Fig. 7.
However, the advantage of selective freezing does not only rely on the fortuitous combination of parameters. Even when the parameters are not chosen optimally, selective freezing still maintains a rather robust performance. For example, along the line of optimal performance for $`f=0.9`$ in Fig. 7, the bending at the kink only causes a modest reduction in the overlap $`M_{\mathrm{sf}}`$ in Fig. 8.
To study the robustness of the performance of selective freezing, we model a situation common in modern communication channels carrying multimedia traffic, which are often bursty in nature. Since burstiness results in intermittent interferences, we consider a noise with two Gaussian components, each with its own characteristics. A random fraction $`f_1`$ of the pixels are influenced by Gaussian noise with signal strength $`a_1`$ and noise variance $`\tau _1^2`$. The rest of the pixels have strength $`a_2`$ and noise variance $`\tau _2^2`$. Hence the distribution of the degraded pixels are
$$P(\tau _i|\xi _i)=f_1\frac{\mathrm{exp}\left[\frac{1}{2\tau _1^2}(\tau _ia_1\xi _i)^2\right]}{\sqrt{2\pi \tau _1^2}}+f_2\frac{\mathrm{exp}\left[\frac{1}{2\tau _2^2}(\tau _ia_2\xi _i)^2\right]}{\sqrt{2\pi \tau _2^2}},$$
(65)
where $`f_2=1f_1`$. The equations for the order parameters can be generalized from the single component case in a straightforward manner.
A case of interest is that the restoration agent operates on the assumption of the characteristics of the majority component of the channel, say the first component. Hence it operates at the ratio $`\beta _m/h=\beta _s\tau _1^2/a_1`$. Suppose the Gaussian noise is partly interrupted to take the characteristics of the second component, but the operation parameters cannot be adjusted soon enough, then there will be a degradation of the quality of the restored images. In the example in Fig. 9, the reduction of the overlap $`M_{\mathrm{sf}}`$ for selective freezing is much more modest than the one-stage process ($`f=0`$).
An alternative situation is that the restoration agent is able to detect the changes in the average signal strengths and noise variance, but still operates on the assumption of a single-component Gaussian channel. Suppose that such simple statistics as $`\mathrm{sgn}\tau _i`$, $`\tau _i`$ and $`\tau _i^2`$ are accessible. Then the parameters $`m_0^{}`$, $`a^{}`$ and $`\tau ^{}`$ estimated by the restoration agent are obtained, for $`\tau _1=\tau _2=\tau `$, from the solutions of
$`m_0^{}\mathrm{erf}{\displaystyle \frac{a^{}}{\sqrt{2}\tau ^{}}}`$ $`=`$ $`\mathrm{sgn}\tau _i=m_0\left[f_1\mathrm{erf}{\displaystyle \frac{a_1}{\sqrt{2}\tau _1}}+f_2\mathrm{erf}{\displaystyle \frac{a_2}{\sqrt{2}\tau _2}}\right],`$ (66)
$`m_0^{}a^{}`$ $`=`$ $`\tau _i=m_0\left[f_1a_1+f_2a_2\right],`$ (67)
$`a^2+\tau ^2`$ $`=`$ $`\tau _i^2=f_1(a_1^2+\tau _1^2)+f_2(a_2^2+\tau _2^2),`$ (68)
and $`\beta _s^{}=\mathrm{tanh}^1m_0^{}/m_0^{}`$. Using these estimated parameters, the performances in Fig. 10 improve over their counterparts based on only the majority component in Fig. 9. Still, one-stage restoration cannot avoid the performance drop when $`h`$ vanishes, whereas correspondingly, selective freezing has a much more gentle drop in performance.
It is interesting to study the more realistic case of two-dimensional images, since we have so far presented analytical results for the mean field model only. As confirmed by the results for Monte carlo simulations in Fig. 11, the overlaps of selective freezing are much more steadier than that of the one-stage dynamics when the decoding temperature changes. This steadiness is most remarkable for a freezing fraction of $`f=0.9`$.
## V DISCUSSIONS
We have introduced a multistage technique for error-correcting codes and image restoration, in which the information extracted from the former stage can be used selectively to improve the performance of the latter one. While the overlap $`M_{\mathrm{sf}}`$ of the selective freezing is bounded by the optimal performance of the one-stage dynamics derived in , it has the advantage of being tolerant to uncertainties in hyperparameter estimation. The performance is especially steady when the fraction of frozen spins, rather than the threshold of their thermodynamic averages, is fixed in the process. This is confirmed by both analytical and simulational results for mean-field and finite-dimensional models. As an example, we have illustrated its advantage of robustness when the noise distribution is composed of more than one Gaussian components, such as in the case of modern communication channels supporting multimedia applications.
We found that selective freezing is most useful when more than one hyperparameters have to be estimated, as illustrated by the example of image restoration, where both $`\beta _m`$ and $`h`$ have to be estimated. In the example of error-correcting codes discussed in Section III, there is only one hyperparameter $`T_m`$, and it is found that selective freezing has performance advantages only when $`T_m`$ is chosen above the Nishimori point. However, more than one hyperparameters are often present in practical applications.
Selective freezing can be generalized to more than two stages, in which spins that remain relatively stable in one stage are progressively frozen in the following one. It is expected that the performance can be even more robust.
While the multistage process described here has a robust performance, it does not raise the critical temperature or the critical noise level for the existence of the ordered phase. Nor can it widen the basin of attraction for the ordered phase. Other multistage processes, proposed in for neural networks, may be able to achieve this. This remains an area for further research.
We have made progress in the theoretical treatment of multistage processes using the cavity method. It allows the thermal averages of spins to be expressed in terms of the cavity fields. Since a cavity field is uncorrelated with the spin in consideration, it can in turn be expressed in terms of the means and covariances of the spin averages, thereby arriving at a set of self-consistent equations for the order parameters. In particular, there appears a trans-susceptibility term since variations of the cavity field in the first stage are correlated with the spin average in the second stage due to the selective nature of the freezing process in the second stage. However, for the ordered phase considered in this paper, the effects of the trans-susceptibility term is not too large except near the phase boundary.
On the other hand, we have a remark about the basic assumption of the cavity method, namely that the addition or removal of a spin causes a small change in the system describable by a perturbative approach. In fact, adding or removing a spin may cause the thermal averages of other spins to change from below to above the thresholds $`\pm \theta `$ (or vice versa). This change, though often small, induces a non-negligible change of the thermal averages from fractional values to the frozen values of $`\pm 1`$ (or vice versa) in the second stage. The perturbative analysis of these changes is only approximate. The situation is reminiscent of similar instabilities in other disordered systems such as the perceptron, and are equivalent to Almeida-Thouless instabilities in the replica method . A full treatment of the problem would require the introduction of a rough energy landscape , or the replica symmetry breaking ansatz in the replica method . Nevertheless, previous experiences on disordered systems showed that the corrections made by a more complete treatment may not be too large in the ordered phase. For example, corresponding analytical and simulational results in Figs. 1 and 6 respectively are close to each other.
In practical implementations of error-correcting codes, algorithms based on belief-propagation methods, rather than Monte Carlo methods, are often employed . It has recently been shown that such decoded messages converge to the solutions of the TAP equations in the corresponding thermodynamic system . Again, the performance of these algorithms are sensitive to the estimation of hyperparameters. We propose that the selective freezing procedure has the potential to make these algorithms more robust.
Incidentally, multistage dynamics has also been applied in the recently popular turbo codes . Messages are coded in sequences with two possible permutations and at each iterative stage, the information derived from decoding one sequence is fed to the other in the form of external fields for each bit. The techniques developed in the present context can be used to study this iterative process.
###### Acknowledgements.
KYMW wishes to thank Tokyo Institute of Technology for hospitality. HN is grateful to Hong Kong University of Science and Technology for hospitality. This work was partially supported by research grant HKUST6157/99P of the Research Grant Council of Hong Kong.
## A Thermal averages of spins
In this appendix we derive Eq. (16) starting from the clustering property Eq. (15). For convenience we illustrate the derivation for $`p=2`$. We separate the Hamiltonian into two parts, one does not contain $`\sigma _1`$ and the other does. Hence
$$H=H^{\backslash 1}\beta \underset{j>1}{}J_{1j}\sigma _1\sigma _j.$$
(A1)
Thus the thermal average can be written as
$$\sigma _1=\frac{\mathrm{Tr}^{\backslash 1}e^{H^{\backslash 1}}\mathrm{Tr}_1\sigma _1\mathrm{exp}\left(\beta \sigma _1\underset{j}{}J_{1j}\sigma _j\right)/\mathrm{Tr}^{\backslash 1}e^{H^{\backslash 1}}}{\mathrm{Tr}^{\backslash 1}e^{H^{\backslash 1}}\mathrm{Tr}_1\mathrm{exp}\left(\beta \sigma _1_jJ_{1j}\sigma _j\right)/\mathrm{Tr}^{\backslash 1}e^{H^{\backslash 1}}}.$$
(A2)
Expanding the exponential function in the denominator and tracing over $`\sigma `$, we get
$$\mathrm{Den}.=2\underset{n\mathrm{even}}{}\frac{\beta ^n}{n!}\underset{j_1\mathrm{}j_n}{}J_{1j_1}\mathrm{}J_{1j_n}\sigma _{j_1}\mathrm{}\sigma _{j_n}^{\backslash 1}.$$
(A3)
Next, we use the clustering property to factorize the thermal average $`\sigma _{j_1}\mathrm{}\sigma _{j_n}^{\backslash 1}`$. For the coupling distribution specified by Eq. (14), only two kinds of contributions are significant in the summation over the indices $`j_1\mathrm{}j_n`$. In the first kind, an index $`j`$ remains distinct from the rest, contributing a factor of $`J_{1j}\sigma _j^{\backslash 1}`$. In the second kind, two indices become paired up. However, when $`j`$ and $`k`$ pair up, the thermal average $`\sigma _j\sigma _k^{\backslash 1}`$ becomes 1 instead of $`(\sigma _j^{\backslash 1})^2`$. Hence the additional contribution due to the pairing is $`J_{1j}^2[1(\sigma _j^{\backslash 1})^2]`$. Other than these, the contributions due to the pairing of three or more indices are smaller by factors of $`N`$.
The denominator can now be considered a summation over $`n`$ and $`m`$, which are respectively the total number of indices and the number of pairs of paired indices appearing in a term. The number of such terms is $`n!/m!2^m(n2m)!`$. Hence
$$\mathrm{Den}.=2\underset{n\mathrm{even}}{}\underset{m=0}{\overset{n/2}{}}\frac{\beta ^n}{n!}\frac{n!}{m!2^m(n2m)!}\left[\underset{j}{}J_{1j}\sigma _j^{\backslash 1}\right]^{n2m}\left\{\underset{j}{}J_{1j}^2[1(\sigma _j^{\backslash 1})^2]\right\}^m,$$
(A4)
which can be simplified to
$$\mathrm{Den}.=2\mathrm{exp}\left\{\frac{1}{2}\beta ^2\underset{j}{}J_{1j}^2[1(\sigma _j^{\backslash 1})^2]\right\}\mathrm{cosh}\left\{\beta \underset{j}{}J_{1j}\sigma _j^{\backslash 1}\right\}.$$
(A5)
Similarly, the numerator can be written as
$$\mathrm{Num}.=2\mathrm{exp}\left\{\frac{1}{2}\beta ^2\underset{j}{}J_{1j}^2[1(\sigma _j^{\backslash 1})^2]\right\}\mathrm{sinh}\left\{\beta \underset{j}{}J_{1j}\sigma _j^{\backslash 1}\right\}.$$
(A6)
Substituting Eq. (A5) and (A6) into Eq. (A2), we arrive at Eq. (16).
## B Change in thermal averages on removal of a spin
In this appendix we derive Eq. (32). For convenience we illustrate the derivation for $`p=2`$. We separate the Hamiltonian into four parts: (a) does not contain spins 1 and $`j`$, (b) contains only spins 1 and $`j`$, (c) contains spin 1 but not $`j`$, (d) contains spin $`j`$ but not 1. This yields
$$H=H^{\backslash 1j}\beta J_{1j}\sigma _1\sigma _j\beta \underset{k1j}{}J_{k1}\sigma _k\sigma _1\beta \underset{k1j}{}J_{kj}\sigma _k\sigma _j.$$
(B1)
The thermal average of $`\sigma _j`$ can then be written as
$$\sigma _j=\frac{\mathrm{Tr}_{1j}\mathrm{Tr}^{\backslash 1j}e^H\sigma _j/\mathrm{Tr}^{\backslash 1j}e^{H^{\backslash 1j}}}{\mathrm{Tr}_{1j}\mathrm{Tr}^{\backslash 1j}e^H/\mathrm{Tr}^{\backslash 1j}e^{H^{\backslash 1j}}}.$$
(B2)
Using the mean-field technique developed in Appendix A, the denominator can be written as
$`\mathrm{Den}.=\mathrm{Tr}_{1j}\mathrm{exp}\{\beta J_{1j}\sigma _1\sigma _j`$ $`+`$ $`\beta {\displaystyle \underset{k1j}{}}\sigma _k^{\backslash 1j}(J_{k1}\sigma _1+J_{kj}\sigma _j)`$ (B3)
$`+`$ $`{\displaystyle \frac{1}{2}}\beta ^2{\displaystyle \underset{k1j}{}}[1(\sigma _k^{\backslash 1j})^2](J_{k1}\sigma _1+J_{kj}\sigma _j)^2\}.`$ (B4)
After collecting terms and discarding negligible ones,
$$\mathrm{Den}.=\mathrm{Tr}_{1j}\mathrm{exp}\{\beta \sigma _1\underset{k1j}{}J_{1k}\sigma _k^{\backslash 1j}+\beta \sigma _j\underset{k1j}{}J_{jk}\sigma _k^{\backslash 1j}+\beta J_{1j}\sigma _1\sigma _j+\beta ^2(1q)J^2\}.$$
(B5)
Together with a similar manipulation of the numerator, we obtain
$$\sigma _j=\mathrm{tanh}\beta \left(h_j^{\backslash 1}+J_{j1}\mathrm{tanh}\beta h_1^{\backslash j}\right),$$
(B6)
whose Taylor expansion yields
$$\sigma _j=\sigma _j^{\backslash 1}+\left(\beta \mathrm{sech}^2\beta h_j^{\backslash 1}\right)\left(J_{j1}\mathrm{tanh}\beta h_1^{\backslash j}\right),$$
(B7)
which becomes Eq. (32) for the case $`p=2`$.
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# The polynomial property (V)
## Abstract.
Given Banach spaces $`E`$ and $`F`$, we denote by $`𝒫(^kE,F)`$ the space of all $`k`$-homogeneous (continuous) polynomials from $`E`$ into $`F`$, and by $`𝒫_{wb}(^kE,F)`$ the subspace of polynomials which are weak-to-norm continuous on bounded sets. It is shown that if $`E`$ has an unconditional finite dimensional expansion of the identity, the following assertions are equivalent: (a) $`𝒫(^kE,F)=𝒫_{wb}(^kE,F)`$; (b) $`𝒫_{wb}(^kE,F)`$ contains no copy of $`c_0`$; (c) $`𝒫(^kE,F)`$ contains no copy of $`\mathrm{}_{\mathrm{}}`$; (d) $`𝒫_{wb}(^kE,F)`$ is complemented in $`𝒫(^kE,F)`$. This result was obtained by Kalton for linear operators. As an application, we show that if $`E`$ has Pełczyński’s property (V) and satisfies $`𝒫(^kE)=𝒫_{wb}(^kE)`$ then, for all $`F`$, every unconditionally converging $`P𝒫(^kE,F)`$ is weakly compact. If $`E`$ has an unconditional finite dimensional expansion of the identity, then the converse is also true.
###### Key words and phrases:
Weakly continuous polynomial, unconditionally converging polynomial, weakly compact polynomial, weakly unconditionally Cauchy series, property (V) The first named author was supported in part by DGICYT Grant PB 97–0349 (Spain) The second named author was supported in part by DGICYT Grant PB 96–0607 (Spain) file ppv.tex
Given two Banach spaces $`E`$ and $`F`$, we denote by $`𝒫(^kE,F)`$ the space of all $`k`$-homogeneous (continuous) polynomials from $`E`$ into $`F`$, and by $`𝒫_{wb}(^kE,F)`$ the subspace of polynomials which are weak-to-norm continuous on bounded sets. This subspace has been studied by many authors: see, for instance, . Clearly, every polynomial in $`𝒫_{wb}(^kE,F)`$ takes bounded sets into relatively compact sets. Observe that $`𝒫(^1E,F)=(E,F)`$, the space of (linear bounded) operators from $`E`$ into $`F`$, and that $`𝒫_{wb}(^1E,F)=𝒦(E,F)`$, the space of compact operators. The spaces $`𝒫(^0E,F)`$ and $`𝒫_{wb}(^0E,F)`$ may be identified with $`F`$.
Kalton studied in the structure of the space $`𝒦(E,F)`$. In the present paper, we obtain versions of his results for the space $`𝒫_{wb}(^kE,F)`$, showing that $`𝒫_{wb}(^kE,F)`$ contains a copy of $`\mathrm{}_{\mathrm{}}`$ if and only if either $`E`$ contains a complemented copy of $`\mathrm{}_1`$ or $`F`$ contains a copy of $`\mathrm{}_{\mathrm{}}`$. We also prove that, for $`E`$ having an unconditional finite dimensional expansion of the identity, the following assertions are equivalent: (a) $`𝒫(^kE,F)=𝒫_{wb}(^kE,F)`$; (b) $`𝒫_{wb}(^kE,F)`$ contains no copy of $`c_0`$; (c) $`𝒫(^kE,F)`$ contains no copy of $`\mathrm{}_{\mathrm{}}`$; (d) $`𝒫_{wb}(^kE,F)`$ is complemented in $`𝒫(^kE,F)`$.
As an application, we prove that, if $`E`$ has property (V) (the definitions are given below) and $`𝒫(^kE)=𝒫_{wb}(^kE)`$, then every $`k`$-homogeneous unconditionally converging polynomial on $`E`$ is weakly compact. If $`E`$ has an unconditional finite dimensional expansion of the identity, then the converse is also true.
Throughout, $`E`$ and $`F`$ will denote Banach spaces, $`B_E`$ is the closed unit ball of $`E`$ and $`S_E`$ is the unit sphere of $`E`$; $`E^{}`$ will be the dual of $`E`$. The set of natural numbers is denoted by $``$. As usual, $`(e_n)`$ stands for the unit vector basis of $`c_0`$.
A formal series $`x_n`$ in $`E`$ is weakly unconditionally Cauchy (w.u.C., for short) if, for every $`\varphi E^{}`$, we have $`|\varphi (x_n)|<+\mathrm{}`$. Equivalent definitions may be seen in \[7, Theorem V.6\]. The series is unconditionally convergent if every subseries converges. Equivalent definitions may be seen in \[8, Theorem 1.9\].
A polynomial $`P𝒫(^kE,F)`$ is unconditionally converging if, for each w.u.C. series $`x_n`$ in $`E`$, the sequence $`\left(P\left(_{i=1}^nx_i\right)\right)_n`$ is convergent in $`F`$. The space of all unconditionally converging polynomials is denoted by $`𝒫_{uc}(^kE,F)`$. This class has been very useful for obtaining polynomial characterizations of Banach space properties (see ). We say that $`P𝒫(^kE,F)`$ is (weakly) compact if $`P(B_E)`$ is relatively (weakly) compact in $`F`$. Every weakly compact polynomial is unconditionally converging. For the general theory of polynomials on Banach spaces, we refer to .
To each polynomial $`P𝒫(^kE,F)`$ we can associate a unique symmetric $`k`$-linear mapping $`\widehat{P}:E\times \stackrel{(k)}{\mathrm{}}\times EF`$ so that $`P(x)=\widehat{P}(x,\mathrm{},x)`$ and an operator $`T_P:E𝒫(^{k1}E,F)`$ given by $`T_P(x)(y):=\widehat{P}(x,y,\stackrel{(k1)}{\mathrm{}},y)`$. It is well known that $`P𝒫_{wb}(^kE,F)`$ if and only if $`T_P`$ is compact \[3, Theorem 2.9\].
Denote by $`_s(E,𝒫(^{k1}E,F))`$ the space of all operators $`C:E𝒫(^{k1}E,F)`$ such that
$$(C(x))^{}(y_1,\mathrm{},y_{k1})=(C(y_1))^{}(x,y_2,\mathrm{},y_{k1})(x,y_1,\mathrm{},y_{k1}E),$$
where $`(C(x))^{}`$ stands for the symmetric $`(k1)`$-linear mapping associated to $`C(x)`$.
###### Proposition 1.
The mapping $`T:𝒫(^kE,F)_s(E,𝒫(^{k1}E,F))`$ given by $`T(P)=T_P`$ is a surjective linear isomorphism.
Proof. Clearly, $`T`$ is well defined, linear and injective. Since
$$T_P=\underset{xB_E}{sup}T_P(x)=\underset{x,yB_E}{sup}\widehat{P}(x,y,\stackrel{(k1)}{\mathrm{}},y)\widehat{P}\frac{k^k}{k!}P$$
\[26, Theorem 2.2\], we have that $`T`$ is continuous. To see that it is surjective, take $`C_s(E,𝒫(^{k1}E,F))`$, and define $`A:E\times \stackrel{(k)}{\mathrm{}}\times EF`$ by $`A(y_1,\mathrm{},y_k):=(C(y_1))^{}(y_2,\mathrm{},y_k)`$, and $`P(x)=A(x,\stackrel{(k)}{\mathrm{}},x)`$. Then
$$T_P(x)(y)=A(x,y,\stackrel{(k1)}{\mathrm{}},y)=C(x)(y)(x,yE).$$
Hence, $`T_P=C`$. $`\mathrm{}`$
The subspace of all operators in $`_s(E,𝒫(^{k1}E,F))`$ which are compact (resp. weakly compact) will be denoted by $`𝒦_s(E,𝒫(^{k1}E,F))`$ (resp. $`𝒲_s(E,𝒫(^{k1}E,F))`$). Given $`C𝒦_s(E,𝒫(^{k1}E,F))`$, the symmetry of $`C`$ easily implies that $`C(E)𝒫_{wb}(^{k1}E,F)`$.
The proof of \[5, Proposition 5.3\] gives:
###### Proposition 2.
For $`nm`$, the space $`𝒫(^mE,F)`$ (resp. $`𝒫_{wb}(^mE,F)`$) is isomorphic to a complemented subspace of $`𝒫(^nE,F)`$ (resp. $`𝒫_{wb}(^nE,F)`$).
###### Theorem 3.
The space $`𝒫_{wb}(^kE,F)`$ contains a copy of $`\mathrm{}_{\mathrm{}}`$ if and only if either $`F`$ contains a copy of $`\mathrm{}_{\mathrm{}}`$ or $`E`$ contains a complemented copy of $`\mathrm{}_1`$.
Proof. If $`𝒫_{wb}(^kE,F)`$ contains a copy of $`\mathrm{}_{\mathrm{}}`$, a fortiori the space $`𝒦(E,𝒫_{wb}(^{k1}E,F))`$ contains it. Therefore \[23, Theorem 4\], either $`E`$ contains a complemented copy of $`\mathrm{}_1`$ or $`𝒫_{wb}(^{k1}E,F)`$ contains a copy of $`\mathrm{}_{\mathrm{}}`$. Repeating the process, we conclude that either $`E`$ contains a complemented copy of $`\mathrm{}_1`$ or $`𝒫_{wb}(^0E,F)F`$ contains a copy of $`\mathrm{}_{\mathrm{}}`$.
Conversely, if $`F`$ contains a copy of $`\mathrm{}_{\mathrm{}}`$, since $`F𝒫_{wb}(^0E,F)`$ is isomorphic to a subspace of $`𝒫_{wb}(^kE,F)`$ (Proposition 2), we obtain that $`𝒫_{wb}(^kE,F)`$ contains a copy of $`\mathrm{}_{\mathrm{}}`$. If $`E`$ contains a complemented copy of $`\mathrm{}_1`$, then $`E^{}`$ contains a copy of $`\mathrm{}_{\mathrm{}}`$; since $`E^{}=𝒫_{wb}(^1E)`$ is isomorphic to a subspace of $`𝒫_{wb}(^1E,F)`$, which is in turn isomorphic to a subspace of $`𝒫_{wb}(^kE,F)`$, we obtain that the latter contains a copy of $`\mathrm{}_{\mathrm{}}`$. $`\mathrm{}`$
The proof of \[23, Lemma 2\] yields:
###### Lemma 4.
Assume $`E`$ is separable, $`𝒫_{wb}(^kE,F)`$ is complemented in $`𝒫(^kE,F)`$, and an operator $`\mathrm{\Phi }:\mathrm{}_{\mathrm{}}𝒫(^kE,F)`$ is given with the following properties:
(a) $`\mathrm{\Phi }(e_n)𝒫_{wb}(^kE,F)`$ for all $`n`$;
(b) the subset $`\{\mathrm{\Phi }(\xi )(x):\xi \mathrm{}_{\mathrm{}},xE\}F`$ is separable.
Then, for every infinite subset $`M`$, there exists an infinite subset $`M_0M`$ with $`\mathrm{\Phi }(\xi )𝒫_{wb}(^kE,F)`$ for all $`\xi \mathrm{}_{\mathrm{}}(M_0)`$.
###### Lemma 5.
Suppose $`E`$ contains a complemented copy of $`\mathrm{}_1`$. Then $`𝒫_{wb}(^kE,F)`$ is uncomplemented in $`𝒫(^kE,F)`$ for every $`F`$ and $`k`$ $`(k>1)`$.
Proof. As in \[23, Lemma 3\], we can reduce the problem to the case $`E=\mathrm{}_1`$.
Fix $`vS_F`$ and define the operator
$$\mathrm{\Phi }:\mathrm{}_{\mathrm{}}𝒫(^k\mathrm{}_1,F)$$
by
$$\mathrm{\Phi }(\xi )(x)=\underset{i=1}{\overset{\mathrm{}}{}}\xi _ix_i^kv\text{for }\xi =(\xi _i)_{i=1}^{\mathrm{}}\mathrm{}_{\mathrm{}}\text{ and }x=(x_i)_{i=1}^{\mathrm{}}\mathrm{}_1.$$
Since
$$\mathrm{\Phi }(\xi )=\underset{xB_\mathrm{}_1}{sup}\underset{i=1}{\overset{\mathrm{}}{}}\xi _ix_i^kv\xi \underset{xB_\mathrm{}_1}{sup}\underset{i=1}{\overset{\mathrm{}}{}}\left|x_i^k\right|=\xi ,$$
$`\mathrm{\Phi }`$ is continuous (easily, it is even an isometric embedding).
We claim that $`\mathrm{\Phi }(\xi )𝒫_{wb}(^k\mathrm{}_1,F)`$ if and only if $`\xi c_0`$. Indeed, let
$$T_{\mathrm{\Phi }(\xi )}:\mathrm{}_1𝒫(^{k1}\mathrm{}_1,F)$$
be the associated operator given by
$$T_{\mathrm{\Phi }(\xi )}(x)(y)=\underset{i=1}{\overset{\mathrm{}}{}}\xi _ix_iy_i^{k1}v\text{for }x=(x_i),y=(y_i)\mathrm{}_1.$$
Since, for $`nm`$,
$`\left(T_{\mathrm{\Phi }(\xi )}(e_n)T_{\mathrm{\Phi }(\xi )}(e_m)\right)(y)=`$
$`T_{\mathrm{\Phi }(\xi )}(e_ne_m)(y)={\displaystyle \underset{i=1}{\overset{\mathrm{}}{}}}\xi _i(\delta _{in}\delta _{im})y_i^{k1}v=\left|\xi _ny_n^{k1}\xi _my_m^{k1}\right|,`$
we have
$$\mathrm{max}\{|\xi _n|,|\xi _m|\}T_{\mathrm{\Phi }(\xi )}(e_n)T_{\mathrm{\Phi }(\xi )}(e_m)|\xi _n|+|\xi _m|,$$
it follows that $`T_{\mathrm{\Phi }(\xi )}`$ is compact if and only if $`\xi c_0`$ (see \[7, Exercise VII.5\]).
By Lemma 4, there is an infinite $`M`$ such that $`\mathrm{\Phi }(\xi )𝒫_{wb}(^k\mathrm{}_1,F)`$ for all $`\xi \mathrm{}_{\mathrm{}}(M)`$, which contradicts the above claim. $`\mathrm{}`$
In the linear case $`(k=1)`$, Kalton needs to assume that $`F`$ is infinite dimensional for the validity of Lemma 5 \[23, Lemma 3\]. Taking $`k>1`$, we can drop this condition. Observe that, if $`k=1`$ and $`dim(F)<\mathrm{}`$, we have
$$𝒫_{wb}(^1E,F)=𝒦(E,F)=(E,F)=𝒫(^1E,F),$$
so the conclusion of the Lemma is not true.
An unconditional finite dimensional expansion of the identity for a Banach space $`E`$ is a sequence of finite dimensional operators $`A_n:EE`$ such that for each $`xE`$,
$$x=\underset{n=1}{\overset{\mathrm{}}{}}A_n(x)$$
unconditionally. This condition is a bit more general than having an unconditional basis .
###### Lemma 6.
Suppose $`E`$ has an unconditional finite dimensional expansion of the identity and let $`P𝒫(^kE,F)`$. Then there is a w.u.C. series $`P_i`$ in $`𝒫_{wb}(^kE,F)`$ such that, for all $`xE`$, $`P(x)=_{m=1}^{\mathrm{}}P_m(x)`$ unconditionally.
Proof. There is a sequence $`(A_i)𝒦(E,E)`$ such that, for every $`xE`$, we have $`x=_{i=1}^{\mathrm{}}A_i(x)`$ unconditionally. Then,
$`P\left({\displaystyle \underset{i=1}{\overset{n}{}}}A_i(x)\right)`$ $`=`$ $`{\displaystyle \underset{i_1,\mathrm{},i_k=1}{\overset{n}{}}}\widehat{P}(A_{i_1}(x),\mathrm{},A_{i_k}(x))`$
$`=`$ $`{\displaystyle \underset{m=1}{\overset{n}{}}}\left({\displaystyle \underset{\mathrm{max}\{i_1,\mathrm{},i_k\}=m}{}}\widehat{P}(A_{i_1}(x),\mathrm{},A_{i_k}(x))\right)`$
$`=`$ $`{\displaystyle \underset{m=1}{\overset{n}{}}}P_m(x)`$
where $`P_m𝒫_{wb}(^kE,F)`$.
Choosing finite subsets $`I_1,\mathrm{},I_k`$ of integers, we have
$`{\displaystyle \underset{i_1I_1,\mathrm{},i_kI_k}{}}\widehat{P}(A_{i_1}(x),\mathrm{},A_{i_k}(x))`$ $`=`$ $`\widehat{P}({\displaystyle \underset{i_1I_1}{}}A_{i_1}(x),\mathrm{},{\displaystyle \underset{i_kI_k}{}}A_{i_k}(x))`$
$``$ $`\widehat{P}{\displaystyle \underset{i_1I_1}{}}A_{i_1}(x)\mathrm{}{\displaystyle \underset{i_kI_k}{}}A_{i_k}(x).`$
Hence, the series
$$\underset{i_1,\mathrm{},i_k=1}{\overset{\mathrm{}}{}}\widehat{P}(A_{i_1}(x),\mathrm{},A_{i_k}(x))$$
is unconditionally convergent for all $`xE`$ \[8, Theorem 1.9\]. Therefore, $`P(x)=_{m=1}^{\mathrm{}}P_m(x)`$ unconditionally.
Moreover, by the uniform boundedness principle \[26, Theorem 2.6\], we have
$$\underset{Ifinite}{sup}\underset{iI}{}P_i<\mathrm{}.$$
So $`P_i`$ is w.u.C. in $`𝒫_{wb}(^kE,F)`$ \[7, Theorem V.6\]. $`\mathrm{}`$
###### Theorem 7.
Suppose $`E`$ has an unconditional finite dimensional expansion of the identity and let $`k`$ $`(k>1)`$. Then the following assertions are equivalent:
(a) $`𝒫(^kE,F)=𝒫_{wb}(^kE,F)`$;
(b) $`𝒫_{wb}(^kE,F)`$ contains no copy of $`c_0`$;
(c) $`𝒫(^kE,F)`$ contains no copy of $`\mathrm{}_{\mathrm{}}`$;
(d) $`𝒫_{wb}(^kE,F)`$ is complemented in $`𝒫(^kE,F)`$.
Proof. (a) $``$ (c): Assume $`𝒫(^kE,F)=𝒫_{wb}(^kE,F)`$ contains a copy of $`\mathrm{}_{\mathrm{}}`$. By Theorem 3, either $`E`$ contains a complemented copy of $`\mathrm{}_1`$ or $`F`$ contains a copy of $`\mathrm{}_{\mathrm{}}`$. Lemma 5 implies that $`E`$ contains no complemented copy of $`\mathrm{}_1`$, so $`F`$ contains a copy of $`\mathrm{}_{\mathrm{}}`$. Take a normalized basic sequence $`(x_i)E`$ and a bounded sequence of coefficient functionals $`(\varphi _n)E^{}`$ ($`\varphi _i(x_j)=\delta _{ij}`$). Define $`P𝒫(^kE,\mathrm{}_{\mathrm{}})`$ by $`P(x):=\left(\varphi _n(x)^k\right)_n`$. Then, for $`ij`$, we have $`P(x_i)P(x_j)|\varphi _i(x_i)^k\varphi _i(x_j)^k|=1`$. Hence, denoting by $`J:\mathrm{}_{\mathrm{}}F`$ an isomorphism, we get $`JP𝒫(^kE,F)\backslash 𝒫_{wb}(^kE,F)`$, a contradiction.
(c) $``$ (b): The proof is the same of the linear case \[23, Theorem 6\] with slight modifications.
(b) $``$ (a): Take $`P𝒫(^kE,F)`$. Consider the w.u.C. series $`P_i`$ given by Lemma 6 with $`P_i𝒫_{wb}(^kE,F)`$. Since this space contains no copy of $`c_0`$, the series is unconditionally convergent \[7, Theorem V.8\]. Clearly, its sum must be $`P`$. Since the space $`𝒫_{wb}(^kE,F)`$ is closed, we conclude that $`P𝒫_{wb}(^kE,F)`$.
(a) $``$ (d) is trivial.
(d) $``$ (a): If $`𝒫_{wb}(^kE,F)`$ is complemented in $`𝒫(^kE,F)`$, Lemma 5 implies that $`E`$ contains no complemented copy of $`\mathrm{}_1`$. Suppose there is $`P𝒫(^kE,F)\backslash 𝒫_{wb}(^kE,F)`$. By Lemma 6, we can find a sequence $`(P_i)𝒫_{wb}(^kE,F)`$ so that $`P(x)=_{i=1}^{\mathrm{}}P_i(x)`$ unconditionally for each $`xE`$, and $`P_i`$ is w.u.C. but is not unconditionally convergent since $`P𝒫_{wb}(^kE,F)`$. Hence, we can find $`ϵ>0`$, and an increasing sequence $`(m_j)`$ of integers such that, for each $`j`$, the polynomial
$$C_j:=\underset{i=m_j+1}{\overset{m_{j+1}}{}}P_i$$
satisfies $`C_j>ϵ`$.
Define now $`\mathrm{\Phi }:\mathrm{}_{\mathrm{}}𝒫(^kE,F)`$ by $`\mathrm{\Phi }(\xi )(x)=\xi _jC_j(x)`$ for $`\xi =(\xi _j)_{j=1}^{\mathrm{}}\mathrm{}_{\mathrm{}}`$ and $`xE`$. Since the series $`P_i(x)`$ is unconditionally convergent, so is the series $`\xi _jC_j(x)`$. The set $`\{\mathrm{\Phi }(\xi )(x):\xi \mathrm{}_{\mathrm{}},xE\}`$ is contained in the closed linear span of $`\{C_j(E):j\}`$ which is separable by the compactness of $`C_j`$. On the other hand, $`\mathrm{\Phi }(e_j)=C_j𝒫_{wb}(^kE,F)`$ for all $`j`$. By Lemma 4, there is an infinite subset $`M`$ such that $`\mathrm{\Phi }(\xi )𝒫_{wb}(^kE,F)`$ for each $`\xi \mathrm{}_{\mathrm{}}(M)`$.
Therefore, for each $`\xi \mathrm{}_{\mathrm{}}(M)`$, the series $`\xi _jC_j`$ is weak subseries convergent. The Orlicz-Pettis theorem then implies that it is unconditionally convergent. In particular,
$$\underset{\begin{array}{c}jM\\ j\mathrm{}\end{array}}{lim}C_j=0,$$
a contradiction. $`\mathrm{}`$
Observe that the unconditional finite dimensional expansion of the identity is used only in (b) $``$ (a) and (d) $``$ (a).
In the linear case $`(k=1)`$, the restriction $`dim(F)=\mathrm{}`$ is required for the validity of Theorem 7 \[23, Theorem 6\].
###### Remark 8.
In order to highlight the difference between Theorem 3 and Theorem 7, let us consider the spaces $`𝒫(^k\mathrm{}_p,\mathrm{}_q)`$ and $`𝒫_{wb}(^k\mathrm{}_p,\mathrm{}_q)`$ for $`1<p,q<\mathrm{}`$ and $`k>1`$.
If $`kq>p`$, then the space $`𝒫(^k\mathrm{}_p,\mathrm{}_q)`$ contains a copy of $`\mathrm{}_{\mathrm{}}`$. Indeed, define $`J:\mathrm{}_{\mathrm{}}𝒫(^k\mathrm{}_p,\mathrm{}_q)`$ by
$$J(\xi )(x)=\left(\xi _ix_i^k\right)_{i=1}^{\mathrm{}}.$$
Since, for all $`n`$,
$$|\xi _n|=J(\xi )(e_n)J(\xi )=\underset{xB_\mathrm{}_p}{sup}J(\xi )(x)\xi ,$$
we get that $`J`$ is an isometric embedding.
However, the space $`𝒫_{wb}(^k\mathrm{}_p,\mathrm{}_q)`$ contains no copy of $`\mathrm{}_{\mathrm{}}`$ since it is separable. In fact, it is the norm closure of the space of all finite type polynomials generated by the mappings of the form $`\varphi ^ny`$, for $`\varphi (\mathrm{}_p)^{}`$, $`y\mathrm{}_q`$ \[4, Proposition 2.7\].
On the other hand, if $`\xi c_0`$, then $`J(\xi )`$ is in the closure of the finite type polynomials, so $`J(\xi )𝒫_{wb}(^k\mathrm{}_p,\mathrm{}_q)`$. Hence, this space contains $`J(c_0)`$, an isometric copy of $`c_0`$.
Recall moreover that, if $`kq<p`$, then the space $`𝒫(^k\mathrm{}_p,\mathrm{}_q)`$ is reflexive \[1, 4.3\].
Every polynomial $`P𝒫(^kE,F)`$ has an extension, called the Aron-Berner extension, to a polynomial $`\overline{P}𝒫(^kE^{},F^{})`$ (see ). If $`P𝒫_{wb}(^kE)`$, then $`\overline{P}𝒫(^kE^{})`$ is weak-star continuous on bounded sets .
Recall that a Banach space has property (V), introduced in , if every unconditionally converging operator on $`E`$ is weakly compact. Every $`C(K)`$ space has property (V) .
We shall need the following result:
###### Theorem 9.
The following assertions are equivalent:
(a) The space $`E`$ has property (V);
(b) for all $`F`$ and $`k`$, the Aron-Berner extension of every $`P𝒫_{uc}(^kE,F)`$ is $`F`$-valued;
(c) There is $`k`$ such that, for all $`F`$, the Aron-Berner extension of every $`P𝒫_{uc}(^kE,F)`$ is $`F`$-valued.
Easily, if a polynomial is weakly compact, then its Aron-Berner extension is $`F`$-valued .
###### Theorem 10.
For $`k`$, consider the assertions:
(a) $`E`$ has property (V) and $`𝒫(^kE)=𝒫_{wb}(^kE)`$;
(b) for each $`F`$, every $`P𝒫_{uc}(^kE,F)`$ is weakly compact.
Then (a) $``$ (b). If, moreover, $`E`$ has an unconditional finite dimensional expansion of the identity, then (b) $``$ (a).
Proof. (a) $``$ (b): Given $`P𝒫_{uc}(^kE,F)`$, by Theorem 9, the range of its Aron-Berner extension $`\overline{P}`$ is contained in $`F`$. Take a net $`(x_\alpha )B_E`$. We can assume that $`(x_\alpha )`$ is weak Cauchy and so it converges in the weak-star topology to some $`zE^{}`$.
Let $`\psi F^{}`$. Then $`\psi P𝒫_{wb}(^kE)`$ and so,
$$\psi P(x_\alpha )=\psi \overline{P}(x_\alpha )=\overline{\psi P}(x_\alpha )\overline{\psi P}(z)=\psi \overline{P}(z).$$
Therefore, the net $`(P(x_\alpha ))`$ is weakly convergent to $`\overline{P}(z)F`$. So, $`P(B_E)`$ is relatively weakly compact.
(b) $``$ (a): By the comment preceding this Theorem, and by Theorem 9, (b) implies that $`E`$ has property (V). Also, every polynomial in $`𝒫(^kE,\mathrm{}_1)`$ is compact. From this, we obtain that $`𝒫(^kE)`$ contains no copy of $`c_0`$ \[10, Corollary 8\]. A fortiori, $`𝒫_{wb}(^kE)`$ contains no copy of $`c_0`$. Since $`E`$ has an unconditional finite dimensional expansion of the identity, we conclude from Theorem 7 that $`𝒫_{wb}(^kE)=𝒫(^kE)`$. $`\mathrm{}`$
We do not know if the condition on the existence of an unconditional finite dimensional expansion of the identity may be removed from Theorem 10. In fact, if (b) is satisfied, since $`𝒫(^{k1}E)`$ contains no copy of $`c_0`$ and $`E`$ has property (V), we have
$`𝒫(^kE)=_s(E,𝒫(^{k1}E))=𝒲_s(E,𝒫(^{k1}E)),`$
$`𝒫_{wb}(^kE)=𝒦_s(E,𝒫(^{k1}E)).`$
So, we only have to show that $`𝒦_s(E,𝒫(^{k1}E))=𝒲_s(E,𝒫(^{k1}E))`$. There are many conditions on $`E`$ (apart from the existence of an unconditional finite dimensional expansion of the identity) that imply this equality , and it is not known if there are Banach spaces $`X`$, $`Y`$ such that $`𝒦(X,Y)𝒲(X,Y)`$ while $`𝒦(X,Y)`$ contains no copy of $`c_0`$ .
Recall that the condition $`𝒫(^kE)=𝒫_{wb}(^kE)`$ implies that $`E`$ contains no copy of $`\mathrm{}_1`$ .
It is proved in that if $`𝒫(^kE)=𝒫_{wb}(^kE)`$ and $`E`$ has property (u) then, for all $`F`$, every $`P𝒫_{uc}(^kE,F)`$ is weakly compact. Part (a) $``$ (b) of Theorem 10 is stronger than the result of . The latter cannot be applied, for instance, to the space $`c_0\widehat{}_\pi c_0`$ which fails property (u) while it does have property (V) . For the definition of property (u) and its relationship to property (V), see . Recall in particular that, if $`E`$ has property (u) and contains no copy of $`\mathrm{}_1`$, then $`E`$ has property (V) \[27, Proposition 2\].
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# Elastic Scattering as a Cause of Quantum Dephasing: The Conductance of Two-Dimensional Imperfect Conductors
## I Introduction
Since the original formulation of the localization problem by Anderson , the question of whether electronic states in disordered systems are localized at any strength of disorder or a mobility edge can be defined by a critical level of disorder has become a central issue. In the former case this yields an insulating-type behaviour of large samples, while in the latter the metallic-type transport of conducting electrons is allowed due to the existence of extended states.
It was subsequently ascertained that the answer to this question depends substantially on the dimensionality of the disordered system. For the case of one-dimensional (1D) conductors in the limit of vanishing temperature one can prove in a mathematically rigorous way (see Ref. and references therein) that the spectrum of the electrons, subject to an arbitrarily small but finite random potential in infinitely long samples, is purely discrete, i.e. all electron states are necessarily localized irrespective of their energy. As a result, the DC conductivity vanishes for such systems, whereas in finite 1D samples the conductance falls exponentially as the length of the conductor grows.
In contrast to the 1D case, no mathematically rigorous theory of localization exists for two- (2D) and tree-dimensional (3D) random systems. Formulation of the problem in terms of the renormalization group (RG) led to the one-parameter scaling hypothesis of localization and appeared to provide a considerable progress in studying systems of dimensionality greater than unity. The main conviction gained from this development was that all the electron states in both 1D and 2D disordered systems are localized at arbitrary small disorder. Hence, the metal-insulator transition (MIT) is usually believed to occur only in 3D systems, while (sufficiently large) 1D and 2D systems are always Anderson insulators.
This opinion, established quite long ago, has recently been challenged after unexpected experimental detection of MIT in dilute 2D electron and hole systems. Originally obtained on Si-MOSFETs, the results have entailed numerous experimental findings of MIT in other dilute 2D systems, leading to the well-founded belief that the effect is of a rather general nature. To elucidate the unconventional experimental facts, different theoretical approaches were put forward including quite disputable ones (see the discussion in Ref. ). Among the variety of options the greatest anticipations in explaining the experiments are mostly associated with accounting for the Coulomb interaction of carriers. Yet the transport theories for correlated particles in the presence of disorder still cannot claim for general acceptance because of the substantial controversy in the interpretation of the role of interaction within different parameter domains which correspond to diffusive and localized regimes.
Traditional understanding of electron localization in 2D disordered systems and numerous experimental facts which indicate the metallic nature of charge transport at low temperatures have brought about active progress in the line of study that is identified with quantum dephasing. The main efforts in this direction are focused on the detection of various phase-breaking mechanisms responsible for the delocalization of quasi-particles primarily localized by disorder. However, in spite of a large body of versions on offer the problem still remains open.
Our intention in this contribution is to re-examine the conventional one-particle approach without appealing to scattering mechanisms other than those characteristic of quenched disorder. Evidence will be given that even within this elementary model the experimental findings of Refs. , at first sight curious, are as a matter of fact quite natural. Numerous attempts to interpret the results of Ref. within the framework of a single-particle approach were made, in particular, by improving the scaling approach. In this study, however, a fundamentally different strategy is chosen which is an alternative to the RG analysis. We prefer to obtain the observables directly, while conclusions (though indirect) about the localization of electron states are made on the basis of the results.
It is instructive to recall that working out, even without a profound spectral analysis, practical asymptotic methods for calculating the disorder-averaged many-particle characteristics (conductivity, density-density correlator, etc.) turned out to be more helpful for the establishment of a highly advanced theory of 1D random systems than the development of rigorous mathematical foundations. The usefulness of such an approach can be attributed to the fact that in the context of the above-mentioned essentially perturbative methods one has managed to trace with the required accuracy the effect of mutual interference of quantum waves corresponding to multiply backscattered current carriers. In such a way, physical results entirely consistent with the anticipations based on mathematical predictions were eventually obtained. The present research was primarily induced by long-standing discontent associated with the lack of arguments of a comparable standard, either in favour of localization or against it, as applied to 2D systems of degenerate electrons subject to a static random potential.
Commonly, the presence of inelastic scattering mechanisms is held responsible as a main cause of preventing quantum interference and, thus, Anderson localization. Among these are the electron-phonon and electron-electron interactions and other conceivable methods of energy interchange between the electron bath and the environment. These can lead to the loss of phase (meaning energy) memory or, in other words, phase coherence of electronic states. Note in this connection that in the 1D case the demand of energy coherence admits a large transfer of momentum for onefold scattering of degenerate electrons in the backward direction. This leads inevitably to considerable local breaking of spatial (instead of temporal) phases of the wavefunctions. It was recognized that such a large violation of spatial phases is quite helpful when deriving a constructive theory of 1D quantum transport. On the basis of such arguments, the selection of efficient subsets of terms in perturbative expansion which are responsible for the effective interference of electronic waves scattered iteratively in a backward direction was suggested. This interference finally results in the formation of true localized states, even when one starts their perturbative construction from plain-wave-like trial functions which belong to a mathematical class different from that of localized functions.
An analogous scenario of the perturbative formation of localized states in 2D systems cannot give the same result, since for sufficiently isotropic scattering the spatial coherence of the wavefunctions is already violated at distances of the order of the quasi-classical mean free path even for weak energy scattering. The coherence is maintained efficiently only within a small phase volume, thus resulting in the relative smallness of the interference corrections, usually known as weak localization corrections for non-one-dimensional systems.
Nevertheless, the analysis of the problem of multidimensional localization with the use of a one-dimensional scenario turns out to be quite productive. As we shall prove below the problem of electron transport in 2D open system of waveguide type can be reduced without any approximations to the set of one-dimensional (though non-Hermitian) problems for the quantum waves propagating in individual conducting channels. The channels will be identified with extended waveguide modes. For the corresponding dynamic equations that are one-dimensional, an opportunity arises for making substantial use of spatial phases of the wavefunctions instead of their temporal parts. This turns out to be preferable from the technical point of view for solving stationary problems of the electron, as well as classical wave, transport. The reason for the usefulness of such a procedure lies in the fact that in one-dimensional problems spatial averaging has been shown to be highly advantageous, leading to the reduction of the perturbative expansions of the physical observables to a summable series.
Being exactly quantum in nature, the waveguide approach used here is, to a certain degree, less obvious than semi-classical ones usually applied in most localization theories. Thus, in its context the clarity of the path-based interpretation is substantially lost. At the same time, the benefit of our method is that one-dimensional ‘channel’ equations enable one to distinguish unambiguously the coherent intra-mode scattering, which is easily interpreted from the standpoint of 1D localization theory, and the inter-mode scattering which corresponds, although not quite directly, to isotropic scattering of semi-classical electrons. The quantum states in different channels are specified by different longitudinal momenta. This difference essentially suppresses interference of primary and scattered electronic waves, if they belong to different channels. As a result, the inter-mode scattering turns out to be an intrinsic origin of the inability for 2D electrons to be localized by weak static disorder.
The ‘one-dimensionalization’ procedure suggested in this paper gives an opportunity to highlight the role of spatial coherence in the interference of electronic waves even in a single-particle approximation. The results obtained here by conventional perturbative methods enable us to conclude about the unrealizability of the strong (Anderson) localization in systems whose spatial dimensionality is greater than unity, irrespective of their size. It should be particularly emphasized that decoherence appearing due to inter-mode scattering is unconnected to genuine inelasticity in the interaction of electrons with disorder. However, the difference in longitudinal energy between the conducting channels could be treated as a source of ‘hidden inelasticity’, if one is accustomed to such an interpretation.
The paper is organized as follows. In the next section, the problem is formulated using linear response theory. In section III we develop a method of exact one-dimensionalization which is a central point of the paper. Then, in section IV, the trial Green functions supremely important for the developed technique are analyzed with the aid of a two-scale perturbation method. A spectral analysis of the electron system is given in section V. In the final two sections we present asymptotic expressions for the conductance and discuss the results. A pair of tedious but important calculations is presented in two appendices.
## II Statement of the problem
A common approach used in studies of random systems of various dimensionality is to take a hyper-cube of a certain linear size and vary the size while searching for the conductance scaling. Such an approach seems to be natural when studying spectral properties of closed systems. At the same time, it is not quite appropriate for solving transport problems as applied to structures of waveguide type, in particular, quantum conductors of arbitrary length.
In this paper, we examine the case of a 2D imperfect rectangular sample of length $`L`$ in the $`x`$-direction and width $`D`$ in the lateral direction $`y`$. Degenerate non-interacting spinless electrons are assumed to be confined between the hard-wall side boundaries $`y=\pm D/2`$, whereas in the direction of current ($`x`$) we suppose the system to be open at the strip ends $`x=\pm L/2`$. The dimensionless conductance $`g(L)`$ (in units $`e^2/\pi \mathrm{}`$) is calculated directly from linear response theory, whence at zero temperature we have
$$g(L)=\frac{4}{L^2}𝑑𝒓𝑑𝒓^{}\frac{G(𝒓,𝒓^{})}{x}\frac{G^{}(𝒓,𝒓^{})}{x^{}}.$$
(1)
Here the integration with respect to $`𝒓=(x,y)`$ is performed over the area occupied by the conductor
$$x(L/2,L/2),y(D/2,D/2).$$
(2)
$`G(𝒓,𝒓^{})`$ is the retarded Green function of the conducting electrons. Within the isotropic Fermi liquid model this function is governed by the equation
$$\left[\mathrm{\Delta }+k_F^2+i0V(𝒓)\right]G(𝒓,𝒓^{})=\delta (𝒓𝒓^{}).$$
(3)
We adopt the system of units with $`\mathrm{}=2m=1`$ ($`m`$ is the electron effective mass), so that $`\mathrm{\Delta }`$ is a two-dimensional Laplace operator, $`k_F`$ is the Fermi wavenumber of the electrons, $`V(𝒓)`$ is the ‘bulk’ static random potential.
The potential $`V(𝒓)`$ in equation (3) will be regarded as a short-range one and not necessarily isotropic. The term ‘short-range’ implies the characteristic spatial interval over which the potential is substantially varied to be small compared with the ‘macroscopic’ lengths of the problem, namely the electron mean free path and the conductor length. Being considered as a statistical variable, the potential $`V(𝒓)`$ will be specified by a zero mean value and the binary correlation function
$$V(𝒓)V(𝒓^{})=𝒬W(𝒓𝒓^{}).$$
(4)
Here the angular brackets denote ensemble averaging and $`W(𝒓)`$ is some function normalized to unity. The explicit form of this function is not so important. In many cases $`W(𝒓)`$ is approximated by the delta-function, $`W(𝒓)=\delta (𝒓)`$. However, the method developed in this contribution permits one to consider not only isotropic and not necessarily local scattering events. In addition, the choice of $`W(𝒓)`$ in the form of a delta-function, apart from restricting the physical applicability of the model, is not quite convenient from the technical viewpoint. Thus, for example, when calculating the corrections to the mode spectrum, equations (V), a formal problem can arise of the divergence of the evanescent mode contribution. It is certainly absent provided the potential $`W(𝒓)`$ is not exactly point-supported — this problem is familiar in quantum mechanics. To get rid of the divergences in the problems of dimension more than one it is merely sufficient to choose the function $`W(𝒓)`$ to be less singular than $`\delta (𝒓)`$. Therefore, in order to focus on the main questions, in place of the correlation equality (4), we shall use the expression below
$$V(𝒓)V(𝒓^{})=𝒬W(xx^{})\delta (yy^{}).$$
(5)
Additionally, this form of equation (4) allows us to consider anisotropic scattering. Similar to $`W(𝒓)`$ from (4), the function $`W(x)`$ in (5) is normalized to unity.
## III The one-dimensionalization procedure
### A The general scheme
It seems intuitively natural for an open system with the prescribed direction of quasi-particle transport to be considered as to some extent a one-dimensional object. However, the realizability of ‘one-dimensionalization’ at the level of dynamic equations, which would be quite important from a mathematical perspective, is not a priori apparent. Here the term ‘one-dimensionalization’ means reduction of the two-dimensional differential equation (3) to a set of one-dimensional equations. Although a system of waveguide type can often be regarded as a collection of one-dimensional quantum channels, the latter are not independent in general. Normally, they are strongly coupled with each other through static or dynamic inhomogeneities present in the problem.
Nevertheless, in what follows we intend to show that just the waveguide nature of a system under consideration enables the mathematical description of the transport problem for a 2D region (2) to be reduced to a set of independent strictly one-dimensional, although non-Hermitian, problems posed on the interval $`x(L/2,L/2)`$, regardless of the strength of the disorder. To perform the reduction one should merely pass to the mode representation, i.e. Fourier transform in the transverse coordinate $`y`$, using some complete set {$`\varphi (y)`$} of eigenfunctions of the transverse free-electron Hamiltonian, namely the Laplace operator in Eq. (3). The conductance (1) acquires the form
$$g(L)=\frac{4}{L^2}_L𝑑x𝑑x^{}\underset{n,n^{}=1}{\overset{N_c}{}}\frac{G_{nn^{}}(x,x^{})}{x}\frac{G_{nn^{}}^{}(x,x^{})}{x^{}},$$
(6)
where $`N_c=[k_FD/\pi ]`$ is the number of conducting channels or, in other words, extended waveguide modes. Equation (3) is then transformed into a set of coupled ordinary differential equations for the Fourier components $`G_{nn^{}}(x,x^{})`$,
$$\left[\frac{^2}{x^2}+k_n^2+i0V_n(x)\right]G_{nn^{}}(x,x^{})\underset{\genfrac{}{}{0pt}{}{m=1}{(mn)}}{\overset{\mathrm{}}{}}U_{nm}(x)G_{mn^{}}(x,x^{})=\delta _{nn^{}}\delta (xx^{}).$$
(7)
Here $`k_n^2=k_F^2(n\pi /D)^2`$ is the longitudinal energy of the $`n`$th mode, $`U_{nm}(x)`$ is the inter-mode matrix element of the potential $`V(𝒓)`$,
$$U_{nm}(x)=_D𝑑y\phi _n(y)V(𝒓)\phi _m(y).$$
(8)
Note the difference between the summation limits in Eqs. (6) and (7). Restriction of the summation in (6) by the number of conducting channels implies, strictly speaking, weakness of the electron scattering, to be specified in Sec. III B. Under the same assumption formula (1) itself is valid, where the products of the Green functions of the same kind (both retarded and advanced) have already been omitted. In equation (7), the summation is naturally performed over a complete set of waveguide modes.
Particular attention should be drawn to the fact that the term containing the diagonal (intra-mode) matrix element $`U_{nn}(x)V_n(x)`$ is initially detached from the sum of Eq. (7), so that the matrix $`U_{nm}`$ hereafter is held off-diagonal in the discrete mode variable. This little technical trick enables one to reduce the problem of finding the overall of the functions $`G_{nn^{}}(x,x^{})`$ to the solution of a subset of purely one-dimensional closed equations for the mode-diagonal functions $`G_{nn}(x,x^{})`$ only. To this end we first introduce the auxiliary trial Green functions $`G_n^{(V)}(x,x^{})`$ ($`n\mathrm{}`$), each obeying the equation
$$\left[\frac{^2}{x^2}+k_n^2+i0V_n(x)\right]G_n^{(V)}(x,x^{})=\delta (xx^{}),$$
(9)
and Sommerfeld’s radiation conditions at the strip ends $`x=\pm L/2`$. These conditions seem to be natural to impose on an open system. For the case of the 1D equation (9) they acquire the form
$$\left(\frac{}{x}ik_n\right)G_n^{(V)}(x,x^{})|_{x=\pm L/2}=0.$$
(10)
It implies that the field of the $`n`$th mode radiated by the point source placed at $`x^{}(L/2,L/2)`$ reaches the corresponding ($`\pm `$) end of the interval and then propagates unscattered with the conserved momentum $`k_n`$ beyond the ends of the conductor. Possible scattering in the leads attached to the strip from the left and right should be taken into account separately.
Although a solution of the stochastic problem (9) and (10) cannot be obtained in quadratures, in the case of weak scattering specified below by inequalities (21) the main features of the solution can be extracted by due consideration with any desired accuracy in the intra-mode potential $`V_n(x)`$. The necessary analysis is presented in the next section. As for now, we merely consider all the functions $`G_n^{(V)}(x,x^{})`$ as a priori known ones whose properties are specified solely by the elastic intra-mode scattering. With these functions chosen as an initial approximation for the exact mode functions $`G_{nn}`$, only the inter-mode scattering associated with the off-diagonal matrix $`U_{nm}(x)`$ will then be taken as a perturbation. To implement this intent, we turn from equation (7) to the consequent integral equation,
$$G_{nn^{}}(x,x^{})=G_n^{(V)}(x,x^{})\delta _{nn^{}}+\underset{m=1}{\overset{\mathrm{}}{}}_L𝑑x_1𝖱_{nm}(x,x_1)G_{mn^{}}(x_1,x^{}).$$
(11)
Here the kernel
$$𝖱_{nm}(x,x_1)=G_n^{(V)}(x,x_1)U_{nm}(x_1)$$
(12)
already contains only those harmonics of the potential $`V(𝒓)`$ which are responsible for the inter-mode scattering. A thorough study of system (11) leads to a conjecture that all non-diagonal mode elements $`G_{mn}`$ ($`mn`$) can be expressed only via the diagonal element $`G_{nn}`$ by means of some linear operator $`\widehat{𝖪}`$,
$$G_{mn}(x,x^{})=_L𝑑x_1𝖪_{mn}(x,x_1)G_{nn}(x_1,x^{}).$$
(13)
To specify this operator, one should separate the term with the diagonal (in the mode variable) Green function on the right-hand side of equation (11) and substitute all non-diagonal Green functions in the form of equation (13). Then, after renaming the mode variables, we arrive at the following equation for the matrix elements of the operator $`\widehat{𝖪}`$:
$$𝖪_{mn}(x,x^{})=𝖱_{mn}(x,x^{})+\underset{\genfrac{}{}{0pt}{}{k=1}{(kn)}}{\overset{\mathrm{}}{}}_L𝑑x_1𝖱_{mk}(x,x_1)𝖪_{kn}(x_1,x^{}).$$
(14)
Equation (14) belongs to a class of multi-channel Lippmann-Schwinger equations that are known to be extremely singular, in contrast to their single-channel counterparts. Nevertheless, by choosing the trial Green function $`G_n^{(V)}`$ as a zero approximation and perturbing it only by the inter-mode potentials, one manages to avoid the above mentioned singularity. Note that mode indices $`m`$ and $`k`$ in Eq. (14) take all the positive integer values except for the value $`n`$. This urges one to interpret the functions appearing in (14) as matrix elements of operators acting in two-dimensional mixed coordinate-mode space ($`x,m`$) which does not include the mode $`n`$ (the notation $`\overline{𝖬}_𝗇`$ will be used for that space). The presence in Eq. (14) of the right-hand index $`n`$, which does not belong to $`\overline{𝖬}_𝗇`$, can be ensured by introducing the projection operator $`𝐏_n`$ that will make the mode index of any operator standing next to it (both from the left or right) equal to $`n`$. With this convention accepted, it follows from equation (14) that the operator $`\widehat{𝖪}`$ implementing relation (13) has the form
$$\widehat{𝖪}=\left(𝟙\widehat{𝖱}\right)^1\widehat{𝖱}𝐏_n.$$
(15)
Here $`\widehat{𝖱}`$ is a 2D operator acting on $`\overline{𝖬}_𝗇`$ and is specified by the matrix elements (12).
As for the mathematical correctness of the operator representation (15), it depends on the existence of the inverse operator $`\left(𝟙\widehat{𝖱}\right)^1`$. This point is discussed in Appendix A, where we provide a proof that detachment of the intra-mode potential $`V_n(x)`$ in Eq. (7) prevents the possible singularity.
Thus, equations (13) and (15) reduce the problem of finding the complete Green function $`G(𝒓,𝒓^{})`$ within the 2D region (2) to calculation of its diagonal mode components only. In order to do that it is necessary to put $`n^{}=n`$ in Eq. (7) and substitute all non-diagonal components $`G_{mn}`$ from equation (13). As a result, the closed equation for the diagonal component $`G_{nn}`$ is deduced,
$$\left[\frac{^2}{x^2}+k_n^2+i0V_n(x)\widehat{𝒯}_n\right]G_{nn}(x,x^{})=\delta (xx^{}).$$
(16)
In equation (16), in addition to the local intra-mode potential $`V_n(x)`$, the operator potential $`\widehat{𝒯}_n`$ has arisen,
$$\widehat{𝒯}_n=𝐏_n\widehat{𝒰}\left(𝟙\widehat{𝖱}\right)^1\widehat{𝖱}𝐏_n=𝐏_n\widehat{𝒰}\left(𝟙\widehat{𝖱}\right)^1𝐏_n.$$
(17)
Here $`\widehat{𝒰}`$ is the inter-mode operator specified on $`\overline{𝖬}_𝗇`$ by the matrix elements
$$\mathbf{<}x,l\mathbf{\left|}\widehat{𝒰}\mathbf{\right|}x^{},m\mathbf{>}=U_{lm}(x)\delta (xx^{}).$$
(18)
The potential $`\widehat{𝒯}_n`$ has quite an interesting interpretation. In the right-hand side of equation (17) there is a conventional $`T`$-matrix enveloped by the projective operators, one of which removes the excitation from mode $`n`$ and the other restores it back to the same mode after the appropriate scattering events within the subspace $`\overline{𝖬}_𝗇`$. Hence, the potential $`\widehat{𝒯}_n`$, although corresponding to effectively intra-mode scattering, actually includes all inter-mode scattering events undergone by the excitation while passing over ‘closed paths’ in the mode space. The intra- and inter-mode scattering mechanisms turn out to be attributed to different potentials in equation (16), facilitating significantly the subsequent interpretation of the results. In what follows the potentials $`V_n(x)`$ and $`\widehat{𝒯}_n`$ will be referred to as those responsible for direct intra-mode and inter-mode scattering, respectively.
Finally, in this subsection we express the conductance (6) through the diagonal elements of the mode Green matrix. After rearranging the terms in Eq. (6) and using relation (13) we obtain
$`g(L)={\displaystyle \frac{4}{L^2}}{\displaystyle \underset{n=1}{\overset{N_c}{}}}`$ $`{\displaystyle }{\displaystyle _L}dxdx^{}[{\displaystyle \frac{G_{nn}(x,x^{})}{x}}{\displaystyle \frac{G_{nn}^{}(x,x^{})}{x^{}}}`$ (19)
$`+{\displaystyle \underset{\genfrac{}{}{0pt}{}{m=1}{(mn)}}{\overset{N_c}{}}}`$ $`{\displaystyle }{\displaystyle _L}dx_1dx_2{\displaystyle \frac{𝖪_{mn}(x,x_1)}{x}}G_{nn}(x_1,x^{})𝖪_{mn}^{}(x,x_2){\displaystyle \frac{G_{nn}^{}(x_2,x^{})}{x^{}}}].`$ (20)
Expression (20) jointly with equation (16) completes, in principle, the ‘one-dimensionalization’ procedure introduced at the beginning of this section. In this form the problem under consideration is convenient for a numerical treatment at any disorder strength, because the solution of the 2D problem governed by equation (3) is reduced to a finite set of purely 1D problems (9) and (16). At the same time, assuming weak electron scattering (in the semi-classical sense), we manage to proceed with our analytical consideration and obtain the results.
### B The weak scattering approximation
It is natural to specify the intensity of electron scattering in terms of characteristic spatial scales inherent to the problem. Henceforth we recognize the scattering as weak provided the following inequalities hold:
$$k_F^1,r_c\mathrm{}.$$
(21)
Here $`r_c`$ is the correlation radius of the potential $`V(𝒓)`$, $`\mathrm{}=2k_F/𝒬`$ is the semiclassical mean free path of electrons evaluated within the model of a $`\delta `$-correlated 2D random potential, i.e. $`W(𝒓)=\delta (𝒓)`$ in Eq. (4).
Estimation of the norm of the operator $`\widehat{𝖱}`$ specified on $`\overline{𝖬}_𝗇`$ by the matrix elements (12) results in
$$\widehat{𝖱}^2\frac{D/L}{k_F\mathrm{}}.$$
(22)
Under conditions (21), this enables us to find an expansion to lowest order in the impurity potential of the operator $`\widehat{𝖪}`$, Eq. (15), so that it becomes approximately equal to $`\widehat{𝖱}`$ almost regardless of the conductor aspect ratio. The exact operator $`\widehat{𝒯}_n`$ from (17) can, in turn, be replaced by the approximate value
$$\widehat{𝒯}_n𝐏_n\widehat{𝒰}\widehat{G}^{(V)}\widehat{𝒰}𝐏_n$$
(23)
where the operator $`\widehat{G}^{(V)}`$ is specified on $`\overline{𝖬}_𝗇`$ by the matrix elements
$$x,k\left|\widehat{G}^{(V)}\right|x^{},m=\delta _{km}G_m^{(V)}(x,x^{}).$$
Besides the reduction of the $`T`$-operator (17) to truncated form (23), a substitution of the approximate matrix elements of the operator $`\widehat{𝖪}`$ brings the second term of conductance (20) to the form
$$_L𝑑x_1𝑑x_2U_{mn}(x_1)U_{mn}(x_2)\frac{G_m^{(V)}(x,x_1)}{x}G_{m}^{(V)}{}_{}{}^{}(x,x_2)G_{nn}(x_1,x^{})\frac{G_{nn}^{}(x_2,x^{})}{x^{}},$$
(24)
which is convenient for performing the ensemble averaging. It will be shown below that all Green functions in (24), and not only the trial ones, may be thought of as independent of the inter-mode potentials $`U_{mn}`$. Yet correlation of those potentials with the intra-mode one, $`V_n(x)`$, governing the trial Green functions, can be disregarded in view of Eq. (44). As a result, after averaging conductance (20) with the use of (24) and (5), the expression, which will be subject to a further analysis, takes the form
$`<g(L)>=`$ $`{\displaystyle \frac{4}{L^2}}{\displaystyle \underset{n=1}{\overset{N_c}{}}}{\displaystyle }{\displaystyle _L}dxdx^{}[<{\displaystyle \frac{G_{nn}(x,x^{})}{x}}{\displaystyle \frac{G_{nn}^{}(x,x^{})}{x^{}}}>`$ (25)
$`+`$ $`{\displaystyle \frac{𝒬}{D}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{m=1}{(mn)}}{\overset{N_c}{}}}{\displaystyle }{\displaystyle _L}dx_1dx_2W(x_1x_2)<G_{m}^{(V)}{}_{}{}^{}(x,x_2){\displaystyle \frac{}{x}}G_m^{(V)}(x,x_1)><G_{nn}(x_1,x^{}){\displaystyle \frac{}{x^{}}}G_{nn}^{}{}_{}{}^{}(x_2,x^{})>].`$ (26)
## IV Analysis of the trial Green functions
The trial Green functions $`G_m^{(V)}(x,x^{})`$ enter the potential $`\widehat{𝒯}_n`$, Eq. (23), thus determining the exact mode functions $`G_{nn}(x,x^{})`$, and the second term of the conductance (26). Although these functions appear as subsidiary mathematical objects, they are of great concern for the problem and therefore deserve particular consideration. The study of the trial functions is also instructive, since it may provide useful insights into the analysis of some misinterpretation regarding 2D localization which existed until recently in the literature.
Note first that the perturbative solution of equation (9) depends substantially on whether the corresponding unperturbed waveguide mode is either extended or evanescent. The Green functions of evanescent modes with $`n>N_c`$ are localized even without any perturbation. To find them in the limit of weak scattering, it is sufficient to restrict oneself to zero-order perturbation in the potential $`V_n(x)`$,
$$G_n^{(V)}(x,x^{})=\frac{1}{2|k_n|}\mathrm{exp}\mathbf{\left(}|k_n||xx^{}|\mathbf{\right)},n>N_c.$$
(27)
The problem is much more involved for the extended modes, $`n<N_c`$. Inasmuch as the function $`G_n^{(V)}(x,x^{})`$ is defined as a solution of the strictly one-dimensional problem (9) and (10), to find it correctly in the context of localization theory the plain-wave-based zero approximation is not quite appropriate. This stems from the fact that in such an approximation it is rather difficult to account for the interference of multiply backscattered waves. Instead we apply the two-scale perturbation method analogous to that used in the theory of non-linear oscillations. This method showed itself well for solving the problem of charge transport in extremely narrow, namely single-mode, surface-corrugated conducting strips. Below an outline of the method is given together with some essential results. The details of their derivation are deferred to Appendix B.
The method used in Ref. is based on the representation of the 1D Green function, which is the solution of the boundary-value problem (9) and (10), via the solutions of the appropriate Cauchy problems,
$$G_n^{(V)}(x,x^{})=𝓦_n^1\left[\psi _+(x|n)\psi _{}(x^{}|n)\mathrm{\Theta }(xx^{})+\psi _+(x^{}|n)\psi _{}(x|n)\mathrm{\Theta }(x^{}x)\right],$$
(28)
Here the functions $`\psi _\pm (x|n)`$ are linearly independent solutions of the homogeneous equation (9) with the radiation conditions analogous to (10) satisfied at the strip ends $`x=\pm L/2`$, according to the ‘sign’ index. The Wronskian of those functions is $`𝓦_n`$, $`\mathrm{\Theta }(x)`$ is Heaviside’s step function. This reformulation of a boundary-value problem in terms of an initial-value problem will allow one to perform averaging over the disorder later on.
It is advantageous to represent the functions $`\psi _\pm (x|n)`$ as superpositions of modulated harmonic waves propagating in opposite directions along the $`x`$-axis,
$$\psi _\pm (x|n)=\pi _\pm (x|n)\mathrm{exp}(\pm ik_nx)i\gamma _\pm (x|n)\mathrm{exp}(ik_nx).$$
(29)
Within the framework of the weak scattering approximation (21), the amplitudes $`\pi _\pm (x|n)`$ and $`\gamma _\pm (x|n)`$ vary slowly as compared with the ‘fast’ exponentials $`\mathrm{exp}(\pm ik_nx)`$, so that the radiation conditions for $`\psi _\pm (x|n)`$ are reformulated as the ‘initial’ conditions for the smooth amplitudes as follows:
$$\pi _\pm (\pm L/2|n)=1,\gamma _\pm (\pm L/2|n)=0.$$
(30)
In view of the smoothness of $`\pi _\pm `$ and $`\gamma _\pm `$, differential equations for them can be derived by means of averaging the equations for $`\psi _\pm (x|n)`$ over an arbitrary-valued spatial interval intermediate between the ‘microscopic’ lengths $`k_n^1`$ and $`r_c`$ on the one hand, and the ‘macroscopic’ lengths on the other hand. Among the latter lengths are the scattering length, to be specified in the course of the solution, and the sample length $`L`$. For weak scattering, the equations for $`\pi _\pm `$ and $`\gamma _\pm `$ are reduced to the coupled first-order ones,
$$\begin{array}{ccc}& & \pi _\pm ^{}(x|n)\pm i\eta _n(x)\pi _\pm (x|n)\pm \zeta _{n\pm }^{}(x)\gamma _\pm (x|n)=0,\hfill \\ & & \gamma _\pm ^{}(x|n)i\eta _n(x)\gamma _\pm (x|n)\pm \zeta _{n\pm }(x)\pi _\pm (x|n)=0.\hfill \end{array}$$
(31)
Here the variable coefficients $`\eta _n(x)`$ and $`\zeta _{n\pm }(x)`$ are random fields associated with the intra-mode potential $`V_n(x)`$ in the following way:
$$\eta _n(x)=\frac{1}{2k_n}\underset{xl}{\overset{x+l}{}}\frac{dt}{2l}V_n(t),\zeta _{n\pm }(x)=\frac{1}{2k_n}\underset{xl}{\overset{x+l}{}}\frac{dt}{2l}\mathrm{e}^{2ik_nt}V_n(t).$$
(32)
For the intermediate property of the averaging interval $`2l`$, the fields $`\eta _n(x)`$ and $`\zeta _{n\pm }(x)`$ are, in fact, nothing but the ‘narrow’ sets of spatial harmonics of the potential $`V_n(x)`$ with the momenta close to zero and $`\pm 2k_n`$, respectively. The real field $`\eta _n(x)`$ is responsible for the ‘forward’ scattering, while the complex one $`\zeta _{n\pm }(x)`$ for the ‘backward’ scattering of the $`n`$th waveguide mode.
Under the assumption of weak scattering, only binary correlators of the random potentials govern the majority of statistical characteristics of physical quantities. It was shown in Ref. that only the two correlation functions, $`\eta _n(x)\eta _n(x^{})`$ and $`\zeta _{n\pm }(x)\zeta _{n\pm }^{}(x^{})`$, of modified random fields (32) may be thought of as non-vanishing. Calculation with the use of model (5) readily gives
$`\eta _n(x)\eta _n(x^{})`$ $`=`$ $`{\displaystyle \frac{1}{L_f^{(V)}(n)}}F_l(xx^{}),`$ (34)
$`\zeta _{n\pm }(x)\zeta _{n\pm }^{}(x^{})`$ $`=`$ $`{\displaystyle \frac{1}{L_b^{(V)}(n)}}F_l(xx^{}).`$ (35)
Here
$$L_f^{(V)}(n)=\frac{2D}{3𝒬}(2k_n)^2\text{and}L_b^{(V)}(n)=\frac{2D}{3𝒬}\frac{(2k_n)^2}{\stackrel{~}{W}(2k_n)}$$
(36)
are the forward and backward mode scattering lengths, respectively; $`\stackrel{~}{W}(q)`$ is the Fourier transform of $`W(x)`$ from Eq. (5). The function
$$F_l(x)=\underset{\mathrm{}}{\overset{\mathrm{}}{}}\frac{dq}{2\pi }\text{e}^{iqx}\frac{\mathrm{sin}^2(ql)}{(ql)^2}=\frac{1}{2l}\left(1\frac{|x|}{2l}\right)\mathrm{\Theta }(2l|x|)$$
(37)
in Eqs. (IV) is sharp at the scale of ‘macroscopic’ lengths, so it can be regarded as the $`\delta `$-function in the ‘distributional’ sense, $`F_l(x)\delta (x)`$.
Equations (31) and correlation relations (IV) enable one to obtain the entire statistical information concerning the trial Green function for any extended mode $`n`$. In Appendix B derivation of all plain moments of that function is presented with the following result:
$`<\left[G_n^{(V)}(x,x^{})\right]^\mu >=\left({\displaystyle \frac{i}{2k_n}}\right)^\mu \mathrm{exp}\left[i\mu k_n|xx^{}|{\displaystyle \frac{\mu }{2}}\left({\displaystyle \frac{\mu }{L_f^{(V)}(n)}}+{\displaystyle \frac{1}{L_b^{(V)}(n)}}\right)|xx^{}|\right],`$ (38)
$`\mu \mathrm{}.`$ (39)
It is noteworthy that from Eqs. (36) we have the following estimate:
$$L_{f,b}^{(V)}(n)N_c\mathrm{}\mathrm{cos}^2\vartheta _n,$$
(40)
where $`\vartheta _n`$ is a ‘sliding angle’ of the $`n`$th mode with respect to the $`x`$-axis, $`|\mathrm{sin}\vartheta _n|=n\pi /k_FD`$. The value of (40) is coincident, in order of magnitude, with the localization length widely believed to be characteristic for multi-mode quasi-one-dimensional (Q1D) waveguide systems (see, e.g., Refs.). However, in the present theory lengths (36) are nothing but the extinction lengths of the auxiliary excitations that are not subjected to inter-mode scattering.
Besides plain moments (39), important characteristics of the random function $`G_n^{(V)}(x,x^{})`$ are the mixed moments $`<\left[G_n^{(V)}(x,x^{})\right]^\mu \left[G_{n}^{(V)}{}_{}{}^{}(x,x^{})\right]^\nu >`$. At $`\mu =\nu `$ all of them are smooth (not oscillatory) functions of the argument $`|xx^{}|`$ whose spatial decrease is determined by the one-dimensional localization length $`\xi _n=4L_b^{(V)}(n)`$. The second term in (26) contains one of the simplest correlators of this type, the ‘density-current’ correlator $`<G^{}G>`$. It was already studied in applications to the problem of classical wave transport. We do not present here the exact expression for this correlator since only the fact of its exponential decrease at the localization length $`\xi _n`$ is of significance for our analysis,
$$<G_{m}^{(V)}{}_{}{}^{}(x,x^{})\frac{}{x}G_m^{(V)}(x,x^{})>\mathrm{exp}\left(\frac{|xx^{}|}{\xi _n}\right).$$
(41)
## V The mode states spectrum: Perturbative treatment
For calculation of the diagonal Green function $`G_{nn}(x,x^{})`$ in equation (16) by means of the perturbation technique, it would make a good sense to reconstruct the operator potential $`\widehat{𝒯}_n`$ so that the mean value equals zero. To that end one has to define the result of the action of the operator $`\widehat{𝒯}_n`$ on the function $`G_{nn}`$. For $`n>N_c`$ (evanescent modes) the function $`G_{nn}`$ can be left in its unperturbed form (27) due to the weakness of scattering. For $`n<N_c`$, with the operator nature of the potential $`\widehat{𝒯}_n`$ and reduced form (23) of the $`T`$-operator taken into account, the following integral has to be calculated:
$$\widehat{𝒯}_nG_{nn}(x,x^{})=_L𝑑x_1𝖳_n(x,x_1)G_{nn}(x_1,x^{}).$$
(42)
Here the kernel $`𝖳_n(x,x_1)`$ is given by
$$𝖳_n(x,x_1)=\underset{\genfrac{}{}{0pt}{}{m=1}{(mn)}}{\overset{\mathrm{}}{}}U_{nm}(x)G_m^{(V)}(x,x_1)U_{mn}(x_1).$$
(43)
On performing the ensemble averaging in Eq. (43), it is justifiable to neglect the correlation between the inter-mode potentials $`U_{nm}(x)`$ and the intra-mode ones, $`V_m(x)`$. Within the model of point-like scatterers such correlation is entirely absent. Yet even in the case of disorder specified by correlation function (5) it is not difficult, using the definition (8) and the hard-wall model of side boundaries, to ascertain the equality
$$U_{nm}(x)V_m(x^{})=0.$$
(44)
This allows us to couple the inter-mode potentials in Eq. (43) only with each other, not affecting the trial functions $`G_m^{(V)}(x,x_1)`$. The averaging procedure thus makes the operator $`\widehat{𝒯}_n`$ effectively local.
Since equation (16) is one-dimensional in the space variable, the exact Green function $`G_{nn}(x,x^{})`$ can be represented in a form similar to that used for the trial function $`G_n^{(V)}(x,x^{})`$, Eq. (28). Specifically, under weak scattering conditions the function $`G_{nn}(x,x^{})`$, being considered as a function of the first argument ($`x`$) only, is composed of two slightly modulated exponential summands of the appearance $`\varphi _\pm (x)=t_\pm (x)\mathrm{exp}(\pm ik_nx)`$. By applying the operator $`\widehat{𝒯}_n`$ to the functions $`\varphi _\pm (x)`$ one can factor the smooth amplitudes $`t_\pm (x)`$ out of the integral thus arriving at the result
$$\widehat{𝒯}_n\varphi _\pm (x)=\mathrm{\Sigma }(k_n)\varphi _\pm (x),$$
(45)
where
$$\mathrm{\Sigma }(k_n)=\frac{𝒬}{D}_L𝑑x_1W(xx_1)\mathrm{exp}\left[ik_n(xx_1)\right]\underset{\genfrac{}{}{0pt}{}{m=1}{(mn)}}{\overset{\mathrm{}}{}}G_m^{(V)}(x,x_1).$$
(46)
For deriving equation (46), the correlation equality was used
$$U_{nm}(x)U_{kn}(x_1)=\frac{𝒬}{D}W(xx_1)\delta _{mk}$$
(47)
which results from definition (8) and the correlation model (5). The sharp function $`W(xx_1)`$ present in Eq. (46) allows us to replace the trial functions $`G_m^{(V)}(x,x_1)`$ by the unperturbed free mode Green functions $`G_m^{(0)}(x,x_1)`$. Then, in the case of even $`W(x)`$, the factor $`\mathrm{\Sigma }(k_n)`$ acquires the simple form
$$\mathrm{\Sigma }(k_n)=\frac{𝒬}{D}\underset{\genfrac{}{}{0pt}{}{m=1}{(mn)}}{\overset{\mathrm{}}{}}\underset{\mathrm{}}{\overset{\mathrm{}}{}}\frac{dq}{2\pi }\stackrel{~}{W}(q+k_n)\stackrel{~}{G}_m^{(0)}(q).$$
(48)
Here $`\stackrel{~}{G}_m^{(0)}(q)`$ is the Fourier transform of the function $`G_m^{(0)}(x)`$, and it is independent of the sign of the momentum $`k_n`$.
From the above analysis it follows that the action of the operator $`\widehat{𝒯}_n`$ on the function $`G_{nn}(x,x^{})`$ is reduced to multiplying by, in general, the complex-valued quantity $`\mathrm{\Sigma }(k_n)`$. Using the explicit form of $`\stackrel{~}{G}_m^{(0)}(q)`$,
$$\stackrel{~}{G}_m^{(0)}(q)=\frac{1}{(k_m+i0)^2q^2},$$
(49)
for both real and imaginary parts of the mode ‘self-energy’ $`\mathrm{\Sigma }(k_n)=\mathrm{\Delta }k_n^2i/\tau _n^{(\phi )}`$ the expressions are obtained as follows:
$`\mathrm{\Delta }k_n^2`$ $`=`$ $`{\displaystyle \frac{𝒬}{D}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{m=1}{(mn)}}{\overset{\mathrm{}}{}}}𝒫{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{dq}{2\pi }}{\displaystyle \frac{\stackrel{~}{W}(q+k_n)}{k_m^2q^2}},`$ (51)
$`{\displaystyle \frac{1}{\tau _n^{(\phi )}}}`$ $`=`$ $`{\displaystyle \frac{𝒬}{4D}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{m=1}{(mn)}}{\overset{Nc}{}}}{\displaystyle \frac{1}{k_m}}\left[\stackrel{~}{W}(k_nk_m)+\stackrel{~}{W}(k_n+k_m)\right].`$ (52)
The symbol $`𝒫`$ in (51) denotes the principal value. Under the conditions of weak scattering the real part $`\mathrm{\Delta }k_n^2`$ is always small, $`|\mathrm{\Delta }k_n^2|k_n^2`$, so it can be disregarded without serious consequences. At the same time, ‘dissipative’ term (52) plays a crucial role for the further analysis and cannot be omitted. As a result, equation (16) takes the form
$$\left[\frac{^2}{x^2}+\kappa _n^2+i0V_n(x)\mathrm{\Delta }\widehat{𝒯}_n\right]G_{nn}(x,x^{})=\delta (xx^{}),$$
(53)
where $`\kappa _n^2=k_n^2+i/\tau _n^{(\phi )}`$ and $`\mathrm{\Delta }\widehat{𝒯}_n=\widehat{𝒯}_n\widehat{𝒯}_n`$. Thus, for the analysis of equation (53) we shall regard the set of renormalized energies $`\kappa _n^2`$ ($`n=1,2,\mathrm{}`$) as representing the new ‘unperturbed’ spectrum of the system, instead of the original spectrum {$`k_n^2`$}. The perturbation theory can now be developed making use of the appropriate zero-mean potentials $`V_n(x)`$ and $`\mathrm{\Delta }\widehat{𝒯}_n`$.
Note the difference in summation regions for Eqs. (51) and (52). The summation in (52) turns out to be restricted by the number of conducting channels (extended waveguide modes), because only for $`nN_c`$ the disorder-averaged trial Green functions in Eq. (46) are essentially complex-valued (see Eq. (39) at $`\mu =1`$). The level broadening $`1/\tau _n^{(\phi )}`$ implies obligatory presence in the conductor of other extended modes besides the $`n`$th mode itself. In the case of an extremely narrow strip with $`N_c=1`$ the sum (52) contains no terms, and thus the system should exhibit true one-dimensional properties. Specifically, the electrons in such systems can be transferred within two regimes only, ballistic and localized, and the conductance should go down exponentially with the length $`L`$ exceeding the localization length $`\xi _1`$.
On increasing the conductor width, as soon as the wire ceases to be single-mode ($`N_c2`$), the situation changes drastically. The $`n`$th-mode spectrum acquires the level broadening (52) and is subjected to both the potentials $`V_n(x)`$ and $`\mathrm{\Delta }\widehat{𝒯}_n`$. The physical reason for spectrum ‘complexification’ with the availability of more than one extended mode in the conductor is actually the randomization of the spatial phase of the electron wavefunction in going between the states with different mode energies and, consequently, mode momenta. It is just the uncertainty of those momenta that destroys the spatial coherence of one-dimensional quantum waves governed by equation (53). This destruction prevents interferential localization of the true mode states of conducting electrons contrary to their trial states.
To evaluate the ‘phase-breaking’ effect of the term (52) note that at any $`N_c>1`$ the estimate $`1/\tau _n^{(\phi )}𝒬`$ is valid. In particular, in the extreme case of a multimode conductor ($`N_c1`$) with point-like scatterers, replacing the summation in Eq. (52) by an integration we obtain
$$1/\tau _n^{(\phi )}𝒬/4.$$
(54)
It is noteworthy that for a large number of conducting sub-bands each level width becomes independent of the mode number $`n`$ and can thus be thought of as a universal dephasing rate inherent to the 2D conductor in general. Although it is commonly believed that the phase breaking stems from inelastic scattering, the level broadening (52) has nothing to do with such a scattering mechanism. The only important reason for the existence of an imaginary part of the mode energy is that the conductor has to possess more than one extended mode.
Besides the phase breaking contribution $`1/\tau _n^{(\phi )}`$, the role of varying potentials $`V_n(x)`$ and $`\mathrm{\Delta }\widehat{𝒯}_n`$ should also be studied. These zero-mean potentials can be assessed through evaluation of the corresponding Born scattering rates $`1/\tau _n^{(V)}`$ and $`1/\tau _n^{(𝒯)}`$. Estimation of the operator norms $`\widehat{V}_n^2`$ and $`\mathrm{\Delta }\widehat{𝒯}_n^2`$ yields
$`{\displaystyle \frac{\tau _n^{(V)}}{\tau _n^{(𝒯)}}}`$ $``$ $`N_c\text{min}(1,{\displaystyle \frac{L}{N_c\mathrm{}}}),`$ (56)
$`{\displaystyle \frac{\tau _n^{(\phi )}}{\tau _n^{(𝒯)}}}`$ $``$ $`{\displaystyle \frac{1}{\mathrm{cos}^2\vartheta _n}}\text{min}(1,{\displaystyle \frac{L}{N_c\mathrm{}}}).`$ (57)
It hence follows that scattering caused by the non-local potential $`\mathrm{\Delta }\widehat{𝒯}_n`$ is more efficient than that attributed to the local intra-mode potential $`V_n(x)`$.
On the other hand, the inter-mode scattering caused by the operator potential $`\widehat{𝒯}_n`$ in Eq. (16) is already taken largely into account through the imaginary renormalization of the mode energy in (53). From (57) it can be seen that the scattering rate $`1/\tau _n^{(𝒯)}`$ is always small compared to the level width (52) provided that the wire is not extremely long in the $`x`$-direction, i.e. if the length $`L`$ does not fall into the interval $`LN_c\mathrm{}`$. Yet even within this interval the quantity $`1/\tau _n^{(𝒯)}`$ cannot exceed the level broadening. If that is the case, the search for strong Anderson localization at any length of the multimode ($`N_c2`$) conducting strip has no sense even without any inelastic scattering mechanisms. Probably, it is necessary to invent some extra conditions to make the interferential localization possible in such conductors, e.g., magnetic field, surface roughness, etc.
## VI Calculation of the conductance
Equations (53) and (9) provide a way to perform a direct calculation of the disorder-averaged conductance (26). In this section, to avoid cumbersome expressions the results will be given for the case of a large number of conducting channels, $`N_c1`$. Nevertheless, all the estimates, as well as the final formulae, are valid for an arbitrary finite number of modes, $`N_c1`$. In what follows, two summands in equation (26) will be considered separately inasmuch as the calculation techniques and the contribution of these terms in the total conductance differ significantly. The first summand, $`g^{(1)}(L)`$, will be conventionally referred to as the diagonal conductance, since it does not clearly contain any quantities characteristic of the inter-mode scattering. The second term in (26), $`g^{(2)}(L)`$, will be conventionally referred to as the non-diagonal conductance.
Note that the effect of phase breaking, which manifests itself strongly through complexification of the excitation spectrum at $`N_c>1`$, enables one to obtain the solution of equation (53) perturbatively in the potentials $`V_n`$ and $`\mathrm{\Delta }\widehat{𝒯}_n`$, i.e. to neglect, in the leading approximation, the interference of quantum waves multiply scattered by those potentials. It follows from estimates (V) that such an approach is absolutely justified for conductors of length $`LN_c\mathrm{}`$. Yet even in the case of a Q1D wire of length $`LN_c\mathrm{}`$, when the potential $`\mathrm{\Delta }\widehat{𝒯}_n`$ in Eq. (53) cannot be disregarded as follows from estimation (57), the effect of this potential will be shown to be negligibly small.
### A The diagonal conductance
The presence of the local potential $`V_n(x)`$ in Eq. (53) does not substantially complicate the calculation of the conductance. By applying the method of Ref. , this potential can be taken into account with just the same accuracy as was done diagrammatically in Ref. . Accounting for this local potential leads to purely interferential corrections, which we are not concerned with in this paper.
As regards the operator potential $`\mathrm{\Delta }\widehat{𝒯}_n`$, to treat it perturbatively keeping in mind the estimate (V) it is advantageous to go over from the differential equation (53) to the consequent integral one,
$$G_{nn}(x,x^{})=G_{nn}^{(0)}(x,x^{})+\left(\widehat{G}_{nn}^{(0)}\mathrm{\Delta }\widehat{𝒯}_n\widehat{G}_{nn}\right)(x,x^{}),$$
(58)
Here $`G_{nn}^{(0)}(x,x^{})`$ obeys equation (53) with $`\mathrm{\Delta }\widehat{𝒯}_n=0`$. In the leading approximation in the parameter $`L/N_c\mathrm{}1`$, the function $`G_{nn}^{(0)}`$ has the simple ‘unperturbed’ form
$$G_{nn}^{(0)}(x,x^{})=\frac{1}{2ik_n}\mathrm{exp}\left\{\left[ik_n1/l_n^{(\phi )}\right]|xx^{}|\right\}$$
(59)
which nonetheless includes most of the inter-mode-scattering effects. In Eq. (59) the mode coherence length $`l_n^{(\phi )}`$ is associated with the $`n`$th level broadening (52), namely $`l_n^{(\phi )}=2k_n\tau _n^{(\phi )}`$. In the limit $`N_c1`$ its value equals
$$l_n^{(\phi )}=4\mathrm{}\mathrm{cos}\vartheta _n.$$
(60)
Substitution of the Green function (59) into the first summand of Eq. (26) readily gives the diagonal part of the conductance
$$<g^{(1)}(L)>=\underset{n=1}{\overset{N_c}{}}\frac{l_n^{(\phi )}}{L}\left[1\frac{l_n^{(\phi )}}{L}\mathrm{exp}\left(\frac{L}{l_n^{(\phi )}}\right)\mathrm{sinh}\frac{L}{l_n^{(\phi )}}\right].$$
(61)
In the limit $`N_c1`$, by replacing the sum in Eq. (61) with the integral and substituting the coherence length in the form (60), we arrive at the following asymptotic expressions for the conductance (61):
$`L\mathrm{}`$ $`<g^{(1)}(L)>N_c\left(1{\displaystyle \frac{\pi }{12}}{\displaystyle \frac{L}{\mathrm{}}}\right),`$ (63)
$`L\mathrm{}`$ $`<g^{(1)}(L)>\pi N_c\mathrm{}/L.`$ (64)
In the limit of $`LN_c\mathrm{}`$, it follows from (57) that the zero-order approximation in $`\mathrm{\Delta }\widehat{𝒯}_n`$ is, strictly speaking, insufficient for calculating the function $`G_{nn}(x,x^{})`$. Nevertheless, since the scattering rate $`1/\tau _n^{(𝒯)}`$ can only be less than or of the order of the level width $`1/\tau _n^{(\phi )}`$, the appropriate correction to the conductance can be reasonably estimated by substituting into $`g^{(1)}(L)`$ the function $`G_{nn}(x,x^{})`$ obtained from (58) to the first order in $`\mathrm{\Delta }\widehat{𝒯}_n`$. A simple but tedious calculation brings about the following estimation of the corresponding correction to the conductance:
$$<\mathrm{\Delta }g^{(1)}(L)>\left(\frac{\mathrm{}}{L}\right)^2.$$
(65)
This quantity is apparently small compared with (64). The result in (65) indicates that one may account for the inter-mode scattering in the problem under consideration by means of the phase breaking factor which shows itself only in smearing of the mode energy levels.
### B The non-diagonal term of the conductance
When calculating the second term $`<g^{(2)}(L)>`$ in Eq. (26) two essentially different correlators have to be evaluated, the first containing the trial mode Green functions and the second composed of the exact ones. As regards the correlator of the trial functions $`G_m^{(V)}`$, from result (39) it is clear that at $`LN_c\mathrm{}`$ those functions can be replaced by the free ones, independently of the mode. Using the function $`G_{nn}`$ in the form (59) we obtain
$$<g^{(2)}(L)>=\frac{𝒬}{4L^2D}\underset{n=1}{\overset{N_c}{}}\frac{\left(l_n^{(\phi )}\right)^3}{k_n}\mathrm{exp}\left(\frac{L}{l_n^{(\phi )}}\right)\left[\frac{L}{l_n^{(\phi )}}\mathrm{cosh}\frac{L}{l_n^{(\phi )}}\mathrm{sinh}\frac{L}{l_n^{(\phi )}}\right]\underset{\genfrac{}{}{0pt}{}{m=1}{(mn)}}{\overset{N_c}{}}\frac{1}{k_m}.$$
(66)
At $`N_c1`$ this expression is substantially simplified giving the asymptotics as follows:
$`L\mathrm{}`$ $`<g^{(2)}(L)>{\displaystyle \frac{\pi }{24}}N_c{\displaystyle \frac{L}{\mathrm{}}},`$ (68)
$`\mathrm{}LN_c\mathrm{}`$ $`<g^{(2)}(L)>\pi N_c\mathrm{}/2L.`$ (69)
In the case of extremely long conductors with $`LN_c\mathrm{}`$, the trial Green functions in (26) cannot be replaced by unperturbed ones so that the effect of 1D localization should be taken into account properly. The ‘density-current’ correlator standing in (26) was studied in Ref. . The result obtained there closely corresponds to analogous results for the ‘density-density’ and ‘current-current’ correlators. The main feature of the correlator is that it decays exponentially at the localization length $`\xi _m=4L_b^{(V)}(m)`$. Keeping this in mind, it is not difficult to estimate the non-diagonal term in (26) as
$$<g^{(2)}(L)>\left(\frac{N_c\mathrm{}}{L}\right)^2.$$
(70)
## VII Results and discussion
While comparing results (VI B) and (70) with those given by equations (VI A) it can be seen that the non-diagonal part of the conductance is not small relative to the diagonal one only if the conductor length falls within the interval $`\mathrm{}LN_c\mathrm{}`$. In this case the non-diagonal term is half the size of its diagonal counterpart and has the opposite sign. Hence in the whole range of the conductor length (with the width being kept constant) the conductance is described by the following asymptotic expressions:
$$\begin{array}{ccc}\hfill \text{(i)}& L<\mathrm{}:\hfill & <g(L)>N_c\hfill \\ \hfill \text{(ii)}& \mathrm{}LN_c\mathrm{}:\hfill & <g(L)>(\pi /2)N_c\mathrm{}/L1\hfill \\ \hfill \text{(iii)}& N_c\mathrm{}L:\hfill & <g(L)>\pi N_c\mathrm{}/L1.\hfill \end{array}$$
(71)
Strictly speaking, result (71) is valid exactly when the number of channels is large, $`N_c1`$. Nevertheless, even in the case of $`N_c1`$ only an insignificant difference occurs, produced by the dependence of the coherence length $`l_n^{(\phi )}`$ on the mode number $`n`$. This dependence differs somewhat in the non-semi-classical limit from that given by equation (60), but the difference reduces to a numerical factor of the order of unity only.
The result given by Eq. (71) allows one to distinguish three regimes of charge transport in a quantum conductor, depending on its aspect ratio. Regime (i) corresponds to entirely ballistic transport, both from the semi-classical and quantum points of view. The result obtained exhibits a natural stepwise dependence of the conductance on the transverse size of the wire.
In regimes (ii) and (iii) the semi-classical motion should be regarded as diffusive. This opinion is consistent with the conventional view of classical diffusion since the mean free path $`\mathrm{}`$ is small compared with the sample length $`L`$, although it can be in arbitrary relation to the conductor width. Furthermore, such an interpretation is supported by the ohmic, i.e. inversely proportional to $`L`$, dependence of the conductance in both of the indicated regimes. At the same time, it should be particularly emphasized that only in regime (ii), commonly called the weak localization regime, is the result given by the classical kinetic theory reproduced exactly. The expression for the diffusion coefficient in regime (iii) differs from that pertinent to regime (ii) by a factor of two.
Regime (iii) is often called localized, because it is usually supposed that in wires of such a length (commonly referred to as Q1D systems) the Anderson localization should manifest itself to a considerable extent, thus leading to an exponential fall in the conductance. However, from the above analysis it follows that conductors with more than one quantum channel interconnected through scattering mechanisms, even elastic ones, should not exhibit an exponential dependence of kinetic coefficients on the sample length. Such a behaviour is characteristic for the single-mode wires only, which is entirely consistent with theoretical predictions for one-dimensional disordered systems.
To associate findings of this paper with convictions that have prevailed hitherto, it is helpful to examine the electron transport in regimes (ii) and (iii) starting with the trial waveguide states governed by the homogeneous equation (9). Those states are certainly fictitious, they would exist provided the inter-mode scattering were disregarded. If so, the system would indeed represent a set of $`N_c`$ independent one-dimensional conducting channels where the true interferential localization should take place as a result of direct intra-mode scattering from the potentials $`V_n`$. For all of the channels, a hierarchy of localization lengths would exist similar to that representative of equation (40). The length region in (ii) corresponds to the condition when the majority of the trial states are extended. In contrast, in regime (iii) all of those states are localized, which is consistent with the expectation of an exponential fall of the conductance with the growth of the sample length.
In reality one certainly cannot disregard the inter-mode scattering in the case of arbitrary quenched disorder. That scattering results in quite strong coupling of the channels or, rather, the trial mode states. The coupling has shown itself through the complexification of true mode spectrum in Eq. (53). Such a complexification suggests that for any extended mode in a multi-mode strip all other extended modes can be thought of as a phase-breaking reservoir destroying quantum interference and hence strong (exponential) localization. Only weak localization corrections due to the local intra-mode potentials can be detected in the both of the diffusive regimes (71). A comprehensive analysis of the matter is beyond the scope of this paper and will be presented elsewhere.
The existence of different diffusion regimes (ii) and (iii) can be interpreted as the dependence of the diffusion coefficient on the conductor aspect ratio. This dependence can hardly be extracted in the framework of the semi-classical approach. It results from the fact that under a gradual transition from regime (ii) to (iii), with the growth of the conductor length, the trial waveguide states undergo sequential localization. This should reduce the probability for the corresponding excitations to leave the conductor through the current terminals, whereas the probability of their scattering into other extended modes should increase. When all the trial states become finally localized, the diffusion coefficient stabilizes at the value corresponding to regime (iii). The dimensionless conductance of such a long wire is less than unity as a consequence of the conventional Ohm’s law. It seems that this smallness has previously been the reason to suppose all genuine states in Q1D conductors to be localized.
###### Acknowledgements.
The author is grateful to N. M. Makarov for stimulating discussions, A. V. Moroz for help in the interpretation of the results and K. Ilyenko for reading the manuscript.
## A On regularity of the operator $`\widehat{𝖪}`$, Eq. (15)
To ascertain that unity is not among the characteristic numbers of the operator $`\widehat{𝖱}`$, let us take advantage of the operator identity
$$\mathrm{log}det(𝟙\widehat{𝖱})=\text{Tr}\mathrm{log}(𝟙\widehat{𝖱})=\text{Tr}\underset{𝕜=\mathrm{𝟙}}{\overset{\mathrm{}}{}}\frac{(\mathrm{𝟙})^𝕜}{𝕜}\widehat{𝖱}^𝕜.$$
(A1)
The last equality in (A1) presumes that the operator norm is limited by $`\widehat{𝖱}<1`$. This is entirely consistent with the estimate (22).
Consider the traces of the first two terms in sum (A1),
$`\text{Tr}\widehat{𝖱}`$ $`=`$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle _L}𝑑xG_n^{(V)}(x,x)U_{nn}(x),`$ (A2)
$`\text{Tr}\widehat{𝖱}^2`$ $`=`$ $`{\displaystyle \underset{n,m=1}{\overset{\mathrm{}}{}}}{\displaystyle _L𝑑x𝑑x^{}G_n^{(V)}(x,x^{})U_{nm}(x^{})G_m^{(V)}(x^{},x)U_{mn}(x)}.`$ (A3)
Since at $`n>N_c`$ the function $`G_n^{(V)}(x,x^{})`$ in the case of weak scattering has the form (27), it is easy to conclude that the divergence of the logarithm in (A1) can arise from the first term, given in (A2), provided $`U_{nn}(x)0`$. The divergence stems from locality, i.e. coincidence of the arguments, of the Green functions. It manifests itself not only on average but also at a given realization of the random potential entering this term. By separating the diagonal, i.e. intra-mode, potential $`V_n(x)U_{nn}(x)`$ and making the matrix $`U_{nm}`$ off-diagonal we prevent the singularity of the operator $`\widehat{𝖪}`$ given by Eq. (15).
## B Statistical moments of the trial Green functions
To perform the ensemble-averaging of the random function $`\mathrm{\Phi }_\mu (x,x^{}|n)=\left[G_n^{(V)}(x,x^{})\right]^\mu `$, $`\mu \mathrm{}`$, in accordance with representation (28) and (29), we first decompose the function $`G_n^{(V)}(x,x^{})`$ into the sum of four terms each containing narrow packets only of spatial harmonics with phases close to $`\pm k_n(x\pm x^{})`$. In doing so one should use the asymptotic expression for the Wronskian $`𝓦_n`$,
$$𝓦_n=2ik_n\left[\pi _+(x|n)\pi _{}(x|n)+\gamma _+(x|n)\gamma _{}(x|n)\right].$$
(B1)
This results from the assumption of smoothness of the amplitude functions in Eq. (29). Then, after substituting (29) and (B1) into (28), the function $`G_n^{(V)}(x,x^{})`$ can be represented in the form of a scalar product
$$G_n^{(V)}(x,x^{})=\left(\begin{array}{cc}\mathrm{e}^{ik_nx}& ;\mathrm{e}^{ik_nx}\end{array}\right)\left(\genfrac{}{}{0pt}{}{\stackrel{~}{G}_1\stackrel{~}{G}_3}{\stackrel{~}{G}_4\stackrel{~}{G}_2}\right)\left(\genfrac{}{}{0pt}{}{\mathrm{e}^{ik_nx^{}}}{\mathrm{e}^{ik_nx^{}}}\right).$$
(B2)
Here $`\stackrel{~}{G}_i(x,x^{}|n)`$ are the smooth amplitudes constructed from the envelopes $`\pi _\pm (x|n)`$ and $`\gamma _\pm (x|n)`$ as follows:
$`\stackrel{~}{G}_1(x,x^{}|n)`$ $`=`$ $`{\displaystyle \frac{i}{2k_n}}A_n(x)\left[\mathrm{\Theta }_+{\displaystyle \frac{\pi _{}(x^{}|n)}{\pi _{}(x|n)}}\mathrm{\Theta }_{}{\displaystyle \frac{\gamma _+(x^{}|n)}{\pi _+(x|n)}}\mathrm{\Gamma }_{}(x|n)\right],`$ (B4)
$`\stackrel{~}{G}_2(x,x^{}|n)`$ $`=`$ $`{\displaystyle \frac{i}{2k_n}}A_n(x)\left[\mathrm{\Theta }_{}{\displaystyle \frac{\pi _+(x^{}|n)}{\pi _+(x|n)}}\mathrm{\Theta }_+\mathrm{\Gamma }_+(x|n){\displaystyle \frac{\gamma _{}(x^{}|n)}{\pi _{}(x|n)}}\right],`$ (B5)
$`\stackrel{~}{G}_3(x,x^{}|n)`$ $`=`$ $`{\displaystyle \frac{1}{2k_n}}A_n(x)\left[\mathrm{\Theta }_+{\displaystyle \frac{\gamma _{}(x^{}|n)}{\pi _{}(x|n)}}+\mathrm{\Theta }_{}{\displaystyle \frac{\pi _+(x^{}|n)}{\pi _+(x|n)}}\mathrm{\Gamma }_{}(x|n)\right],`$ (B6)
$`\stackrel{~}{G}_4(x,x^{}|n)`$ $`=`$ $`{\displaystyle \frac{1}{2k_n}}A_n(x)\left[\mathrm{\Theta }_{}{\displaystyle \frac{\gamma _+(x^{}|n)}{\pi _+(x|n)}}+\mathrm{\Theta }_+\mathrm{\Gamma }_+(x|n){\displaystyle \frac{\pi _{}(x^{}|n)}{\pi _{}(x|n)}}\right].`$ (B7)
The notation used in Eqs. (B) is
$$A_n(x)=\left[1+\mathrm{\Gamma }_+(x|n)\mathrm{\Gamma }_{}(x|n)\right]^1,\mathrm{\Gamma }_\pm (x|n)=\frac{\gamma _\pm (x|n)}{\pi _\pm (x|n)},\mathrm{\Theta }_\pm =\mathrm{\Theta }[\pm (xx^{})].$$
Before averaging the functions (B) over the random potential, let us note some useful features of the dynamic system (31). Since $`\pi _\pm (x|n)`$ and $`\gamma _\pm (x|n)`$ are the causal functionals of the fields $`\eta _n(x)`$, $`\zeta _{n\pm }(x)`$ and $`\zeta _{n\pm }^{}(x)`$, they are determined by the values of those fields on the intervals $`(x,L/2]`$ and $`[L/2,x)`$ for the functionals labelled by the indexes ($`+`$) and ($``$), correspondingly. The Green function elements (B), which will be subjected to ensemble averaging, are constructed in such a fashion that supports of the random functions entering the functionals of ‘plus’ and ‘minus’ type do not meet. Due to the random fields being effectively $`\delta `$-correlated, see Eqs. (IV) and (37), averaging of the functionals with different sign indexes can be performed independently.
It also follows from equations (31) and conditions (30) that all terms of functional series for $`\pi _\pm (x|n)`$ contain an equal number of fields $`\zeta _{n\pm }`$ and $`\zeta _{n\pm }^{}`$, whereas $`\gamma _\pm (x|n)`$ has an extra functional factor $`\zeta _{n\pm }`$. Since for weak scattering all fields $`\eta _n(x)`$, $`\zeta _{n\pm }(x)`$ and $`\zeta _{n\pm }^{}(x)`$ are approximately Gaussian random processes, only the first summands in square brackets of Eqs. (B4) and (B5) remain non-zero after averaging, while the quantities $`<\stackrel{~}{G}_{3,4}>`$ vanish. By the same arguments, the factor $`A_n(x)`$ in (B) can be replaced by unity.
In view of the statistical independence of the functions $`\eta _n(x)`$ and $`\zeta _{n\pm }(x)`$, it is convenient to average over the real field $`\eta _n(x)`$ already at the initial stage. To that end, it is advantageous to perform the following phase transformation of the amplitudes $`\pi _\pm `$ and $`\gamma _\pm `$:
$`\pi _\pm (x|n)`$ $`=`$ $`\stackrel{~}{\pi }_\pm (x|n)\mathrm{exp}\left[\pm i{\displaystyle _x^{\pm L/2}}\eta _n(x_1)𝑑x_1\right],`$ (B8)
$`\gamma _\pm (x|n)`$ $`=`$ $`\stackrel{~}{\gamma }_\pm (x|n)\mathrm{exp}\left[i{\displaystyle _x^{\pm L/2}}\eta _n(x_1)𝑑x_1\right].`$ (B10)
The new amplitudes $`\stackrel{~}{\pi }_\pm `$ and $`\stackrel{~}{\gamma }_\pm `$ obey the equations
$$\begin{array}{ccc}& & \stackrel{~}{\pi }_\pm ^{}(x|n)\pm \stackrel{~}{\zeta }_{n\pm }^{}(x)\stackrel{~}{\gamma }_\pm (x|n)=0,\hfill \\ & & \stackrel{~}{\gamma }_\pm ^{}(x|n)\pm \stackrel{~}{\zeta }_{n\pm }(x)\stackrel{~}{\pi }_\pm (x|n)=0,\hfill \end{array}$$
(B11)
where the random field $`\stackrel{~}{\zeta }_{n\pm }(x)`$ is related to $`\zeta _{n\pm }(x)`$ by the equality
$$\stackrel{~}{\zeta }_{n\pm }(x)=\zeta _{n\pm }(x)\mathrm{exp}\left[\pm 2i_x^{\pm L/2}\eta _n(x_1)𝑑x_1\right].$$
(B12)
This latter condition does not modify correlation properties of the backscattering fields, Eqs. (IV). Then performing a Fourier transformation of the function $`\mathrm{\Phi }_\mu (x,x^{}|n)`$ over $`x^{}`$, we arrive at the expression conveniently decomposed into the sum of ‘plus’ and ‘minus’ functionals,
$$\stackrel{~}{\mathrm{\Phi }}_\mu (x,q|n)=\left(\frac{1}{2k_n}\right)^\mu \mathrm{e}^{iqx}\left[\stackrel{~}{\mathrm{\Phi }}_\mu ^{(+)}(x,q|n)+\stackrel{~}{\mathrm{\Phi }}_\mu ^{()}(x,q|n)\right].$$
(B13)
Here the functions $`\stackrel{~}{\mathrm{\Phi }}_\mu ^{(\pm }`$ are given by
$$\stackrel{~}{\mathrm{\Phi }}_\mu ^{(\pm )}(x,q|n)=\pm _x^{\pm L/2}𝑑x_1\left[\frac{\stackrel{~}{\pi }_\pm (x_1|n)}{\stackrel{~}{\pi }_\pm (x|n)}\right]^\mu \mathrm{exp}\left[iq(xx_1)+i\mu k_n|xx_1|\pm i\mu _{x_1}^x\eta _n(x_2)𝑑x_2\right].$$
(B14)
Averaging functions (B14) over the random field $`\eta _n(x)`$ with the use of (34) readily yields
$$\stackrel{~}{\mathrm{\Phi }}_\mu ^{(\pm )}(x,q|n)=\pm _x^{\pm L/2}𝑑x_1\left[\frac{\stackrel{~}{\pi }_\pm (x_1|n)}{\stackrel{~}{\pi }_\pm (x|n)}\right]^\mu \mathrm{exp}\left\{iq(xx_1)+\left[i\mu k_n\frac{\mu ^2}{L_f^{(V)}(n)}\right]|xx_1|\right\}.$$
(B15)
To then perform averaging over the fields $`\zeta _{n\pm }(x)`$ it is convenient to use the dynamic equations for the functions $`\stackrel{~}{\mathrm{\Phi }}_\mu ^{(\pm )}(x,q|n)`$ and $`\stackrel{~}{\mathrm{\Gamma }}_\pm (x|n)=\stackrel{~}{\gamma }_\pm (x|n)/\stackrel{~}{\pi }_\pm (x|n)`$. They read
$$\frac{d\stackrel{~}{\mathrm{\Phi }}_\mu ^{(\pm )}(x,q|n)}{dx}=1\left[\frac{\mu ^2}{2L_f^{(V)}(n)}i\mu k_niq\right]\stackrel{~}{\mathrm{\Phi }}_\mu ^{(\pm )}(x,q|n)\mu \stackrel{~}{\zeta }_{n\pm }^{}(x)\stackrel{~}{\mathrm{\Gamma }}_\pm (x|n)\stackrel{~}{\mathrm{\Phi }}_\mu ^{(\pm )}(x,q|n),$$
(B16)
$$\pm \frac{d\stackrel{~}{\mathrm{\Gamma }}_\pm (x|n)}{dx}=\stackrel{~}{\zeta }_{n\pm }(x)+\stackrel{~}{\zeta }_{n\pm }^{}(x)\stackrel{~}{\mathrm{\Gamma }}_\pm ^2(x|n).$$
(B17)
These equations stem from definitions (B14) and system (B11) along with the obvious ‘initial’ conditions
$$\stackrel{~}{\mathrm{\Phi }}_\mu ^{(\pm )}(\pm L/2,q|n)=0,\stackrel{~}{\mathrm{\Gamma }}_\pm (\pm L/2|n)=0.$$
(B18)
Averaging of (B16) with the use of Furutsu-Novikov formula gives the equation
$$\frac{d<\stackrel{~}{\mathrm{\Phi }}_\mu ^{(\pm )}(x,q|n)>}{dx}=1\left[\frac{\mu }{2}\left(\frac{\mu }{L_f^{(V)}(n)}+\frac{1}{L_b^{(V)}(n)}\right)i\mu k_niq\right]<\stackrel{~}{\mathrm{\Phi }}_\mu ^{(\pm )}(x,q|n)>,$$
(B19)
from which the result (39) arises immediately.
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# Measurement Theory and General Relativity
## 1 Introduction
The quantum theory of measurement deals with observers and measuring devices that are all inertial. The universality of gravitational interaction implies, however, that gravitational fields cannot be ignored in general. Moreover, most measurements are performed in laboratories on the Earth, which — among other motions — rotates about its proper axis; in fact, measurements are generally performed by devices and observers that are accelerated. It is therefore necessary to investigate the assumptions that underlie the extension of physics to accelerated systems and gravitational fields. This amounts to a determination of the physical foundations of Einstein’s theory of gravitation inasmuch as this theory is in agreement with all observational data available at present . A critical examination of general relativity from the standpoint of measurement theory leads to certain basic limitations general relativity!limitations that are the main subject of this paper.
## 2 Physical Elements of General Relativity
The basic concepts of general relativity can be uniquely determined starting from the consideration of what observers would measure in physical experiments. This results in the four building blocks of general relativity that are described below.
(i) The fundamental laws of microphysics have been formulated with respect to inertial observers. The measurements of inertial observers in Minkowski spacetime are connected via inhomogeneous Lorentz transformations (i.e. Poincaré transformations). An inertial observer is an observer at rest in an inertial reference system; in fact, such an observer can be thought of as carrying a natural orthonormal tetrad frame $`\lambda _{(\alpha )}^\mu `$ along its worldline. Here $`\lambda _{(0)}^\mu =\mathrm{d}x^\mu /\mathrm{d}\tau `$ is the vector tangent to the worldline (“time axis”) and $`\lambda _{(i)}^\mu `$, $`i=1,2,3`$, are the natural spatial axes of the frame so that $`\lambda _{(\alpha )}^\mu =\delta _\alpha ^\mu `$. Thus Maxwell’s equations in this inertial frame refer to the fields actually measured by these standard observers, i.e. $`F_{\mu \nu }\lambda _{(\alpha )}^\mu \lambda _{(\beta )}^\nu (𝐄,𝐁)`$. One can consider other inertial observers as being at rest in other inertial systems in uniform motion with respect to the original reference system described above. To express the measurements of the other observers, one could transform to their rest frames; alternatively, one could consider physics in the original inertial system and simply describe all measurements with respect to a single system of inertial coordinates $`x^\alpha =(ct,𝐱)`$. In the latter case, which is adopted here for the sake of convenience, one can describe the determination of the electromagnetic field by a moving inertial observer as the projection of the field on the observer’s frame,
$$\widehat{F}_{(\alpha )(\beta )}=F_{\mu \nu }\widehat{\lambda }_{(\alpha )}^\mu \widehat{\lambda }_{(\beta )}^\nu .$$
(1)
Let us now suppose that inertial observers choose to employ arbitrary smooth spacetime coordinates $`x_{}^{}{}_{}{}^{\mu }=x_{}^{}{}_{}{}^{\mu }(x^\alpha )`$. It turns out that — so long as the observers remain inertial — this extension is purely mathematical in nature and can be accomplished without introducing any new physical assumption into the theory. Consider, for instance, the Lorentz force law for a particle of mass $`m`$ and charge $`q`$,
$$m\frac{\mathrm{d}^2x^\mu }{\mathrm{d}\tau ^2}=qF^\mu {}_{\nu }{}^{}\frac{\mathrm{d}x^\nu }{\mathrm{d}\tau }.$$
(2)
Here $`\mathrm{d}\tau `$ is the invariant spacetime interval measured along the path of the particle by the standard inertial observers, i.e. $`\mathrm{d}\tau =c\mathrm{d}t/\gamma `$ and $`\gamma `$ is the Lorentz factor. Assuming the invariance of this interval under the change of coordinates, $`\mathrm{d}\tau ^2=\eta _{\mu \nu }\mathrm{d}x^\mu \mathrm{d}x^\nu =g_{\alpha \beta }^{}\mathrm{d}x_{}^{}{}_{}{}^{\alpha }\mathrm{d}x_{}^{}{}_{}{}^{\beta }`$ with
$$g_{\alpha \beta }^{}=\eta _{\mu \nu }\frac{x^\mu }{x_{}^{}{}_{}{}^{\alpha }}\frac{x^\nu }{x_{}^{}{}_{}{}^{\beta }},$$
(3)
one can simply write equation (2) as
$$m\left[\frac{\mathrm{d}^2x_{}^{}{}_{}{}^{\rho }}{\mathrm{d}\tau ^2}+\mathrm{\Gamma }_{\alpha \beta }^\rho (x^{})\frac{\mathrm{d}x_{}^{}{}_{}{}^{\alpha }}{\mathrm{d}\tau }\frac{\mathrm{d}x_{}^{}{}_{}{}^{\beta }}{\mathrm{d}\tau }\right]=qF_{}^{}{}_{}{}^{\rho }{}_{\sigma }{}^{}\frac{\mathrm{d}x_{}^{}{}_{}{}^{\sigma }}{\mathrm{d}\tau },$$
(4)
with the Christoffel connection
$$\mathrm{\Gamma }_{\alpha \beta }^\rho =\frac{^2x^\mu }{x_{}^{}{}_{}{}^{\alpha }x_{}^{}{}_{}{}^{\beta }}\frac{x_{}^{}{}_{}{}^{\rho }}{x^\mu },$$
(5)
and the auxiliary field variables
$$F_{}^{}{}_{}{}^{\rho \sigma }(x^{})=\frac{x_{}^{}{}_{}{}^{\rho }}{x^\mu }\frac{x_{}^{}{}_{}{}^{\sigma }}{x^\nu }F^{\mu \nu }(x).$$
(6)
In Euclidean space, one can always introduce curvilinear coordinates for the sake of convenience; similarly, one can introduce arbitrary (smooth and admissible) coordinates in Minkowski spacetime. In this way, tensors under the inhomogeneous Lorentz group become tensors under general coordinate transformations.
(ii) To extend measurements to accelerated observersaccelerated observer in Minkowski spacetime, a physical hypothesis is required that would connect the measurement of accelerated and inertial observers. In the standard approach to the theory of relativity, the assumption is that an accelerated observer is at each instant physically equivalent to a hypothetical momentarily comoving inertial observer. Thus an accelerated observer passes through an infinite sequence of such hypothetical inertial observers. Mathematically, this basic assumption is equivalent to replacing a curve by its tangent vector at each point as illustrated in Figure 1.
This assumption is clearly valid for Newtonian point particles, since at each instant the accelerated particle and the momentarily comoving inertial particle have the same state, i.e. the same position and velocity. Moreover, it can be naturally extended to all pointlike phenomena; that is, the assumption is also valid if all phenomena are thought of in terms of pointlike *coincidences* of Newtonian point particles and null rays. However, in more general cases involving intrinsic temporal and spatial scales the above assumption will be referred to as “the hypothesis of locality” . Imagine, for instance, an accelerated measuring device; clearly, it is affected by internal inertial effects. If these inertial effects integrate to a perceptible influence on the outcome of a measurement, the hypothesis of locality is violated. On the other hand, if the timescale of the measurement is so short that the influence of the inertial effects is negligible, then the device is “standard”, i.e. its acceleration can be locally ignored. The hypothesis of locality applied to a clock implies that a standard clock will measure proper time $`\tau `$ along its path; therefore, the hypothesis of locality is the generalization of the “clock hypothesis” to all standard measuring devices . Moreover, the local equivalence of an accelerated observer with an infinite sequence of comoving inertial observers endows the accelerated observer with the continuously varying tetrad system of the inertial observers. This variation can be characterized by a translational acceleration $`𝐠(\tau )`$ and a rotation of the spatial frame with frequency $`𝛀(\tau )`$; alternatively, one may associate acceleration scales (such as $`c^2/g`$ and $`c/\mathrm{\Omega }`$) with the motion of the observer .
The extension of measurements to all observers that can use arbitrary coordinates in Minkowski spacetime implies that one can formulate physical laws in a *generally covariant* form. To extend this covariance further to curved spacetime manifolds, Einstein’s principle of equivalenceequivalence principle is indispensable.
(iii) Einstein’s principle of equivalence embodies the universality of the gravitational interaction and is the cornerstone of general relativity. This principle generalizes a result of Newtonian gravitation that is directly based upon the principle of equivalence of inertial and gravitational masses. Einstein postulated a certain equivalence between an observer in a gravitational field and an accelerated observer in Minkowski spacetime. This heuristic principle, when combined with the hypothesis of locality, implies that an observer in a gravitational field is locally inertial. Thus gravitation has to do with the way local inertial frames are connected to each other. The simplest possibility is through the pseudo-Riemannian curvature of the spacetime manifold; therefore, in general relativity the gravitational field is identified with the spacetime curvature.
(iv) The correspondence between general relativity and Newton’s theory of gravitation is established via the gravitational field equation. That is, within the framework of Riemannian geometry the gravitational field equations are the simplest generalizations of Poisson’s equation, $`^2\mathrm{\Phi }_\text{N}=4\pi G\rho `$, for the Newtonian potential $`\mathrm{\Phi }_\text{N}`$. In general relativity, the Newtonian potential is generalized and replaced by the ten components of the metric tensor $`g_{\mu \nu }`$; similarly, the acceleration of gravity is replaced by the Christoffel connection $`\mathrm{\Gamma }_{\alpha \beta }^\mu `$ and the tidal matrix $`^2\mathrm{\Phi }_\text{N}/x^ix^j`$ is replaced by the Riemann curvature tensor $`R_{\mu \nu \rho \sigma }`$. In Newtonian gravitation, the trace of the tidal matrix is connected to the local density of matter $`\rho `$ by Newton’s constant of gravitation. Similarly, in general relativity the trace of the Riemann tensor is connected to the energy-momentum tensor of matter,
$$R_{\mu \nu }\frac{1}{2}g_{\mu \nu }g^{\alpha \beta }R_{\alpha \beta }=\frac{8\pi G}{c^4}T_{\mu \nu }.$$
(7)
## 3 Measurements of Accelerated Observers
The primary measurements of an observer are those of duration and distance. In general relativity, the hypothesis of locality is indispensable for the interpretation of the results of measurements by accelerated observers. In particular, we define “standard” measuring devices to be those that are compatible with the locality assumption. Thus a standard clock measures proper time along its trajectory; similarly, a standard measuring rod is usually assumed to provide a proper measure of distance. At each instant of time, the accelerated observer is momentarily equivalent to a hypothetical comoving inertial observer; therefore, both observers have the instantaneous Euclidean space in common. It would appear then that placing standard measuring rods one next to the other and so on should lead to the proper measurement of spatial distances by accelerated observers.
An important issue is the extent to which such measurements of time and distance can lead to the establishment of an admissible coordinate system around the accelerated observer. Well-known investigations have led to the result that such coordinate systems have limited spatial extent given by the acceleration lengths (e.g. $`c^2/g`$ and $`c/\mathrm{\Omega }`$), since these are the only length scales in the problem. The method of construction of accelerated coordinate systems could even be nonlocal; however, limitations would still exist as recently pointed out by Marzlin . It might therefore appear that (local and nonlocal) coordinate systems could in general be constructed in a cylindrical region around the worldline of the accelerated observer. However, this conclusion is ultimately based upon the use of standard measuring rods whose existence turns out to be in conflict with the hypothesis of locality. A fundamental problem associated with length measurements is the following: a standard measuring rod, however small, has nevertheless a nonzero spatial extent whereas the hypothesis of locality is only pointwise valid. This implies a rather basic limitation on the measurement of length by accelerated observers and can be illustrated by the following thought experiment. Imagine two observers $`O_1`$ and $`O_2`$ at rest in an inertial frame. For $`t0`$, their coordinates $`x^\alpha =(ct,𝐱)`$ are $`(ct,0,0,0)`$ and $`(ct,L,0,0)`$, respectively. At $`t=0`$, they are accelerated from rest along the $`x`$-direction *in exactly the same way* so that at time $`t>0`$ each has a velocity $`𝐯=v\widehat{𝐱}`$. The distance between $`O_1`$ and $`O_2`$ as measured by observers at rest in the inertial frame is always $`L`$, since
$$x_1(t)=_0^tvdt\text{and}x_2(t)=L+_0^tvdt$$
(8)
for $`t>0`$ and $`x_2(t)x_1(t)=L`$. What is the distance between $`O_1`$ and $`O_2`$ as measured by comoving observers? It turns out that the hypothesis of locality provides a unique answer to this question only in the limit $`L0`$. To show this, let us first note that at a given time $`\widehat{t}>0`$, $`O_1`$ and $`O_2`$ have the same speed $`\widehat{v}=c\beta `$. The hypothesis of locality implies that the accelerated observers pass through an infinite sequence of momentarily comoving inertial observers. Thus imagine the Lorentz transformation between the inertial frame $`x^\alpha =(ct,𝐱)`$ and the “instantaneous” inertial rest frame $`x_{}^{}{}_{}{}^{\alpha }=(ct^{},𝐱^{})`$ of the observers at $`\widehat{t}`$ given by
$$c(t\widehat{t})=\gamma (ct^{}+\beta x^{}),x\widehat{x}=\gamma (x^{}+c\beta t^{}),y=y^{},z=z^{},$$
(9)
where $`\gamma =(1\beta ^2)^{1/2}`$ is the Lorentz factor at $`\widehat{t}`$. The events with coordinates $`O_1:(c\widehat{t},x_1,0,0)`$ and $`O_2:(c\widehat{t},x_2,0,0)`$ in the original inertial frame have coordinates $`O_1:(ct_{}^{}{}_{1}{}^{},x_1^{},0,0)`$ and $`O_2:(ct_{}^{}{}_{2}{}^{},x_2^{},0,0)`$ in the instantaneous inertial frame. It follows from the Lorentz transformation (9) that $`L^{}=x_2^{}x_1^{}=\gamma L`$. This has a simple physical interpretation: The Lorentz-FitzGerald contracted distance between $`O_1`$ and $`O_2`$ is always $`L`$, hence the “actual” distance between $`O_1`$ and $`O_2`$ must be larger by the Lorentz $`\gamma `$-factor. One can imagine that the distance between $`O_1`$ and $`O_2`$ is populated by a large number of hypothetical accelerated observers moving in exactly the same way as $`O_1`$ and $`O_2`$ and carrying infinitesimal measuring rods that are placed side by side to measure the distance under consideration.
It must be equally correct to replace the infinite sequence of inertial systems $`x_{}^{}{}_{}{}^{\alpha }=(ct^{},𝐱^{})`$ by a continuously moving frame. To this end, we must choose the worldline $`\overline{x}^\mu (\tau )`$ of one of the accelerated observers — such as $`O_1`$, $`O_2`$, or any of the hypothetical observers in between the two — and note that at any instant of proper time $`\tau `$ along the worldline, this fiducial observer is in a Euclidean space with Cartesian coordinates $`𝐗`$ in accordance with the hypothesis of locality. The connection between the coordinates $`x^\mu `$ in the original inertial frame and the new coordinates $`X^\mu `$ is given by $`X^0=\tau `$ and
$$x^\mu =\overline{x}^\mu (X^0)+X^i\overline{\lambda }_{(i)}^\mu ,$$
(10)
where $`\overline{\lambda }_{(i)}^\mu `$ is the natural tetrad frame along the worldline of the reference observer. Specifically, the fiducial observer is instantaneously inertial by the hypothesis of locality and hence assigns coordinates $`X^0=\tau `$ and $`X^i=\sigma \xi _\mu \overline{\lambda }_{(i)}^\mu `$ to spacetime events. Here $`\xi ^\mu `$ is a unit spacelike vector normal to $`\overline{\lambda }_{(0)}^\mu `$ at $`\overline{x}^\mu (\tau )`$ along a straight line that connects $`\overline{x}^\mu (\tau )`$ to an event with coordinates $`x^\mu `$ in the original background inertial frame, $`\xi _\mu \overline{\lambda }_{(i)}^\mu `$ are direction cosines and $`\sigma =\left|𝐗\right|`$ is the proper length of this spacelike line segment. To develop this approach further, it is necessary to specify the motion explicitly. Thus we assume that $`O_1`$ and $`O_2`$ are uniformly accelerated with acceleration $`g`$ and we choose $`O_1`$ to be the fiducial observer. The natural orthonormal nonrotating tetrad frame along the worldline of $`O_1`$ is given by
$`\overline{\lambda }_{(0)}^\mu `$ $`=`$ $`(\gamma ,\beta \gamma ,0,0),`$ (11)
$`\overline{\lambda }_{(1)}^\mu `$ $`=`$ $`(\beta \gamma ,\gamma ,0,0),`$ (12)
$`\overline{\lambda }_{(2)}^\mu `$ $`=`$ $`(0,0,1,0),`$ (13)
$`\overline{\lambda }_{(3)}^\mu `$ $`=`$ $`(0,0,0,1),`$ (14)
just as for the Lorentz transformation (9). Then the inertial frame $`x^\alpha =`$ $`(ct,x,y,z)`$ and the Fermi frame $`X^\alpha =(cT,X,Y,Z)`$ are connected by
$`ct`$ $`=`$ $`\left(X+{\displaystyle \frac{c^2}{g}}\right)\mathrm{sinh}\left(gT/c\right),`$ (15)
$`x`$ $`=`$ $`\left(X+{\displaystyle \frac{c^2}{g}}\right)\mathrm{cosh}\left(gT/c\right){\displaystyle \frac{c^2}{g}},`$ (16)
$`y=Y`$ and $`z=Z`$. The spatial origin of the new coordinate system is occupied by $`O_1`$ such that $`\overline{x}^\mu (\tau )=(\beta \gamma ,\gamma 1,0,0)`$, where $`\beta =\mathrm{tanh}(\tau _1/)`$, $`\gamma =\mathrm{cosh}(\tau _1/)`$, $`=c^2/g`$ is the acceleration length and $`\tau _1`$ is the proper time along $`O_1`$. As before, at any given time $`\widehat{t}>0`$ the events $`O_1:(c\widehat{t},x_1,0,0)`$ and $`O_2:(c\widehat{t},x_2,0,0)`$ now correspond to $`O_1:(\tau _1,X_1,0,0)`$ and $`O_2:(\tau _2,X_2,0,0)`$, where $`x_2x_1=L`$ and $`X_1=0`$ by construction. The distance between $`O_1`$ and $`O_2`$ in this Fermi frame is then given by $`L_\text{F}=X_2X_1=X_2`$. It follows from equations (15) and (16) that
$`c\widehat{t}`$ $`=`$ $`\mathrm{sinh}(\tau _1/),`$ (17)
$`x_1`$ $`=`$ $`\left[\mathrm{cosh}(\tau _1/)1\right],`$ (18)
$`c\widehat{t}`$ $`=`$ $`(X_2+)\mathrm{sinh}(\tau _2/),`$ (19)
$`x_2`$ $`=`$ $`(X_2+)\mathrm{cosh}(\tau _2/).`$ (20)
Equations (19) and (20) can be written as
$$(X_2+)^2=(x_2+)^2c^2\widehat{t}^2,$$
(21)
where $`x_2=x_1+L`$ and $`x_1`$ and $`\widehat{t}`$ are given by equations (18) and (17), respectively. Thus one finds that
$$L_\text{F}=\left[(1+2ϵ\gamma +ϵ^2)^{1/2}1\right],$$
(22)
where $`ϵ=L/=gL/c^2`$ and $`\gamma =(1+g^2\widehat{t}^2/c^2)^{1/2}`$. The length in the Fermi frame $`L_\text{F}`$ must be compared with the corresponding result from the instantaneous Lorentz frame $`L^{}=\gamma L`$; indeed, the ratio $`L_\text{F}/L^{}`$ approaches unity only in the limit $`ϵ0`$. This is a remarkable result that has far-reaching consequences. Let us note that for $`ϵ1`$,
$$L_\text{F}/L^{}1\frac{1}{2}\beta ^2\gamma ϵ$$
(23)
to first order in $`ϵ`$; however, over a long time $`c/g`$ the quantity $`\beta ^2\gamma ϵ`$ may not remain small compared to unity. Moreover, $`L_\text{F}/L^{}0`$ as $`g\widehat{t}/c\mathrm{}`$ and hence $`\gamma \mathrm{}`$. It follows from these considerations that consistency is achieved for $`\gamma ϵ0`$; hence, the acceleration length and time, i.e. $`c^2/g`$ and $`c/g`$, respectively, place severe limitations on the domain of applicability of the hypothesis of locality. Furthermore, let us suppose that the Fermi frame is established along $`O_2`$ instead of $`O_1`$. Then the resulting distance would be different from $`L_\text{F}`$; however, all such lengths agree in the $`ϵ0`$ limit.
It is interesting to mention here another measure of distance from $`O_1`$ and $`O_2`$ using light signals. Let $`O_1`$ send a signal at $`\tau _1^{}`$ that reaches $`O_2`$ at $`\tau _2`$ and is immediately returned to $`O_1`$. The return signal reaches $`O_1`$ at $`\tau _1^+`$, where $`\tau _2=(\tau _1^{}+\tau _1^+)/2`$. Observer $`O_1`$ would then determine the distance to $`O_2`$ via $`L_{\text{ph}}=c(\tau _1^+\tau _1^{})/2`$, which works out to be
$$L_{\text{ph}}=\mathrm{ln}(1+L_\text{F}/).$$
(24)
It is clear by symmetry that if $`O_2`$ initiates a light signal to $`O_1`$, etc., then the resulting light travel time would be different, since in equation (24) the Fermi length would be the one determined on the basis of $`O_2`$ as the fiducial observer. Nevertheless, for $`\gamma ϵ`$ negligibly small all these length measurements agree with each other.
The simple example that has been worked out here can be generalized to arbitrary but identical velocity for $`O_1`$ and $`O_2`$. The comparison of the instantaneous local inertial frame with the continuously moving geodesic frame leads to the conclusion that the basic length and time scales under consideration must in general be negligible compared to the relevant acceleration scales. This has significant consequences for the comparison of theory and experiment in general relativity ; in particular, the physical significance of Fermi coordinates is in general further limited to the immediate neighborhood of the observer and wave equations are meaningful only within this domain.
It follows from these considerations that the physical dimensions of any standard measuring device must be negligible compared to the relevant acceleration length $``$ and the duration of the measurement must in general be negligible compared to $`/c`$. These are not significant limitations for typical accelerations in the laboratory; for instance, for the Earth’s acceleration of gravity $`c^2/g1`$lyr. Moreover, observers at rest on the Earth typically refer their measurements to rotating Earth-based coordinates; hence, this coordinate system is mathematically valid up to a “light cylinder” at a radius of $`=c/\mathrm{\Omega }28`$AU. But physically valid length measurements can extend over a neighborhood of the observer with a radius much smaller than $``$. In fact, this “light cylinder” has no bearing on astronomical observations, since observers simply take into account the absolute rotation of the Earth and reduce astronomical data by taking due account of aberration and Doppler effects.
The standard “classical” measuring device of mass $`\mu `$ has wave characteristics, given by its Compton wavelength $`\mathrm{}/\mu c`$ and period $`\mathrm{}/\mu c^2`$, that must be negligible in comparison with the scales of length and time that characterize the device as a consequence of the quasi-classical approximation. For instance, a clock of mass $`\mu `$ must have a resolution exceeding $`\mathrm{}/\mu c^2`$; similarly, the mass of a clock with resolution $`\theta `$ must exceed $`\mathrm{}/\theta c^2`$. These assertions follow from the application of the uncertainty principle to measurements performed by a standard device . When such quantum limitations are combined with the classical limitations discussed above, on finds that $`\mathrm{}/\mu c`$; therefore, the translational acceleration of a standard classical measuring device must be much less than $`\mu c^3/\mathrm{}`$ and its rotational frequency must be much less than $`\mu c^2/\mathrm{}`$. The idea of the existence of a maximal proper acceleration is due to Caianiello .
## 4 Measurements in Gravitational Fields
The physical results of the previous section can be extended to local measurements in a gravitational field via an interpretation of the Einstein principle of equivalenceequivalence principle—( in terms of the gravitational Larmor theorem. Larmor theorem!gravitational—(
A century ago, Larmor established a local equivalence between magnetism and rotation for all particles with the same charge to mass ratio $`(q/m)`$. That is, charged particle phenomena in a magnetic field correspond to those in a frame rotating with the Larmor frequency $`𝛀_\text{L}=q𝐁/2mc`$. This local relation is valid to first order in field strength for slowly varying fields and slowly moving charged particles. Such a correspondence also exists for electric fields and linearly accelerated frames. It turns out that Larmor’s theorem can be generalized in a natural way to the case of gravitational fields.
The close analogy between Coulomb’s law of electricity and Newton’s law of gravitation leads to an interpretation of Newtonian gravity in terms of nonrelativistic theory of the gravitoelectric field. Moreover, any theory that combines Newtonian gravity with Lorentz invariance in a consistent manner is expected to contain a gravitomagnetic field as well. In fact, in general relativity the exterior spacetime metric for a rotating mass may be expressed in the linear approximation as
$$\mathrm{d}s^2=c^2(1\frac{2}{c^2}\mathrm{\Phi }_\text{N})\mathrm{d}t^2+(1+\frac{2}{c^2}\mathrm{\Phi }_\text{N})\delta _{ij}\mathrm{d}x^i\mathrm{d}x^j\frac{4}{c}(𝐀_g\mathrm{d}𝐱)\mathrm{d}t,$$
(25)
where $`\mathrm{\Phi }_\text{N}=GM/r`$ is the Newtonian potential and $`𝐀_g=G𝐉\times 𝐫/cr^3`$ is the gravitomagnetic vector potential. The gravitoelectric and gravitomagnetic fields are then given by $`𝐄_g=\mathrm{\Phi }_\text{N}`$ and $`𝐁_g=\times 𝐀_g`$, respectively.
It is possible to formulate a gravitational Larmor theorem by postulating that the gravitoelectric and gravitomagnetic charges are given by $`q_E=m`$ and $`q_B=2m`$, respectively. In fact, $`q_B/q_E=2`$ since gravitation is a spin-2 field. Thus $`𝛀_\text{L}=𝐁_g/c`$, which is consistent with the fact that an ideal gyroscope at a given position in space would precess in the gravitomagnetic field with a frequency $`𝛀_\text{P}=𝐁_g/c`$. The general form of the gravitational Larmor theorem Larmor theorem!gravitational—) is then an interpretation of the Einstein principle of equivalenceequivalence principle—) for linear gravitational fields in a finite neighborhood of an observer; for instance, in the gravitational field of the Earth an observer can be approximately inertial within the “Einstein elevator” if the “elevator” falls freely with acceleration $`gGM/r^2`$ while rotating with frequency $`\mathrm{\Omega }GJ/c^2r^3`$. It follows that the relevant gravitoelectromagnetic acceleration lengths are given by $`c^2/g`$ and $`c/\mathrm{\Omega }`$ in this case and the restrictions discussed in the previous section would then apply to the measurements of an observer in a gravitational field as well. These limitations are generally expected to be important for the post-Newtonian corrections of high order in relativistic gravitational systems.
General relativity has found applications mostly in astronomical systems, where Newtonian results have been extended to the relativistic domain. In particular, small post-Newtonian corrections are usually included in the equations of motion. Suppose, for instance, that one is interested in the distance between the members of a relativistic binary system. It follows from our considerations that such a length — which corresponds in the Newtonian theory to the Euclidean distance — may not be well defined. However, the resulting discrepancy could be masked by other parameters; that is, this circumstance may be difficult to ascertain experimentally since the comparison of data with the theory generally involves parameters that are not independently available and whose particular values need to be determined from the data.
Let us next consider tidal accelerations within the “Einstein elevator”. For a device of dimension $`\widehat{\delta }`$, the tidal acceleration $`\widehat{g}`$ is given by the Jacobi equation and can be estimated by $`\widehat{g}K\widehat{\delta }`$, where $`K`$ is a typical component of the tidal matrix $`K_{ij}=c^2R_{\mu \nu \rho \sigma }\lambda _{(0)}^\mu \lambda _{(i)}^\nu \lambda _{(0)}^\rho \lambda _{(j)}^\sigma `$. According to the results of the previous section $`\widehat{g}\mu c^3/\mathrm{}`$, where $`\mu `$ is the mass of the device. Imagine, for instance, such a device on a star of mass $`M`$ and radius $`R`$ that is undergoing “complete” spherical gravitational collapse. In this case , $`KGM/R^3`$ and $`\widehat{\delta }c^2/\widehat{g}`$ imply that $`\widehat{\delta }^2c^2R^3/GM`$. On the other hand, the requirements that $`\mu M`$ and $`\widehat{\delta }\mathrm{}/\mu c`$ result in
$$R^3\frac{GM}{c^2}\left(\frac{\mathrm{}}{Mc}\right)^2=\frac{\mathrm{}}{Mc}L_\text{P}^2,$$
(26)
where $`L_\text{P}=(\mathrm{}G/c^3)^{1/2}`$ is the Planck length ($`10^{33}`$cm) that is the geometric mean of the gravitational radius $`GM/c^2`$ and the Compton wavelength $`\mathrm{}/Mc`$ for any physical system . Thus collapse to a classical point singularity is meaningless on the basis of these considerations.
## 5 Wave Phenomena
Classical waves have intrinsic scales and are thus expected to be in conflict with the hypothesis of locality; indeed, for an electromagnetic wave of (reduced) wavelength $`\mathrm{¯}\lambda `$ the expected deviation from the hypothesis of locality is expected to be of the form $`\mathrm{¯}\lambda /`$. More specifically, let us consider the problem of determination of the period of an incident electromagnetic wave by an accelerated observer. The observer needs to measure at least a few oscillations of the wave before a reasonable determination of the period can be made; therefore, the curvature of the observer’s worldline cannot be neglected unless $`\mathrm{¯}\lambda /`$ is too small to be observationally significant. It follows that the instantaneous Doppler and aberration formulas are in general valid only in the eikonal limit $`\mathrm{¯}\lambda /0`$. The issues involved here can be illustrated by a simple thought experiment. Let us consider an observer rotating with uniform speed $`c\beta `$ and frequency $`\mathrm{\Omega }`$ in the positive sense around the origin on a circle of radius $`r=c\beta /\mathrm{\Omega }`$ in the $`(x,y)`$-plane. A plane electromagnetic wave of frequency $`\omega `$ is incident along the $`z`$-axis and the rotating observer measures its frequency. According to the hypothesis of locality, the observer is at each instant momentarily inertial and hence $`\omega ^{}=\gamma \omega `$ according to the transverse Doppler effect.Doppler!effect This is illustrated in Figure 2.
On the other hand, if we assume that the hypothesis of locality applies to the field measurement,
$$F_{(\alpha )(\beta )}(\tau )=F_{\mu \nu }\lambda _{(\alpha )}^\mu \lambda _{(\beta )}^\nu ,$$
(27)
and the instantaneously determined electromagnetic field $`F_{(\alpha )(\beta )}(\tau )`$ is then Fourier analyzed over proper time — which is definitely a nonlocal procedure — to determine its frequency content, then we find that $`\omega ^{}=\gamma (\omega \mathrm{\Omega })`$. Thus $`\omega ^{}=\gamma \omega (1\mathrm{¯}\lambda /)`$, where $`=c/\mathrm{\Omega }`$; hence, the instantaneous Doppler result is recovered for $`\mathrm{¯}\lambda 0`$. The upper (lower) sign here refers to right (left) circularly polarized incident wave. Apart from the Lorentz factor $`\gamma `$ that refers to the time dilation involved here, the result for $`\omega ^{}`$ has a simple physical interpretation: The electromagnetic field rotates with frequency $`\omega `$ ($`\omega `$) about the $`z`$-axis for an incident right (left) circularly polarized wave, so that the field rotates with respect to the observer with frequency $`\omega \mathrm{\Omega }`$ ($`\omega \mathrm{\Omega }`$). Thus the helicity of the radiation couples to the rotation of the observer, i.e. $`\mathrm{}\omega ^{}=\gamma (\mathrm{}\omega 𝐬𝛀)`$; in fact, this is an example of the general phenomenon of spin-rotation coupling . For instance, for experiments on the Earth the “nonrelativistic” Hamiltonian for a spin-$`\frac{1}{2}`$ particle should be supplemented by
$$_{\text{SR}}=𝐬𝛀+𝐬𝛀_\text{P},$$
(28)
where $`𝛀`$ is the frequency of Earth’s rotation and $`𝛀_\text{P}`$ is the gravitomagnetic precession frequency. The second term in equation (28) illustrates the gravitational Larmor theorem.Larmor theorem!gravitational It is interesting to note that $`\mathrm{}\mathrm{\Omega }10^{19}`$eV and $`\mathrm{}\mathrm{\Omega }_\text{P}10^{29}`$eV; in fact, recent experiments have demonstrated the existence of the first term in (28). Moreover, the position dependence of the second term in (28) indicates the existence of a gravitomagnetic Stern-Gerlach force $`(𝐬𝛀_\text{P})`$ that is purely spin dependent and violates the universality of free fall. For instance, neutrons in different spin states in general fall differently in the gravitational field of a rotating mass; similarly, the gravitational deflection of polarized light is affected by the rotation of the mass. That is, in addition to, and about, the Einstein deflection angle $`\mathrm{\Delta }=4GM/c^2D`$, there is a splitting due to the helicity-rotation coupling by a much smaller angle $`\delta =4\mathrm{¯}\lambda GJ/c^3D^3`$, where $`D`$ is the impact parameter for radiation propagating normal to the rotation axis and over a pole of the rotating mass . As $`\mathrm{¯}\lambda /0`$, $`\delta 0`$ and hence the standard result for a null geodesic is recovered.
To explain all of the experimental tests of general relativity, it is sufficient to consider all wave phenomena only in the JWKB limit. That is, geometric “optics” is all that is required; no gravitational effect involving wave “optics” has ever been detected thus far. An interesting opportunity for detecting such effects would come about if the quasinormal modes (QNMs) of black holes could be observed. The infinite set of QNMs corresponds to damped oscillations of a black hole that come about as the black hole divests itself of the energy of the external perturbation and returns to a stationary state; therefore, these ringing modes of black holes appear as $`𝒜\mathrm{exp}(i\omega t)`$ at late times far from a black hole. Here $`𝒜`$ is the amplitude of the oscillation that depends on the strength of the perturbation as well as the black hole response, while $`\omega =\omega _0i\mathrm{\Gamma }`$ with $`\mathrm{\Gamma }0`$ is purely a function of mass $`M`$, angular momentum $`J`$ and charge $`Q`$ of the black hole, i.e. $`\omega =\omega _{jmn}(M,J,Q)`$, where $`j`$, $`m`$ and $`n`$ are parameters characterizing the total angular momentum of the radiation field, its component along the $`z`$-axis and the mode number, respectively . The mode number $`n=0,1,2,\mathrm{}`$, generally refers to the fundamental, first excited state, etc., of the perturbed black hole with $`j`$ and $`m`$; in fact, $`\mathrm{\Gamma }`$ increases with $`n`$ so that the higher excited states are more strongly damped. The fundamental least-damped gravitational mode with $`j=2`$ and $`n=0`$ for a Schwarzschild black hole is given by
$`\omega _0/2\pi `$ $``$ $`10^4(M_{}/M)\text{Hz},`$ (29)
$`\mathrm{\Gamma }^1`$ $``$ $`6\times 10^5(M/M_{})\text{sec},`$ (30)
so that even this mode is rather highly damped and would therefore be very difficult to observe. The damping problem improves by an order of magnitude if the black hole rotates rapidly; however, the observational difficulties would still be considerable. The observation of such a mode would be very significant physically since, among other things, near an oscillating black hole $`_gGM/c^2`$ and with $`\mathrm{¯}\lambda =c/\omega _0`$, we have $`\mathrm{¯}\lambda /_g1`$, so that wave “optics” can be explored in the gravitational field of a black hole.
It is necessary to examine the justification for the local field assumption (27), since it leads — in the thought experiment of Figure 2 — to the result that a normally incident right circularly polarized wave with $`\omega =\mathrm{\Omega }`$ would stand completely still with respect to the observer. This circumstance is in contradiction with expectations based on elementary notions of relativity theory . In fact, at $`\omega =\mathrm{\Omega }`$ one has $`\mathrm{¯}\lambda /=1`$ and it is possible to argue that the hypothesis of locality must be violated. To this end, imagine an accelerated charged particle in the nonrelativistic approximation. The particle radiates electromagnetic waves with characteristic wavelength $`\mathrm{¯}\lambda `$; therefore, it is expected that such a particle would not be locally inertial and that (27) is violated. Indeed, the equation of motion of the particle is given by
$`m{\displaystyle \frac{\mathrm{d}^2𝐱}{\mathrm{d}t^2}}{\displaystyle \frac{2}{3}}{\displaystyle \frac{q^2}{c^3}}{\displaystyle \frac{\mathrm{d}^3𝐱}{\mathrm{d}t^3}}+\mathrm{}=𝐟.`$ (31)
The radiation reaction term — due originally to Abraham and Lorentz — ensures that the particle is not pointwise inertial, since its position and velocity are not sufficient to determine the state of the radiating particle.
These classical considerations must naturally extend to the quantum domain as well, since quantum theory is based on the notion of wave-particle duality. That is, we expect that the hypothesis of locality would be violated in the quantum regime. Consider, for instance, the determination of muon lifetime by Bailey *et al.* involving muons (in a storage ring at CERN) undergoing centripetal acceleration of $`g=\gamma ^2v^2/r10^{21}\text{cm}\text{sec}^2`$. If $`\tau _\mu ^0`$ is the lifetime of the muon at rest, then the hypothesis of locality would imply that the lifetime in the storage ring would be $`\tau _\mu =\gamma \tau _\mu ^0`$. In the experiment, $`r7`$m, $`\gamma 29`$ and time dilation is verified at the level of $`10^3`$. On the other hand, the deviation from the hypothesis of locality is expected to be of the form $`\mathrm{¯}\lambda /10^{13}`$, where $`\mathrm{¯}\lambda =\mathrm{}/mc`$ is the Compton wavelength of the muon and $`=c^2/g1`$cm is the translational acceleration length. But the functional form of this deviation is not specified by our general intuitive considerations. In any case, $`\mathrm{¯}\lambda /`$ is about ten orders of magnitude below the level of experimental accuracy . In fact, the decay of the muon has been considered in this case by Straumann and Eisele by replacing the accelerated muon by the stationary state of a muon in a Landau level with very high quantum number . It can be shown that the decay of such a state results in
$`\tau _\mu \gamma \tau _\mu ^0\left[1+{\displaystyle \frac{2}{3}}(\mathrm{¯}\lambda /)^2\right],`$ (32)
so that the deviation from the hypothesis of locality is very small ($`10^{25}`$) in this case but definitely nonzero.
## 6 Discussion
General relativity is a consistent theory of pointlike coincidences involving classical point particles and rays of radiation. The theory is robust and can be naturally extended to include wave phenomena (“minimal coupling”); however, general relativity is expected to have limited significance in this regime. From a basic standpoint, the main difficulty is the hypothesis of locality.
An accelerated observer passes through a continuous infinity of hypothetical inertial observers; therefore, the most general linear connection between the field measured by the accelerated observer $`_{\alpha \beta }(\tau )`$ and the locally measured field $`F_{\alpha \beta }^{}(\tau )=F_{(\alpha )(\beta )}(\tau )`$ that is consistent with causality is
$$_{\alpha \beta }(\tau )=F_{\alpha \beta }^{}(\tau )+_0^\tau 𝒦_{\alpha \beta }{}_{}{}^{\gamma \delta }(\tau ,\tau ^{})F_{\gamma \delta }^{}(\tau ^{})\mathrm{d}\tau ^{}.$$
(33)
Here the observer is inertial for $`\tau 0`$ and the absence of the kernel $`𝒦`$ would be equivalent to the hypothesis of locality; moreover, if $`𝒦`$ is directly connected with acceleration, then the deviation from the hypothesis of locality is generally of order $`\mathrm{¯}\lambda /`$. Assuming that $`𝒦`$ is a convolution-type kernel (i.e. it depends only on $`\tau \tau ^{}`$), it is possible to determine $`𝒦`$ uniquely based on the assumption that no observer can ever stay at rest with respect to a basic radiation field. This is simply a generalization of the well-known result of Lorentz invariance, so that the *motion* of an electromagnetic wave would then become independent of the observer. We extend the observer independence of wave notion to all basic radiation fields and elevate this notion to the status of a fundamental physical principle . Writing equation (27) as $`F^{}=\mathrm{\Lambda }F`$, our basic assumption implies that the *resolvent* kernel $``$ is given by
$$=\frac{\mathrm{d}\mathrm{\Lambda }(\tau )}{\mathrm{d}\tau }\mathrm{\Lambda }^1(0).$$
(34)
It follows that for a scalar field ($`\mathrm{\Lambda }=1`$), $`=0`$ and hence $`𝒦=0`$; therefore, an observer can in principle stay at rest with respect to a scalar field. This is contrary to our basic assumption, which then excludes fundamental scalar fields. In this way, a nonlocal theory of accelerated observers has been developed that is in agreement with all available observational data . Moreover, novel inertial effects are predicted by the nonlocal theory. For instance, let us recall the thought experiment (cf. Figure 2) involving plane electromagnetic radiation of frequency $`\omega `$ normally incident on an observer rotating counterclockwise with $`\mathrm{\Omega }\omega `$; the nonlocal theory predicts that the field amplitude measured by the observer is larger by a factor of $`1+\mathrm{\Omega }/\omega `$ for positive helicity radiation and smaller by a factor of $`1\mathrm{\Omega }/\omega `$ for negative helicity radiation. For radio waves with $`\mathrm{¯}\lambda 1\text{cm}`$ and an observer rotating at a frequency of 50 Hz, we have $`\mathrm{\Omega }/\omega =\mathrm{¯}\lambda /10^8`$.
Finally, it should be mentioned that no thermal ambience is encountered for an accelerated observer on the basis of the approach adopted in this paper. This is consistent with the absence of any experimental evidence for such a thermal ambience at present . That is, either (27) or (33) can be used to determine the quantum radiation field according to an accelerated observer once the quantum field in the inertial frame is given. Indeed, the nonlocal theory has been developed based on the assumption that no quanta are created or destroyed merely because an observer accelerates (“quantum invariance condition”).
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# Quantum Tunneling and Quantum-Classical Transitions in Large Spin Systems
## Abstract
We have studied quantum tunneling of large spins in a biaxial spin system and in a single-axial spin system with a transverse magnetic field. The asymptotically exact eigenvalues and eigenstates of the spin systems are obtained by solving the Mathieu equation. When the height $`4q`$ of energy barrier formed by the anisotropy and the magnetic field exceeds some critical values $`4q_{1T}(r)`$, the ground states and the excited states split due to tunneling of large spins. We also have presented the phase diagram of these spin systems in the extremely weak spin-phonon coupling limit. The crossover from thermal to quantum regime is second order. With further decreasing temperature, the first order phase transition occurs from a thermally-assisted resonant tunneling to thermal regime as $`q<q_{1T}(s)`$. Our theory agrees with recent experimental observations well.
In recent years quantum tunneling of magnetization (QTM) in small magnetic particles has been widely investigated in both theory and experiment because of its fundamental interest in exploring the transition between classical and quantum physics . It gives evidence of quantum mechanical behavior in macroscopic systems. Such phenomena have been observed experimentally in ferritin , the molecular nanomagnets Mn<sub>12</sub>-ac , Fe<sub>8</sub> , BaFeCoTiO , etc at very low temperatures. Theoretical calculations of such effects were performed for the ferromagnets and antiferromagnets without and with the excess spin by using path integral and WKB methods, etc.. However, these approaches depend on the semiclassical treatment and did not present the energy spectrum of the spin systems, which is necessary to give the complete solutions of spin tunneling problems. Especially, the quantization of spin levels becomes very important to explain well the experiments of quantum tunneling at very low temperatures . In very recent Refs. \[10-12\], the spin systems arouse new interest again because they provide examples exhibitting first order phase transitions, which were not known previously. The energy spectrum of the spin systems can also help us to understand the origin of the first order phase transition. In this paper, we investigate a biaxial spin system without the magnetic field and a single-axial spin system with a transverse magnetic field which are generic for spin tunneling problems studied by different methods \[1, 8, 10-12\] in the framework of the Schr$`\stackrel{..}{o}`$dinger’s picture of quantum mechanics, so that the energy spectrum as well as the physical properties of the spin systems are obtained.
We first consider QTM in the biaxial spin system without the magnetic field described by the Hamiltonian
$$H=AS_z^2BS_x^2,$$
$`(1)`$
where the anisotropy constants $`A`$ and $`B`$ are taken to be positive. So the ground state of the system corresponds to spin $`𝐒`$ pointing the positive or negative x axis. To diagonalize the Hamiltonian (1), we first write its matrix representation in the basis $`|s,m>`$, which has the properties: $`𝐒^2|s,m>=s(s+1)|s,m>,S_z|s,m>=m|s,m>`$ and $`S_\pm |s,m>(S_x\pm iS_y)|s,m>=\sqrt{s(s+1)m(m\pm 1)}|s,m\pm 1>`$, where $`s`$ is the spin quantum number to be taken as an integer , $`m=s,s+1,\mathrm{},s`$ and the unit of $`\mathrm{}=1`$ is used. Let $`E`$ and $`\mathrm{\Psi }_m`$ be the eigenenergies and the eigenstates of $`H`$, respectively, then we have
$$\begin{array}{c}\sqrt{[s(s+1)(m1)^2]^2(m1)^2}\mathrm{\Psi }_{m2}\\ +\sqrt{[s(s+1)(m+1)^2]^2(m+1)^2}\mathrm{\Psi }_{m+2}\\ +4[\frac{E}{B}+\frac{1}{2}s(s+1)(\lambda +\frac{1}{2})m^2]\mathrm{\Psi }_m=0,\end{array}$$
$`(2)`$
where $`\lambda =\frac{A}{B}`$. For a large spin system, $`s`$ is a sufficiently large number, i.e. $`s1`$. Define $`\mathrm{\Phi }_m=(1)^{[\frac{m}{2}]}\mathrm{\Psi }_m`$, where $`[\frac{m}{2}]`$ denotes the integer part of $`\frac{m}{2}`$,then Eq. (2) becomes
$$(1x^2)\frac{d^2\mathrm{\Phi }}{dx^2}2x\frac{d\mathrm{\Phi }}{dx}+[\frac{E}{B}\frac{1}{4}+\lambda s(s+1)x^2\frac{1}{4}\frac{1}{1x^2}]\mathrm{\Phi }=0$$
$`(3)`$
in the large-$`s`$ limit. Here $`x=\frac{m}{\sqrt{s(s+1)}}`$ and only the leading terms are remained. Taking the transformations $`\mathrm{\Phi }=(1x^2)^{\frac{1}{4}}y(x)`$ and $`x=\mathrm{sin}t`$ and substituting them into Eq. (3), we finally obtain the well-known Mathieu equation
$$\frac{d^2y}{dt^2}+[\mathrm{\Lambda }(q)2q\mathrm{cos}(2t)]y=0,$$
$`(4)`$
which describes the motion of a particle with the mass $`\frac{1}{2}`$ in the cosine potential. Here the characteristic values $`\mathrm{\Lambda }(q)=\frac{E}{B}+2q`$ and $`q=\frac{1}{4}\lambda s(s+1)`$. Obviously, when $`q=0`$, $`\mathrm{\Lambda }_{|m|}(0)=m^2`$, then $`E_m=Bm^2`$, which are the exact eigenenergies of $`(1)`$ in the eigenstates $`|s,m>_x`$. Now we have completed the mapping of the spin problem onto a particle problem. The energy spectrum of the Hamiltonian (1) is determined completely by the Mathieu equation.
The Mathieu equation (4) has been studied in a large set of literature due to its physically basic importance. It is known that there exist periodic solutions of period $`n\pi `$, where $`n`$ is any positive integer. However, the solutions relative to quantum tunneling problem discussed here are only those even and odd ones with periods $`\pi `$ and $`2\pi `$, i.e., $`y=`$ce$`{}_{r}{}^{}(t,q)`$ and se$`{}_{r}{}^{}(t,q)`$, $`r=0,1,\mathrm{},s`$. Assume that $`a_r`$ and $`b_r`$ are the characteristic values associated with the even and odd periodic solutions, respectively, then these characteristic values form a countable sequence, i.e. $`a_0<b_1<a_1<b_2<a_2<\mathrm{}<b_s<a_s`$ for $`q>0`$ and $`a_{2i}(q)=a_{2i}(q)`$, $`b_{2i}(q)=b_{2i}(q)`$, $`a_{2i+1}(q)=b_{2i+1}(q)`$ and $`b_{2i+1}(q)=a_{2i+1}(q)`$. We note that the approximately degenerate characteristic levels $`(a_r,b_r)`$ are lifted remarkably as the energy barrier parameter $`q`$ exceeds the critical value $`q_{1T}(r)`$. The higher the characteristic level is, the larger the critical value is, i.e. $`q_{1T}(i)>q_{1T}(j)`$ for $`i>j`$. The value of $`q_{1T}(r)`$ can be estimated by letting the height of the energy barrier be equal to the energy $`\mathrm{\Lambda }_r(0)`$ of the particle in the vanishing energy barrier, i.e. $`4q_{1T}(r)=r^2`$, which is a good approximation for large $`r`$. This shows clearly that when $`q>q_{1T}(r)`$, the splitting of the characteristic levels $`(a_r,b_r)`$ is induced by tunneling of the particle. As $`q`$ is larger than $`q_{2T}(r)\frac{[s(s+1)]^2}{4(sr+1)^2}`$ , the levels $`b_r`$ and $`a_{r1}`$ degenerate approximately due to very high energy barrier. Therefore, in order to observe experimentally tunneling splitting of level $`r`$, $`q`$ must locate between $`q_{1T}(r)`$ and $`q_{2T}(r)`$. We also note that tunneling of the particle mainly occurs in those characteristic levels ranging from the highest level to one with characteristic value $`a_ss^2+0(s)`$ as $`q<q_{1T}(s)`$, which is important to evaluating approximately the critical temperature from thermal to quantum regime below.
For the Hamiltonian (1), its lower excited states correspond to the higher characteristic levels of the Mathieu equation. In order to see the tunneling splitting of the degenerate ground states, $`q`$ must exceed the critical value $`q_{1T}(s)\frac{1}{4}s^2`$, i.e. $`\lambda _{1T}(s)1\frac{1}{s}`$, which coincides with the result reported previously $`[11]`$ when $`s`$ is very large. When $`q>q_{2T}(m)`$, we obtain tunneling splitting of the ground states as well as the excited states $`[14]`$
$$\mathrm{\Delta }E_m=B(a_mb_m)=B(2\sqrt{\lambda s(s+1)}m)+0(q^{\frac{1}{2}}).$$
$`(5)`$
When $`q\mathrm{}`$, the ground state $`a_s`$ is singlet while all the excited states are twofold degenerate, which coincide with the degeneracy of the Hamiltonian (1) with $`B0`$.
Our another example is devoted to the single-axial spin system with a transverse magnetic field, which Hamiltonian has form
$$H=hS_zBS_x^2,$$
$`(6)`$
where $`B>0`$ and $`h>0`$. The Hamiltonian (6) can be also diagonalized by the same method mentioned above. Its energy spectrum is determined by the Mathieu quation (4), too. However, in this case, $`x=\mathrm{cos}(2t)`$, $`\mathrm{\Lambda }(q)=\frac{4E}{B}`$ and $`q=\frac{2h\sqrt{s(s+1)}}{B}`$. As $`h=0`$, then $`\mathrm{\Lambda }_r(0)=r^2`$ and $`E_m=Bm^2`$. So the ground states and the excited states of the Hamiltonian (6) correspond to the characteristic levels $`(a_{2s},b_{2s})`$ and $`(a_{2r},b_{2r})(r<s)`$, respectively. According to that $`4q_{1T}(m)=(2m)^2`$, we have the critical magnetic field $`h_{1T}(m)=\frac{Bm^2}{2\sqrt{s(s+1)}},m1`$. As $`q<h_{1T}(m)`$, there is no tunneling between the approximately degenerate energy levels $`E_m`$ (i. e. $`a_{2m}`$ and $`b_{2m}`$) of the Hamiltonian (6). As $`h>h_{1T}(s)=\frac{Bs^2}{2\sqrt{s(s+1)}}`$, the ground states occur splitting, which is the same with that in Ref. for very large $`s`$. We also note that $`h`$ must be smaller than $`h_c=2Bs`$, at which the degenerate ground states coincide. As $`q1`$, we have
$$\mathrm{\Delta }E_m=\frac{1}{4}B(a_{2m}b_{2m})=B\left(\sqrt{\frac{2h\sqrt{s(s+1)}}{B}}\frac{1}{2}m\right)+0(q^{\frac{1}{2}}).$$
$`(7)`$
When $`q\mathrm{}`$, all the energy levels $`a_{2m}`$ and $`b_{2m}`$ are singlet, which agree with the energy spectrum of the Hamiltonian (6) with $`B0`$. Because there is no Kramers’ degeneracy in the Hamiltonian (6), tunneling splitting also holds for half odd integer spin $`s`$ .
Up to now, we have constructed the energy spectrum of the pure spin systems (1) and (6). However, the experiments of observing quantum tunneling were performed at low temperatures. So the influence of phonons or thermal activation on tunneling must be considered. Here, we assume the spin-phonon coupling is so weak that the energy spectrum of the spin systems is not changed too much. The transition between the ground state and the excited states can be completed by absorbing or emitting phonons. With the help of the Debye theory of a solid , we obtain approximately the critical temperature $`T_c`$ satifying
$$\begin{array}{ccc}\hfill 3k_BT_cD(\frac{T_c}{\theta })& & Bs^2\hfill \end{array}$$
$`(8)`$
as $`q<q_{1T}(s)`$ and
$$\begin{array}{cccc}\hfill 3k_BT_cD(\frac{T_c}{\theta })& & B(a_sb_r)\hfill & \text{for (1)}\hfill \\ & & \frac{1}{4}B(a_{2s}b_{2r})\hfill & \text{for (6)}\hfill \end{array}$$
$`(9)`$
as $`q_{2T}(r)qq_{2T}(r+1)`$ and $`q_{2T}(2r)qq_{2T}(2r+1)`$, respectively. Here the Debye function $`D(x)\frac{1}{5}\pi ^4x^3`$ at low temperature, $`k_B`$ and $`\theta `$ are the Boltzmann constant and the Debye temperature, respectively. Obviously, the crossover from the high temperature (thermal) phase to the low temperature (quantum) phase is second order because tunneling splitting of large spin continuously transforms into zero. As $`q<q_{1T}(s)`$, in order to see tunneling, large spin must absorb at least energy $`B(s^24q)`$ and $`B(s^2q)`$ for the models (1) and (6), respectively, provided mainly by phonons or thermal activation. So we have the first order phase transition temperature $`T_0`$,
$$\begin{array}{cccc}\hfill 3k_BT_0D(\frac{T_0}{\theta })& & B(s^24q)\hfill & \text{for (1).}\hfill \\ & & B(s^2q)\hfill & \text{for (6).}\hfill \end{array}$$
$`(10)`$
The phase diagram of the spin systems $`(1)`$ and $`(6)`$ is shown in Fig. $`1`$.
Here we compare our theoretical results with recent experiments of quantum tunneling. The Hamiltonian $`(1)`$ is believed to be a good description for the molecular nanomagnet Fe<sub>8</sub> with spin $`s=10`$ . The parameters $`A=0.092`$K and $`B=0.224`$K, then $`q=11.295<q_{1T}(s)=25`$. Therefore, with decreasing temperature, the spin system enters the thermally-assisted resonant tunneling regime II from thermal regime I, as observed down to $`0.067`$K. However, tunneling can not be seen possibly if the temperature is further lowered to below $`T_0`$, which would be of interest to be verfied experimentally. The Hamiltonian (6) has been found to be a good approximation for the molecular nanomagnet Mn<sub>12</sub>-ac with spin $`s=10`$ . For this nanomagnet, $`B=0.68`$K, then $`h_{1T}(s)=3.24`$K. The parameter $`h=g_{}\mu _BH`$ was changed from 6.38K to 8.93K in experiment, i.e. the magnetic field $`H`$ was applied between 5T and 7T. So the spin system fell into quantum regime III, where the two states of the fundamental doublet of the high-spin molecule are lifted. The transitions between the two states induced by a two-phonon process have been observed from 0.1K to 0.02K when $`H=6.06`$T . Obviously, as the spin systems just set in the thermal regime IV or $`q>q_{2T}(s)`$, tunneling does not be observed in experiments even at very low temperatures .
In summary, we have investigated quantum tunneling of large spins based on simple models (1) and (6). The asymptotically exact spectrum of the spin systems is completely determined in the full range of the magnetic anisotropy and the magnetic field. The degeneracy of all the energy levels is removed due to QTM. The phase diagram presented here coincides with the experimental observations. To our knowledge, the thermal regime IV in Fig. 2 was not predicted by previous theories. We think that the first order phase transition is not limited to the spin systems (1) and (6) and also exists generally in other spin systems, depending on whether the energy barrier and the ground states as well as the low-lying excited states cross. The method of solving the spin Hamiltonians used above is also applied to the other large spin systems, such as the other symmetric ferromagnets, antiferromagnets and etc., which is in progress.
This work is supported in part by grants from the Hong Kong Research Grants Council (RGC), the Hong Kong Baptist University Faculty Research Grant (FRG), the Sichuan Youth Science and Technology Foundation, the NSF of the Sichuan Educational Commission and the NSF of China (No. 19677101).
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# The evolution of the stellar populations in low surface brightness galaxies
## 1 Introduction
Deep searches for field galaxies in the local universe have revealed the existence of a large number of galaxies with low surface brightnesses that are hard to detect against the night sky (see the reviews by Impey & Bothun 1997 and Bothun, Impey & McGaugh 1997). Many different kinds of LSB galaxies have so far been found, including giant LSB galaxies (e.g. Sprayberry et al. 1995) and red LSB galaxies (O’Neil et al. 1997), but the most common type seems to be “blue LSB galaxies”: late-type, disk-dominated spirals with central surface brightnesses $`\mu _0(B)\begin{array}{c}>\\ \end{array}23`$ mag arcsec<sup>-2</sup>.
In this paper, we will concentrate on these late-type LSB spirals. Observations show that they are neither exclusively dwarf systems nor just the faded counterparts of “normal” high surface brightness (HSB) spirals (de Blok et al. 1996, hereafter dB96). In many cases, LSB galaxies follow the trends in galaxy properties found along the Hubble sequence towards very late types. These trends include increasingly blue colours (e.g. Rönnback 1993; McGaugh & Bothun 1994; de Blok et al. 1995 \[hereafter dB95\]; Bell et al. 1999), decreasing oxygen abundances in the gas (e.g. McGaugh 1994; Rönnback & Bergvall 1995; de Blok & van der Hulst 1998a), and decreasing Hi surface densities (from type Sc onwards; e.g. van der Hulst et al. 1993, hereafter vdH93). Despite the low gas densities, LSB spirals rank among the most gas-rich disk galaxies at a given total luminosity as their Hi disks in general are extended (Zwaan et al. 1995; dB96). The fact that LSB galaxies still have large reservoirs of gas together with their low abundances suggest that the amount of star formation in the past cannot have been very large (vdH93; Van Zee et al. 1997). Clearly, LSB galaxies are not the faded remnants of HSB spirals.
The unevolved nature of LSB spirals can be interpreted in various ways. For instance, LSB spirals could be systems whose stellar population is young and in which the main phase of star formation is still to occur. Alternatively, the stellar population in LSB spirals could be similar to those in HSB galaxies but with a young population dominating the luminosity. This has been explored and modelled in studies by e.g. Knezek (1993), Jimenez et al. (1998), dB96, Gerritsen & de Blok (1999) and more recently in two papers by Bell et al. (1999) and Bell & de Jong (1999). These studies all indicate a scenario in which LSB galaxies are unevolved systems with low surface densities, low metallicities and a constant or increasing star formation rate (cf. Padoan et al. 1997, however, for a different opinion).
In this paper we address these various scenarios for the evolution of LSB galaxies by detailed modelling of their spectro-photometric and chemical properties. Model results are compared directly with the observed colours, gas phase abundances, gas contents, and current star formation rates of LSB galaxies. Using simple model star formarion histories we confirm the results of previous studies and show that the ratio of young stars to old stars in LSB galaxies is larger than typically found in HSB galaxies.
In reality, however, star formation is unlikely to proceed as smoothly as the various population synthesis models suggest. Gerritsen & de Blok (1999) point out that small surges in the star formation rate are likely to be very important in determining the observed colours of LSB galaxies. In this paper we will pay special attention to the influence of small bursts of star formation on the optical colours, and show that they can have a significant influence on the inferred properties of LSB galaxies.
This paper is organized as follows. We briefly compare observational data of LSB galaxies with those of spirals and dwarf galaxies in Sect. 2. In Sect. 3, we describe the ingredients of the galactic evolution model we used. Some general properties of the model are presented in Sect. 4. The comparison with observational data is presented and discussed in Sect. 5. The impact of small amplitude star formation bursts in LSB galaxies is investigated in Sect. 6 and predicted star formation rates are compared with the observations in Sect. 7. A short discussion of where LSB galaxies fit in the grand scheme of galaxy formation and evolution is presented in Sect. 8. A short summary is given in Sect. 9.
## 2 Observational characteristics of LSB galaxies
### 2.1 Sample selection
We use the sample described in dB96. This consists of 24 late-type LSB galaxies (inclinations up to $`60^{}`$), taken from the lists by Schombert et al. (1992) and the UGC (Nilson 1973). The sample is representative for the LSB galaxies generally found in the field. We selected a subsample of 16 LSB galaxies for which good data are available. For these systems, optical data have been taken from dB95, Hi data from dB96, and abundance data from McGaugh & Bothun (1994) and de Blok & van der Hulst (1998a).
The galaxy identification and Johnson $`UBV`$ and Kron-Cousins $`RI`$ are given in columns (1) to (6) in Table 1 (a Hubble constant of $`H_0=100`$ km s<sup>-1</sup> Mpc<sup>-1</sup> was used). Column (7) lists the $`B`$-band luminosities; column (8) the Hi-mass-to-light ratio; column (9) and (10) the Hi and dynamical masses respectively. The gas fraction $`\mu _{\mathrm{rot}}`$ in column (11) is given by $`M_{\mathrm{gas}}`$/$`(M_{\mathrm{gas}}+M_{,\mathrm{max}})`$ where $`M_{,\mathrm{max}}`$ denotes the mass of the stellar component obtained from maximum disk fitting of the rotation curve (dB96) (see Sect. 2.6). The mean \[O/H\] abundances are listed in column (12).
### 2.2 Magnitudes and colours
We compare in Fig. 1$`a`$ and $`b`$ magnitudes and broadband colours of LSB galaxies to those of late-type HSB face-on spirals (de Jong & van der Kruit 1994) and dwarf galaxies (Melisse & Israel 1994; Gallagher & Hunter 1986, 1987). The usual distinction between HSB and LSB galaxies is purely an artificial one since normal galaxies along the Hubble sequence show a continuous range in central surface brightness, i.e. ranging from values around the Freeman value for early type systems to the very faint values observed for the late-type LSB galaxies in our sample (e.g. de Jong 1995; dB95). For the purpose of this paper we will continue to use the distinction as a shorthand, where “HSB galaxy” should be interpreted as “typical Hubble sequence Sc galaxy.” Fig. 1$`a`$ and $`b`$ show that LSB galaxies are generally bluer and fainter than their HSB counterparts (e.g. McGaugh 1992; vdH93; dB95).
### 2.3 Abundances
Estimates of the ISM abundances in LSB galaxies predominantly rely on abundance determinations of their Hii regions. Oxygen abundances of Hii regions listed in Table 1 are taken from McGaugh (1994) and de Blok & van der Hulst (1998a).
We assume that the Hii-region abundances on average are a reasonable indicator of the ISM abundances within a given LSB spiral. The intrinsic scatter in \[O/H\] among different Hii regions within a given LSB galaxy is usually of the order 0.2 dex (e.g. McGaugh 1994). Bell & de Jong (1999) find that the gas metal abundances in LSB galaxies generally trace the stellar metallicities quite well, with the stars being $`0.4`$ dex more metal-poor.
In Fig. 1c we compare mean \[O/H\] abundances of Hii regions in LSB galaxies with those in HSB spirals (Zaritsky et al. 1994) and dwarf galaxies (Melisse & Israel 1994; Gallagher & Hunter 1986, 1987). On average, LSB galaxies generally follow the correlation between the characteristic gas-phase abundance and luminosity as found for HSB spirals (Zaritsky et al. 1994), but with a large scatter. The range in \[O/H\] at a given $`B`$ magnitude is nearly one dex. This scatter is probably related to evolutionary differences among the LSB galaxies of a given $`B`$ magnitude (e.g. in the ratio of old to young stellar populations) and/or the Hii regions they contain.
### 2.4 Extinction
Observational studies by (amongst others) Bosma et al. (1992), Byun (1992), Tully & Verheijen (1997) show that late-type spirals appear transparent throughout their disks. This supports the idea that face-on extinction in LSB galaxies is relatively low, i.e. $`E_{BV}\begin{array}{c}<\\ \end{array}0.1`$ (e.g. McGaugh 1994). This is consistent with findings by Tully & Verheijen (1998), who studied a large sample of galaxies in the UMa cluster and found the LSB galaxies in their sample to be transparent, whereas non-negligible extinction was present in the HSB galaxies.
In this study we will ignore the possible effects of dust on the colours and magnitudes. See Sect. 5.3 for a discussion.
### 2.5 Gas masses
We compare in Fig. 1d the observed atomic gas masses $`M_\mathrm{g}1.4M_{\mathrm{HI}}`$ (corrected for helium) of LSB galaxies with those of HSB galaxies and dwarfs. Strictly speaking we should also take the molecular gas mass into account. However, so far no CO-emission has been detected in late-type LSB galaxies (Schombert et al. 1990; de Blok & van der Hulst 1998b, but see Knezek 1993 for a discussion on giant LSB galaxies). The amount of molecular gas in LSB galaxies is probably small, even though relatively high CO/H<sub>2</sub> conversion factors may apply in these low metallicity galaxies (Wilson 1995, Mihos et al. 1999). We conclude that, on average, LSB galaxies are considerably more gas-rich (up to a factor $`3`$) than HSB galaxies of the same $`B`$ magnitude (see also Fig. 1g).
### 2.6 Gas fractions and mass-to-light ratios
We show in Figs. 1g and h the distribution of the mass-to-light ratio $`M_{\mathrm{HI}}/L_B`$ vs $`(BV)`$ and vs $`\mu _{rot}`$, respectively. LSB galaxies exhibit considerably higher $`M_{\mathrm{HI}}/L_B`$ ratios than HSB spirals (cf. Table 1 and McGaugh & de Blok 1997).
We would however like to know the ratio of gas mass to stellar mass, i.e. the gas fraction $`\mu _1M_{\mathrm{gas}}`$ / $`(M_{\mathrm{gas}}+M_{\mathrm{stars}}`$ of a galaxy. This involves the conversion of the observed galaxy luminosity to stellar mass. However, the mass-to-light ratio of the underlying stellar population is generally not well known and, in fact, is an important quantity to determine. Assuming a fixed value of the mass-to-light ratio will however introduce artificial trends since this ratio is expected to vary among galaxies having different star formation histories (e.g. Larson & Tinsley 1978).
As a partial solution we have determined the gas fraction $`\mu _{\mathrm{rot}}M_{\mathrm{gas}}/(M_{\mathrm{gas}}+M_{,\mathrm{max}})`$ where $`M_{,\mathrm{max}}`$ denotes the mass of the stellar component obtained from maximum disk fitting of the rotation curve (dB96). This method likely overestimates the contribution of the luminous stellar disk to the observed total mass distribution by about a factor of 2 (e.g. Kuijken & Gilmore 1989; Bottema 1995) and, therefore, also provides a lower limit to the actual gas fraction. In summary, while $`\mu _1`$ is a model quantity, $`\mu _{\mathrm{rot}}`$ is determined from observations.
Figs. 1e and 1f show the distribution of the present-day gas fraction $`\mu _{\mathrm{rot}}`$ and total gas mass vs $`(BV)`$. It can be seen that LSB galaxies and dwarfs exhibit larger gas fractions than HSB spirals. Taking into account the factor of $`2`$ overestimate, values for LSB galaxies are typically $`\mu _1\begin{array}{c}>\\ \end{array}0.5`$ and $`\mu _10.1`$ for HSB spirals.
We conclude that the available evidence supports of the view that LSB galaxies are relatively unevolved systems compared to HSB spirals. This agrees well with the fact that LSB galaxies usually display properties intermediate to those of HSB spirals and dwarf galaxies.
## 3 Model description and assumptions
We here summarise the galactic evolution model used (for a more extensive description see van den Hoek 1997, hereafter vdH97). We concentrate on the stellar contribution to the total galaxy luminosity in a given passband (other contributions are neglected). For a given star formation history (SFH), we compute the chemical enrichment of a model galaxy by successive generations of evolving stars. To derive the stellar luminosity in a given passband at a given age, we use a metallicity dependent set of theoretical stellar isochrones as well as a library of spectro-photometric data.
### 3.1 Chemical evolution model
We start from a model galaxy initially devoid of stars. We follow the chemical enrichment during its evolution assuming stars to be formed according to a given star formation rate (SFR) and initial mass function (e.g. a power law IMF: d$`N`$/d$`mM(m)m^\gamma `$). Both stellar and interstellar abundances as a function of galactic evolution time $`t_{ev}`$ are computed assuming that the stellar ejecta are returned and homogeneously mixed to the ISM at the end of the individual stellar lifetimes (i.e. relaxing the instantaneous recycling approximation; see Searle & Sargent 1972). A description of the set of galactic chemical evolution equations used can be found in e.g. Tinsley (1980) and Twarog (1980), see also vdH97. This model is a closed-box model.
We follow the stellar enrichment of the star forming galaxy in terms of the characteristic element contributions of Asymptotic Giant Branch (AGB) stars, SNII and SNIa. This treatment is justified by the specific abundance patterns observed within the ejecta of these stellar groups (see e.g. Trimble 1991; Russell & Dopita 1992). A detailed description of the metallicity dependent stellar lifetimes, element yields, and remnant masses is given by vdH97 and van den Hoek & Groenewegen (1997). We compute the abundances of H, He, O, and Fe, as well as the heavy element integrated metal-abundance $`Z`$ (for elements more massive than helium), during the evolution of the model galaxy. Both the SFH, IMF and resulting element abundances as a function of galactic evolution time are used as input for the spectro-photometric evolution model described below.
Boundary conditions to the chemical evolution model are the galaxy total mass $`M_{\mathrm{tot}}`$, its evolution time $`t_{\mathrm{ev}}`$, and the initial gas abundances. Note that solutions of the galactic chemical evolution equations are independent of the ratio of the SFR normalisation and $`M_{\mathrm{tot}}`$. Primordial helium and hydrogen abundances are adopted as $`Y_\mathrm{p}=0.232`$ and $`X=0.768`$ (cf. Pagel & Kazlauskas 1992). Initial abundances for elements heavier than helium are initially set to zero.
We list the main input parameters in Table 2, i.e. the adopted IMF-slope, minimum and maximum stellar mass limits at birth as well as the progenitor mass ranges for stars ending their lives as AGB star, SNIa, and SNII, respectively. For simplicity, we assume the stellar yields of SNIb,c to be similar to those of SNII. Furthermore, we assume a fraction $`\nu ^{\mathrm{SNIa}}=0.015`$ of all white dwarf progenitors with initial masses between $`2.5`$ and 8 $`M_{}`$ to end as SNIa. These and other particular choices for the enrichment by massive stars are based on similar models recently applied to the chemical evolution of the Galactic disk (e.g. Groenewegen et al. 1995; vdH97).
### 3.2 Spectro-photometric evolution model
The total luminosity of a galaxy in a wavelength interval $`\mathrm{\Delta }\lambda `$ is determined by: 1) the stellar contribution $`L_{}`$, 2) the ISM contribution $`L_{\mathrm{ism}}`$ (e.g. Hii-regions, high-energy stellar outflow phenomena, etc.), and 3) the contribution due to absorption and scattering $`L_{\mathrm{ext}}`$:
$$L_{\mathrm{gal}}^{\mathrm{\Delta }\lambda }(t)=L_{}+L_{\mathrm{ism}}L_{\mathrm{ext}}$$
(1)
where each term in general is a complex function of galactic evolution time. We concentrate on the stellar contribution and neglect the last two terms in Eq. (1). Then, the luminosity in a wavelength interval $`\mathrm{\Delta }\lambda `$ at galactic evolution time $`t=T`$ can be written as:
$`L_{\mathrm{gal}}^{\mathrm{\Delta }\lambda }(t=T)=`$ (2)
$`{\displaystyle _{t=0}^T}{\displaystyle _{m_\mathrm{l}}^{m_\mathrm{o}(Tt)}}`$ $`L_{}^{\mathrm{\Delta }\lambda }(m,Z(t),Tt)S(t)M(m)\mathrm{d}m\mathrm{d}t`$
where $`m_\mathrm{l}`$ denotes the lower stellar mass limit at birth, $`m_\mathrm{o}(t)`$ the turnoff mass for stars evolving to their remnant stage at evolution time $`t_{\mathrm{ev}}`$, and $`L_{}^{\mathrm{\Delta }\lambda }`$ the luminosity of a star with initial mass $`m`$, initial metallicity $`Z(t)`$, and age ($`Tt`$). We assume a separable SFR: $`S(m,t)=S(t)M(m)`$ where $`S(t)`$ is the star formation rate by number \[yr<sup>-1</sup>\] and $`M(m)`$ the IMF \[$`M_{}`$<sup>-1</sup>\]. By convention, we normalise the IMF as $`M(m)dm=1`$ where the integration is over the entire stellar mass range \[$`m_\mathrm{l}`$, $`m_\mathrm{u}`$\] at birth (cf. Table 2).
Starting from the chemical evolution model described above, we compute the star formation history $`S(m,t)`$, gas-to-total mass-ratio $`\mu (t)`$, and age-metallicity relations (AMR) $`Z_i(t)`$ for different elements $`i`$. Thus, at each galactic evolution time $`t_{\mathrm{ev}}`$ the ages and metallicities of previously formed stellar generations are known. To derive the stellar passband luminosity $`L_{}^{\mathrm{\Delta }\lambda }`$ we use a set of theoretical stellar isochrones, as well as a library of spectro-photometric data. Stellar evolution tracks provide the stellar bolometric luminosity $`L_{}^{\mathrm{bol}}`$, effective temperature $`T_{\mathrm{eff}}`$, and gravity $`g`$, as a function of stellar age for stars with initial mass $`m`$ born with metallicity $`Z_{}`$. We compute Eq. (2) using a spectro-photometric library containing the stellar passband luminosities $`L_{}^{\mathrm{\Delta }\lambda }`$ tabulated as a function of $`T_{\mathrm{eff}}`$, $`\mathrm{log}g`$, and $`Z_{}`$ (see below).
The turnoff mass $`m_\mathrm{o}(Tt)`$ occuring in Eq. (2) depends on the metallicity $`Z(t)`$ of stars formed at galactic evolution time $`t_{\mathrm{ev}}`$. For instance, the turnoff mass for stars born with metallicity $`Z=10^3`$ at a galactic age of $`t_{\mathrm{ev}}=14`$ Gyr is $`m_\mathrm{o}0.8`$ $`M_{}`$ (e.g. Schaller et al. 1992). This value differs considerably from $`m_\mathrm{o}0.95`$ $`M_{}`$ for stars born with metallicity $`Z=Z_{}`$. Such differences in $`m_\mathrm{o}`$ affect the detailed spectro-photometric evolution of a galaxy by constraining the mass-range of stars in a given evolutionary phase (e.g. horizontal branch) at a given galactic evolution time. In the models described below, we explicitly take into account the dependence of $`m_\mathrm{o}(t)`$ on the initial stellar metallicity $`Z_{}`$ (see vdH97)
### 3.3 Stellar evolution tracks and spectro-photometric data
We use the theoretical stellar evolution tracks from the Geneva group (e.g. Schaller et al. 1992; Schaerer et al. 1993). These uniform grids cover large ranges in initial stellar mass and metallicity, i.e. $`m=0.05120`$ $`M_{}`$ and $`Z=0.040.001`$, respectively. These tracks imply a revised solar metallicity of $`Z_{}=0.0188`$ with $`Y_{}=0.299`$, and $`\mathrm{\Delta }Y/\mathrm{\Delta }Z=3.0`$ for a primordial He-abundance of $`Y_\mathrm{p}=0.232`$. For stars with $`m>7`$ $`M_{}`$, these tracks were computed until the end of central C-burning, for stars with $`m=25`$ $`M_{}`$ up to the early-AGB, and for $`m<1.7`$ $`M_{}`$ up to the He-flash. For stars with $`m\begin{array}{c}<\\ \end{array}0.8`$ $`M_{}`$, we used the stellar isochrone program from the Geneva group (Maeder & Meynet, priv. comm.; see also vdH97 ).
To cover the latest stellar evolutionary phases (i.e. horizontal branch (HB), early-AGB, and AGB) for stars with $`m\begin{array}{c}<\\ \end{array}8`$ $`M_{}`$, we extended the tracks from Schaller et al. with those from Lattanzio (1991; HB and early-AGB) for $`m12`$ $`M_{}`$, and from Groenewegen & de Jong (1993; early-AGB and AGB) for $`m18`$ $`M_{}`$. These tracks roughly cover the same metallicity range as the tracks from the Geneva group. Corresponding isochrones were computed on a logarithmic grid of stellar ages, covering galactic evolution times up to $`t_{\mathrm{ev}}14`$ Gyr. Isochrones are linearly interpolated in $`m`$, log $`Z`$, and log $`t`$.
The spectro-photometric data library that we use is based on the Revised Yale Isochrones and has been described extensively by Green et al. (1987). These data include stellar $`UBVRI`$ Johnson-Cousins magnitudes covering the following ranges in $`T_{\mathrm{eff}}[\mathrm{K}]=2800`$ to 20000, log $`g`$ \[cm s<sup>-2</sup>\] $`=0.5`$ to 6, and log ($`Z/Z_{})2.5`$ to $`+0.5`$. Corresponding spectro-photometric data for stars with $`T_{\mathrm{eff}}>20000`$ K have been adopted from Kurucz (1979) at solar metallicity, covering $`T_{\mathrm{eff}}=2000050000`$ K.
Although a detailed description of the tuning and calibration of the adopted photometric model is beyond the scope of this paper, we note that the model has been checked against various observations including integrated colours and magnitudes, luminosity functions, and colour-magnitude diagrams of Galactic and Magellanic Cloud open and globular clusters covering a wide range in age and metallicity (see vdH97).
## 4 General properties of the model
Before discussing the modelling of LSB galaxies, we first describe the general behaviour of the models assuming a few different star formation scenarios.
We start from a model galaxy with initial mass $`M_\mathrm{g}(t=0)=10^{10}`$ $`M_{}`$, initially metal-free and devoid of stars. The chemical and photometric evolution of this galaxy is followed during an evolution time $`t_{\mathrm{ev}}=14`$ Gyr, assuming one of the theoretical star formation histories discussed below. Unless stated otherwise, we assume that stars are formed according to a Salpeter (1955) IMF (i.e. $`\gamma =2.35`$) with stellar mass limits at birth between 0.1 and 60 $`M_{}`$ (cf. Table 2).
The following star formation histories are considered: 1) constant, 2) exponentially decreasing, 3) linearly decreasing, 4) exponentially increasing, and 5) linearly increasing. Normalised SFRs and resulting age-ISM metallicity relations are shown in Fig. 2. For each model, the amplitude of the SFR is chosen such that a present-day gas-to-total mass-ratio $`\mu _1=0.1`$ is achieved.
In columns (2) to (6) of Table 3, we list the functional form of the SFR, average past and current SFRs, the ratio of current and average past SFRs $`\alpha `$, and the present-day mass averaged ISM oxygen abundance. The average past SFR is roughly the same for all models ending at $`\mu _1=0.1`$ and is $``$SFR$`=0.9`$ $`M_{}`$yr<sup>-1</sup>. In contrast, present-day SFRs range from SFR<sub>1</sub> = 0.07 to 3 $`M_{}`$yr<sup>-1</sup> and in fact determine the contribution by young stars to the integrated light of the model galaxy. Current oxygen abundances predicted are \[O/H\]$`{}_{1}{}^{}+0.25`$ and are mainly determined by $`\mu _1`$, the IMF, and the assumed mass limits for SNII (cf. Table 3). Column (7) gives the number of stars formed, Column (8) their total $`B`$ band luminosity. Column (10) shows the resulting Hi mass-to-light ratio. Column (11) gives the slope of the IMF used, while Column (11) refers to additional Notes.
Let us consider the photometric evolution of the constant and exponentially decreasing SFR models in some more detail. Fig. 3 shows the evolution of the total number of MS and post-MS stars for exponentially decreasing SFR model A1. The current total number of MS stars is roughly $`410^{10}`$, compared to $`10^8`$ post-MS stars (all phases) and $`310^4`$ for AGB stars only. The present-day mean luminosity of stars on the MS is $`10^1L_{}`$ compared to $`10^4L_{}`$ for AGB stars. The current bolometric galaxy luminosity is determined mainly by MS stars ($`L_{\mathrm{MS}}410^9L_{}`$). In particular, AGB stars ($`L_{\mathrm{AGB}}10^8L_{}`$) are relatively unimportant.
We show in Fig. 4 the $`B`$ and $`I`$-band magnitudes of stars in different evolutionary phases (again for the exponential model). As for the bolometric galaxy luminosity, MS stars generally dominate in the $`B`$-band. However, in the $`I`$-band, RGB and HB stars are nearly as important as MS stars, at least in later stages of galactic evolution. Due to the cooling of old, low-mass MS stars as well as the increasing contribution by RGB and HB stars with age, the current total $`I`$-band magnitude is considerably brighter than that in the B-band. This qualitative model behaviour is insensitive to the adopted star formation history (for e-folding times larger than $`35`$ Gyr) but instead is determined by the assumed IMF and the stellar input data used. Thus, constant star formation models exhibit a similar behaviour apart from being brighter by about one magnitude in all passbands at late times (cf. Fig. 4).
Fig. 5 illustrates the sensitivity of broadband colours to the star formation history for constant and exponentially decreasing SFR models. The colours considered here redden with galactic age. In general, differences between colours such as for different SFR models are less than the variations of these colours with age for a given model. Assuming a galactic age of e.g. 8 instead of 14 Gyr has limited effect on the resulting galaxy colours (e.g. less than 0.1 mag in $`BV`$), even though absolute magnitudes are substantially altered. Both age and extinction effects can result in substantial reddening of the colours of a stellar population in almost the same manner and it is difficult to disentangle their effects on the basis of photometry data alone. Clearly, galaxy colours alone are not well suited to discriminate between different SFR models, even when internal extinction is low and other reddening effects are negligible (see below).
## 5 Results
As our emphasis will be on a comparison between the evolution of HSB and LSB galaxies, we will first focus on the relatively well-known evolution of HSB galaxies, using a model chosen to emphasise the observed properties of HSB galaxies. Using this as a starting point, we explore other models to find the best model for LSB galaxies.
### 5.1 HSB galaxies: $`\mu _1=0.1`$ models
Models using exponentially decreasing SFRs and different e-folding times for different Hubble types (e.g. Larson & Tinsley 1978; Guiderdoni & Rocca-Volmerange 1987; Kennicutt 1989; Bruzual & Charlot 1993; Fritze-v. Alvensleben & Gerhard 1994) are usually found to be in good agreement with observations of HSB galaxies. Therefore we will mainly focus on the exponential SFR models discussed above. As the gas-richness is one of the obvious differences between HSB and LSB galaxies, we will first focus on models with $`\mu _1=0.1`$ to describe HSB galaxies.
Fig. 6 shows the results of the exponentially decreasing SFR model: SFR $`\mathrm{exp}(t/\tau _{\mathrm{sfr}})`$ with $`\tau _{\mathrm{sfr}}=5`$ Gyr ending at $`\mu _1=0.1`$ at $`t=t_{ev}=14`$ Gyr for several initial galaxy masses $`M_\mathrm{g}(t=0)10^810^{10}`$ $`M_{}`$. The precise value of $`\mu _1`$ was achieved by scaling the amplitude of the SFR accordingly. Galaxy colours and abundances are not affected by this scaling while absolute magnitudes and final gas masses scale with $`M_\mathrm{g}(t=0)`$ (as indicated in Fig. 6).
Present-day ($`BV`$) and ($`RI`$) colours for this model are 0.6 and 0.5, respectively. HSB galaxies with $`(BV)\begin{array}{c}>\\ \end{array}0.6`$ and $`(RI)\begin{array}{c}>\\ \end{array}0.5`$ can be explained only when substantial amounts of internal dust extinction are incorporated. The reason is that even though a single stellar population born with metallicity $`Z=0.02`$ may become as red as $`(BV)1.1`$ and $`(RI)0.8`$, the luminosity contribution by young stellar populations (with ages less than a few Gyr) results in substantial blueing of the galaxy colours. For the same reason, values of $`(BV)\begin{array}{c}>\\ \end{array}0.65`$ are not predicted by dust-free models independent of the adopted (exponential) SFR (as long as $`\tau \begin{array}{c}>\\ \end{array}4`$ Gyr) or IMF. Therefore, considerable reddening of the younger stellar populations is required to explain the colours of HSB galaxies in this model. Figs. 6c and d show that both the \[O/H\] abundance ratios and present-day gas masses of HSB galaxies can be explained reasonably well for $`M_\mathrm{g}(t=0)>10^{10}`$ $`M_{}`$. These conclusions agree with results from many previous models (with different stellar input data) for HSB galaxies (e.g. Larson & Tinsley 1978; Guiderdoni & Rocca-Volmerange 1987).
It is clear from Fig. 6 that this model is not able to explain the properties of the LSB galaxies in our sample. The model colours are too red while the predicted \[O/H\] abundances are too high. As LSB galaxies have higher gas-fractions we thus need to broaden the range of gas-fractions.
### 5.2 LSB galaxies: $`\mu _10.1`$ models
We now consider models with gas fractions $`\mu _1=`$ 0.025, 0.1, 0.3, 0.5, 0.7, and 0.9. First we will examine the constraints which individual properties such as abundance and gas-fraction impose on the models, before combining these in a final model.
#### 5.2.1 Colours and magnitudes
Results for a model with initial mass $`M_\mathrm{g}(t=0)=10^{10}`$ $`M_{}`$ are shown in Fig. 7. While the age distribution of the stellar populations is identical in all the models shown, present-day ($`BV`$) and ($`RI`$) colours are found to decrease by 0.2 and 0.1 mag, respectively, when going from models ending at $`\mu _1=0.025`$ to $`0.9`$. This bluing effect is due to the decrease of stellar metallicities for models ending at increasingly higher gas fractions.
The LSB galaxies can be distinguished in two major groups by means of their colours. First, LSB galaxies with $`(BV)\begin{array}{c}>\\ \end{array}0.5`$ which can best be modelled with exponentially decaying SFR models ending at values $`\mu _1=0.50.7`$ at age 14 Gyr. These results can be shifted towards brighter or fainter $`B`$ magnitudes by changing the initial gas masses accordingly. This leaves the resulting colours unaltered (cf. Fig. 6).
Second, relatively blue LSB galaxies with $`(BV)\begin{array}{c}<\\ \end{array}0.4`$ cannot be fitted by exponentially decreasing SFR models with $`\tau \begin{array}{c}<\\ \end{array}5`$ alone (assuming $`t_{\mathrm{ev}}=14`$ Gyr), regardless of their current gas fraction $`\mu _1`$. There are several possible explanations: a relatively young “surplus” stellar population may influence the galaxy colours resulting from an underlying exponentially decreasing SFR. Alternatively, these LSB galaxies may have started forming stars recently (e.g. $`58`$ Gyr instead of 14 Gyr, cf. Fig. 7a), or they may have more slowly decaying or constant SFRs which generally result in present-day colours of $`(BV)\begin{array}{c}<\\ \end{array}0.4`$ and $`(RI)\begin{array}{c}<\\ \end{array}0.35`$ (but the latter option appears to be excluded because of their abundances and gas fractions; see below).
#### 5.2.2 Abundances
The observed range in \[O/H\] abundances for LSB galaxies is well explained by exponentially decaying SFR models with $`\mu _1\begin{array}{c}>\\ \end{array}0.3`$ (cf. Fig. 7c). Constant SFR models ending at $`\mu _1\begin{array}{c}>\\ \end{array}0.3`$ are also consistent as the abundances of elements predominantly produced in massive stars are in general determined by the present-day gas fraction $`\mu _1`$ and are insensitive to the detailed underlying star formation history (e.g. Tinsley 1980). However, these latter models appear to be excluded by the measured gas masses (see 5.2.4 and 5.2.5) (for a closed box model). Though the abundances of metal-poor LSB galaxies with \[O/H\] $`\begin{array}{c}<\\ \end{array}1`$ can be well explained by models ending at $`\mu _1=0.9`$, it is likely that the exponential models presented here are an over-simplification for these gas-dominated systems. It is more likely that they have experienced a low and sporadic star formation rate.
#### 5.2.3 Gas content
Present-day gas masses observed in LSB galaxies can be well fitted by exponentially decaying SFR models ending at $`\mu _10.5`$ (Fig. 7d). This is consistent with the range of $`\mu _1\begin{array}{c}>\\ \end{array}0.3`$ derived from the \[O/H\] data. Exponentially decreasing SFR models ending at $`\mu _1\begin{array}{c}>\\ \end{array}0.6`$ are inconsistent with the observations unless we assume that LSB galaxies have started forming stars only a few Gyr ago.
Fig. 8 displays the gas mass vs $`(BV)`$ for constant and exponentially decaying SFR models. Exponentially decreasing SFR models are able to explain simultaneously values of $`\mu _10.5\pm 0.2`$, $`(BV)\begin{array}{c}>\\ \end{array}0.5`$, and $`M_\mathrm{g}\begin{array}{c}>\\ \end{array}10^9`$ $`M_{}`$, as observed for the majority of the LSB galaxies in our sample. Over the entire range of possible gas fractions, constant (or increasing) SFR models are only able to explain the bluest galaxies.
Fig. 8 demonstrates that exponentially decreasing SFR models ending at $`\mu _10.5`$ are in good agreement with the observed hydrogen mass-to-light ratios $`M_{\mathrm{HI}}/L_B`$, in contrast to constant (or slowly decreasing) SFR models (expect for the bluest systems).
The conclusion that constant and slowly decreasing SFR models are inconsistent with the observed photometric and chemical properties of LSB galaxies is based on the assumption of negligible amounts of dust extinction in these systems. However, if a considerable fraction of the LSB spirals in our sample would suffer from extinctions of $`E_{(BV)}=0.10.25`$, the $`M_{\mathrm{HI}}/L_B`$ ratios predicted by constant SFR models would increase by a factor $`1.42.5`$ after correction for internal extinction. In this manner, constant and slowly decreasing SFR models ending at gas fractions $`\mu _1=0.30.7`$ after $`14`$ Gyr could also explain relatively red LSB galaxies with $`(BV)\begin{array}{c}>\\ \end{array}0.5`$ and $`M_{\mathrm{HI}}/L_B`$ ratios as large as $`1.5`$. Even though we argued in Sect. 2 that internal extinction in LSB galaxies is unlikely to exceed $`E_{(BV)}0.1`$, knowing the dust content of LSB galaxies is obviously of crucial importance in modelling their evolution.
#### 5.2.4 Summary: the best model
To summarise, we find that exponentially decreasing SFR models ending at $`\mu _1=0.30.7`$ are in agreement with the colours, magnitudes, \[O/H\] abundances, gas contents, and mass-to-light ratios observed for LSB galaxies with $`(BV)\begin{array}{c}>\\ \end{array}0.45`$. Blue LSB galaxies with $`(BV)\begin{array}{c}<\\ \end{array}0.45`$ cannot be fitted by exponentially decreasing SFR models without an additional light contribution from a young stellar population. Alternatively, such LSB galaxies may have experienced constant SFRs, or may be much younger than 14 Gyr.
### 5.3 HSB/LSB: Just an extinction effect?
Since the models require an internal extinction in HSB galaxies up to $`E_{(BV)}0.5`$ (corresponding to $`A_V1.5`$ mag), independent of the SFH, one could argue that extinction may be the main explanation for the difference in colour observed between LSB and HSB galaxies.
As shown in this paper the models indicate that the underlying stellar populations in LSB and HSB galaxies are still distinctly different due to differences in the chemical evolution of these systems. Therefore, extinction effects, although important, cannot be the entire explanation for the colour differences observed.
## 6 The effects of small star formation bursts
### 6.1 Luminosity contribution by HII regions in LSB galaxies
The presence of Hii regions in all the LSB galaxies in our sample suggests that recent star formation is a common phenomenon in these systems, and may thus influence the observed properties. In the previous section we found that for some LSB galaxies a simple exponential SFR was unable to explain the observed properties and that an additional “burst” of young stars was needed. Here we examine the effects of such bursts on the colours of LSB galaxies. To this end, we selected all Hii regions that could be identified by eye, either from H$`\alpha `$ or $`R`$-band CCD images, and added their total luminosity in the $`I`$-band. We restrict ourselves to the $`I`$-band data as to be least affected by extinction that may be present in or near Hii regions (McGaugh 1994). We define $`\eta _I`$ as the ratio of this Hii region integrated luminosity and the total luminosity of the corresponding LSB galaxy. As we are probably incomplete at low Hii region luminosities, this ratio gives a lower limit for the Hii region contribution to the total luminosity.
In columns (1) and (2) of Table 4, we list the LSB galaxy identification and number of Hii regions selected. The number of Hii regions identified within individual LSB galaxy ranges from a few to $`25`$. For the ensemble of Hii regions in each LSB galaxy, we tabulate the absolute $`I`$ magnitudes as well as the corresponding ratios $`\eta _I`$ of the Hii region integrated luminosity and total LSB galaxy luminosity, in columns (3) and (4). Mean ($`BV`$) colours for the Hii regions and for the LSB galaxy as a whole are given in columns (5) and (6), respectively.
For most of the LSB galaxies in our sample, the contribution of the Hii regions to the total light emitted by LSB galaxies does not exceed $`\eta _I=0.050.1`$. However, as Hii regions are likely to contain an increased amount of dust (McGaugh 1994) the actual contributions may be higher by a factor of $`2`$ in $`I`$. Thus, the values of $`\eta _I`$ provide lower limits to the actual luminosity contributions of the Hii regions. For some LSB galaxies, e.g. F568-V1 and F577-V1, the Hii region contribution is found to be as high as $`\eta _I=0.2`$. These systems contain a only modest number of Hii regions so that their Hii regions on average may be larger and/or brighter than those present in several other LSB galaxies.
Figure 10 shows the resulting Hii region contributions $`\eta _I`$ for the SFR models discussed in Sect. 4. If we assume a maximum age $`\tau _{\mathrm{HII}}=5`$ Myr for the Hii regions observed, then this implies that stars more massive than $`25`$ $`M_{}`$ are associated with these regions, according to the stellar evolution tracks from the Geneva group (see below).
Corrections for dust extinction within the Hii regions and LSB galaxy as a whole, respectively, will shift the observations in the directions as indicated in Fig. 10 (assuming a mean Galactic extinction curve). From the $`I`$ band observations we can conclude that the values of $`\eta _I`$ predicted by smoothly varying SFR models are systematically too low for most of the LSB galaxies in our sample.
This is true in particular for F568-V1, F577-V1, and U628 for which values of $`\eta _{I,R}\begin{array}{c}>\\ \end{array}0.20.25`$ suggest that star formation has been recently enhanced by factors $`510`$ relative to the SFRs predicted by exponentially decreasing or constant SFR models.
### 6.2 Effects of small amplitude star formation bursts
We investigate whether the observed values of $`\eta _I0.25`$ in the $`I`$ band, as observed for several LSB galaxies discussed above, can be explained by small-amplitude bursts of star formation.
We assume a Gaussian star formation burst profile with a given amplitude $`A_b`$ and dispersion $`\sigma `$. We follow its evolution during 1 Gyr with a time resolution of $`0.1`$ Myr at time of burst maximum and of $`2`$ Myr at roughly 5$`\sigma `$ from burst maximum. We superimpose the star formation burst on an exponential SFR model, as discussed in Sect. 5.1.
We initially assume a galactic evolution time at burst maximum of $`t_\mathrm{b}=13`$ Gyr, burst maximum amplitude $`A_\mathrm{b}=8`$ $`M_{}`$ yr<sup>-1</sup>, burst duration $`\mathrm{\Delta }t_\mathrm{b}=5`$ Myr, maximum Hii region lifetime $`\tau _{\mathrm{HII}}=5`$ Myr, and an initial galaxy mass of $`10^{10}`$ $`M_{}`$. For simplicity, we neglect any influence of the burst on the chemical evolution of the model galaxy.
We show the result in Fig. 11a. At $`t=t_b=13`$ Gyr the contribution by young stars increases rapidly. Simultaneously, the $`(RI)`$ colours become significantly bluer. After burst maximum, the contribution by young stars decreases again and galaxy colours start to redden until the effect of the burst becomes negligible and colours and magnitudes evolve as prior to the burst.
In this manner, a characteristic burst loop is completed as shown in Fig. 11a. The shape of this loop is determined by: i) the burst amplitude, ii) the extinction within the Hii regions, iii) the contribution by the old stellar population to the integrated galaxy light, iv) the duration (and profile) of the burst, v) the maximum time during which young stars produced by the Hii regions can still be distinguished from the surrounding field population.
Burst amplitude and IMF: Fig. 11a demonstrates that values of $`\eta _I0.25`$ can be well explained by bursts with amplitude $`A_\mathrm{b}=8`$ $`M_{}`$yr<sup>-1</sup> superimposed on an exponentially decreasing SFR model ending at $`\mu _1=0.1`$. Decreasing the burst amplitude by factors 2.5 and 5, respectively, results in the smaller loops shown in Fig. 11a. The actual burst amplitude required to explain the observations depends on many quantities as described below
Dust extinction: Fig. 11b illustrates the effect of dust extinction on $`\eta _I`$ and $`(RI)`$. A selective extinction of $`E_{BV}=(0.25,0.5,1)`$ mag for the Hii regions results in a reduction $`\eta _I`$ by a factor (1.6, 2.5, 6.5) and a reddening $`E_{RI}=`$ (0.14, 0.27, 0.56) mag, assuming a Galactic extinction curve. For values of $`E_{BV}\begin{array}{c}>\\ \end{array}0.5`$ mag, we find that the bluing effect on $`(RI)`$ by young massive stars formed during the burst is neutralised almost entirely by extinction. For intense bursts reddening of the galaxy colours may even occur. If variations in the mean extinction of the ensembles of Hii regions among LSB galaxies are small (e.g. less than a factor two), it is difficult to see how extinction alone can provide an adequate explanation for the large variations observed in $`\eta _I`$.
Global star formation history: From Fig. 11c it is clear that the burst effect on the galaxy magnitudes and colours increases when the contribution by the old stellar population is decreased. Thus, colours and magnitudes are less affected by bursts imposed on constant or even increasing SFRs compared to those imposed on exponentially decreasing SFR models: the smaller loop sizes just reflect that the mean age of the underlying stellar population is relatively low.
Current gas fraction: Fig. 11d demonstrates how the present-day gas fraction $`\mu _1`$ affects the effect of the burst. The luminosity contribution by the old stellar population decreases for increasing values of $`\mu _1`$, and thus, as discussed above, we find that the effect of a star formation burst for galaxies with $`\mu _1=0.9`$ (i.e. unevolved systems) is as large as that of a ten times stronger burst for $`\mu _1=0.1`$ (i.e. highly evolved systems). Thus, the burst amplitude required to explain the observations strongly depends on the present-day gas-to-total mass-ratio.
Burst duration: Fig. 11e shows that the duration of the burst affects the impact of the burst as well. For $`\mathrm{\Delta }t_\mathrm{b}=1`$, 5, and 10 Myr, the variation in $`\eta _I`$ is $`0.25`$ while the resulting galaxy $`(RI)`$ colours become bluer. Burst durations in excess of $`\mathrm{\Delta }t_\mathrm{b}5`$ Myr are unlikely since this would require dust extinctions $`E_{(BV)}\begin{array}{c}>\\ \end{array}1`$ mag in order to provide agreement with the observed $`(RI)`$ colours (cf. Fig. 11b). Such large extinction in Hii regions are probably excluded by the observations (e.g. McGaugh 1994). Thus, relatively narrow burst profiles are needed to explain extreme values $`\eta _I0.2`$ (as in F568-V1; cf. Table 4).
Maximum age of Hii regions: Fig. 11f shows that when the maximum lifetime of the Hii regions is increased from $`\tau _{\mathrm{HII}}=5`$ to 50 Myr, partial agreement with the observations can be achieved without invoking star formation bursts. We note that the resulting Hii region contributions do not increase linearly with $`\tau _{\mathrm{HII}}`$ as short lived massive stars dominate the luminosity contribution of all stars formed during the past $`\tau _{\mathrm{HII}}`$ yr.
Values of $`\tau _{\mathrm{HII}}\begin{array}{c}>\\ \end{array}50`$ Myr would mean that stars down to masses of $`m_{\mathrm{ion}}\begin{array}{c}<\\ \end{array}7`$ $`M_{}`$ would contribute to the Hii regions identified (e.g. Schaerer et al. 1993; $`Z=0.001`$). Observational estimates for $`m_{\mathrm{ion}}`$ are usually in the range $`1015`$ $`M_{}`$ (Wilcots 1994; García-Vargas 1995) and correspond to $`\tau _{\mathrm{HII}}\begin{array}{c}>\\ \end{array}15`$ Myr. Even though these values would imply that our adopted value of $`\tau _{\mathrm{HII}}=5`$ Myr is too low, this excludes extreme values of $`\tau _{\mathrm{HII}}=200`$ Myr which would be required to explain the observed range in $`\eta _I`$ exclusively in terms of variations in $`\tau _{\mathrm{HII}}`$ and/or $`m_{\mathrm{ion}}`$. Therefore, variations in $`\tau _{\mathrm{HII}}`$ (or equivalently $`m_{\mathrm{ion}}`$) and/or extinction may provide an explanation only for part of the variations in the Hii region integrated luminosity contributions observed among LSB galaxies.
In conclusion, the observations suggest that small amplitude, short bursts of star formation are important in at least several of the LSB galaxies for which accurate photometry data is available. Such recent episodes of enhanced star formation may play an important role in affecting the colours of the blue LSB galaxies.
### 6.3 Quantitative effect of small amplitude bursts on galaxy colours and magnitudes
Table 5 lists the effect of a 5 Myr star formation burst on the galaxy integrated magnitudes and colours for various burst amplitudes. For an exponentially decreasing SFR (model A1 in Table 3; assuming $`M_\mathrm{g}(0)=10^{10}`$ $`M_{}`$, $`\mu _1=0.1`$, Salpeter IMF), we find that a burst with amplitude $`A_\mathrm{b}=0.8`$ $`M_{}`$ yr<sup>-1</sup> results in maximum colour variations $`\mathrm{\Delta }(BV)`$ and $`\mathrm{\Delta }(RI)`$ of $`0.18`$ and $`0.04`$ mag, respectively. The effect of increasing the burst amplitude by a factor 10 to $`A_\mathrm{b}=8`$ $`M_{}`$ yr<sup>-1</sup> results in corresponding shifts of $`0.56`$ and $`0.26`$ mag, respectively. This effect is similar to that when the initial galaxy mass is reduced by a factor ten while leaving the burst amplitude unaltered (i.e. for. $`M_\mathrm{g}(0)=10^9`$ $`M_{}`$ and $`A_\mathrm{b}=0.8`$ $`M_{}`$ yr<sup>-1</sup>).
For bursts superimposed on exponentially decreasing SFRs models ending at $`\mu _1=0.1`$, the colour and magnitude shifts predicted are consistent with the observations in case of burst amplitudes $`A_\mathrm{b}\begin{array}{c}<\\ \end{array}3`$ $`M_{}`$ yr<sup>-1</sup>, assuming a typical extinction of $`E_{(BV)}=0.3`$ mag in the Hii regions. In fact, the effect of the burst is determined mainly by the total luminosity of the young stellar populations formed according to the continuous SFR during the last Gyr or so. Since the amplitude of the SFR scales with $`(1\mu _1)`$ for models ending at different gas fractions $`\mu _1`$ the impact of the burst for models ending at $`\mu _1=0.5`$ is about twice that given in Table 5. Similarly, for models ending at $`\mu _1\begin{array}{c}>\\ \end{array}0.9`$, the burst effect becomes roughly ten times stronger compared to the $`\mu _1=0.1`$ case (cf. Fig. 11). The effect of the burst is substantially reduced when going from exponentially decreasing to constant SFRs (cf. Table 5).
## 7 Present-day star formation rates in LSB galaxies
### 7.1 Theoretical star formation rates
#### 7.1.1 Continuous SFR
Theoretical SFRs can be derived from the models discussed in the previous section (exponentially decreasing SFR, Salpeter IMF), using:
$$\mathrm{SFR}^{\mathrm{cont}}A(\mu _1)\frac{M^{\mathrm{tot}}}{10^{10}}M_{}\mathrm{yr}^1$$
(3)
where $`A(\mu _1)`$ is the model SFR amplitude required to end at a gas fraction $`\mu _1`$ at a galactic evolution time of 14 Gyr (assuming an initial mass of $`10^{10}`$ $`M_{}`$) and $`M^{\mathrm{tot}}`$ the total galaxy mass as obtained from $`M_{\mathrm{HI}}`$ and $`\mu _1`$.
We list the required values for $`A`$ to have the models end at certain values of $`\mu _1`$ in Table 6. Using these values we find that LSB galaxies show present-day SFRs (without bursts) between $`0.01\begin{array}{c}<\\ \end{array}\mathrm{SFR}^{\mathrm{cont}}\begin{array}{c}<\\ \end{array}0.15`$$`M_{}`$yr<sup>-1</sup>. For a typical LSB galaxy (i.e. $`M^{\mathrm{tot}}=10^{10}`$ $`M_{}`$, $`\mu _1=0.5`$) we estimate SFR$`{}_{}{}^{\mathrm{cont}}=0.1`$ $`M_{}`$yr<sup>-1</sup>.
#### 7.1.2 Burst SFR
As we discussed in Sect. 6, the effect of a starburst on the galaxy integrated magnitudes and colours depends on many quantities. However, assuming a modest Hii region extinction of $`E_{(BV)}=0.25`$ mag and lifetime $`\tau _{\mathrm{HII}}=20`$ Myr, crude estimates for the maximum burst amplitude can be derived from:
$$\mathrm{SFR}^{\mathrm{burst}}20\eta _\mathrm{I}\frac{M_{\mathrm{HI}}}{10^{10}}\left(\frac{1}{\mu _{\mathrm{rot}}}1\right)M_{}\mathrm{yr}^1$$
(4)
We find that $`I`$-band contributions of $`\eta _\mathrm{I}=0.2`$ are best explained by a 5 Myr burst with amplitude SFR$`{}_{}{}^{\mathrm{burst}}=0.8`$ $`M_{}`$yr<sup>-1</sup>, assuming present-day values of $`M_{\mathrm{HI}}=210^9`$ $`M_{}`$ and $`\mu _{\mathrm{rot}}=0.5`$.
Theoretical estimates for the total current SFRs in LSB galaxies are found from $`\mathrm{SFR}^{\mathrm{tot}}=\mathrm{SFR}^{\mathrm{cont}}+\mathrm{SFR}^{\mathrm{burst}}`$ and are listed in the last two columns of Table 7. Present-day SFRs for LSB galaxies experiencing small amplitude bursts range from SFR$`{}_{}{}^{\mathrm{tot}}0.02`$ to 0.8 $`M_{}`$ yr<sup>-1</sup>.
### 7.2 Empirical star formation rates
We have also derived current SFRs for LSB galaxies using the empirical method presented by Ryder & Dopita (1994; hereafter RD) based on CCD surface photometry of galactic disks. These authors found a relationship between the local H$`\alpha `$ and $`I`$-band surface brightness in the disks of a sample of 34 of southern spiral galaxies. From this relation, RD derived a constraint on the present-day SFR integrated over the entire stellar mass spectrum as:
$$\mathrm{log}\mathrm{SFR}_1=0.26\mu _I+0.92\mathrm{log}\sigma _{\mathrm{HI}}+5.3$$
$$M_{}\mathrm{pc}^2\mathrm{Gyr}^1$$
(5) where $`\mu _I`$ is the $`I`$-band surface brightness and $`\sigma _{\mathrm{HI}}`$ is the global mean Hi surface density \[$`M_{}`$pc<sup>-2</sup>\] within the star-forming disk. The relation between SFR<sub>1</sub> and $`\mu _I`$ is normalised by a term related to the mean surface density $`\sigma _{\mathrm{HI}}`$ and by a constant which is partly related to the conversion of the massive star formation rate to total SFR depending on the adopted IMF (cf. Kennicutt 1983). It is unclear from the RD sample whether the relation is valid also for the lowest (stellar) $`I`$-band surface densities that are observed among LSB galaxies. However, since this relation holds over a wide range in surface brightness and massive star formation in the disks of spirals it appears rather insensitive to galactic dynamics, extinction, and molecular gas content (RD), we expect this relation to be valid also in case of LSB galaxies. At the faintest surface brightnesses (i.e. $`\mu _I\begin{array}{c}>\\ \end{array}25.6`$ mag arcsec<sup>-2</sup>), the relation may be flattening off although the effects of sky subtraction and small number statistics leave this open to question. If flattening indeed occurs, Eq. (5) provides lower limits to the actual SFRs in LSB galaxies.
Using Eq. (3), we estimate global present-day SFRs for all LSB galaxies in our sample with measured $`I`$-band magnitudes and related data. For these LSB galaxies we list the distance, apparent $`I`$-band magnitude, and radius $`R_{27}`$ of the 27 mag arcsec<sup>-2</sup> $`B`$-band isophote, in columns (2) to (4) in Table 7. This radius corresponds to the optical edge of the LSB galaxy and is more representative for the radius within which the old disk stellar population in LSB galaxies is contained than is $`R_{25}`$ as used by RD for HSB galaxies. Accordingly, we define an effective I band surface brightness as:
$$\mu _I^{\mathrm{eff}}=m_I+2.5\mathrm{log}(\pi R_{27}^2)$$
(6)
and use this in Eq. (5). We tabulate the outermost radius of the measured Hi rotation curve $`R_{\mathrm{HI}}`$, effective $`I`$ band surface brightness $`\mu _I^{\mathrm{eff}}`$, total Hi mass derived within $`R_{\mathrm{HI}}`$, and the mean global surface Hi density $`\sigma _{\mathrm{HI}}`$ in columns (5) to (8), respectively. For a consistent comparison between effective surface brightness and global Hi surface density, one ought to measure them out to the same radius (e.g. $`R_{27}`$). However, since $`M_{\mathrm{HI}}`$ has been measured within $`R_{\mathrm{HI}}`$, we expect the former values to be more representative of the average Hi surface density in the star forming part of the disk. The mean global Hi surface densities $`\sigma _{\mathrm{HI}}`$ derived using $`R_{\mathrm{HI}}`$ vary between 2 and 5.5 $`M_{}`$pc<sup>-2</sup> and are substantially smaller (i.e. by $`2060`$ %) than those derived using $`R_{27}`$ instead.
The empirically derived mean present-day SFRs for the LSB galaxies from our sample range from about SFR$`{}_{1}{}^{}=0.02`$ to $`0.2`$ $`M_{}`$yr<sup>-1</sup> (cf. column 10 of Table 7). Errors arising from the Hi normalisation are estimated to be within a factor of two. This is illustrated when the same SFRs are derived assuming a fixed Hi surface density of 2 $`M_{}`$pc<sup>-2</sup> for all LSB galaxies (cf. SFR$`{}_{}{}^{\mathrm{fix}}{}_{1}{}^{}`$ in Table 7).
The empirically derived current SFRs in LSB galaxies are in good agreement with the theoretically derived SFRs ranging from SFR$`{}_{}{}^{\mathrm{cont}}0.01`$ to 0.15 $`M_{}`$yr<sup>-1</sup> (cf. Table 7). As discussed before, the theoretically derived present-day SFRs of individual LSB galaxies lie probably between SFR<sup>cont</sup> and SFR<sup>tot</sup> where the latter values include the contribution of small amplitude bursts (SFR$`{}_{}{}^{\mathrm{tot}}=`$SFR$`{}_{}{}^{\mathrm{cont}}+`$SFR<sup>burst</sup>). In several cases, the values of SFR<sup>tot</sup> are considerably larger than the empirical values which suggests that the burst contribution is overestimated by the models and/or that the SFRs derived empirically do not trace well local enhancements of star formation at the faint surface brightnesses of LSB galaxies.
The present-day SFRs in LSB galaxies are thus found to be considerably below the $`510`$ $`M_{}`$yr<sup>-1</sup> derived for their HSB counterparts (e.g. Kennicutt 1992 and references therein) but significantly larger than the $`0.001`$ $`M_{}`$yr<sup>-1</sup> observed typically in dwarf irregular galaxies (e.g. Hunter & Gallagher 1986).
## 8 Discussion
In the previous sections we have focussed on the difference between HSB galaxies and the type of LSB galaxies we have in our sample. As noted in the introduction (and also shown in Fig. 1) HSB and LSB galaxies, however, do not form separate groups. In fact there is a broad range of galaxy properties and the purpose of the discussion is to clarify what role the star formation history plays in determining the final character of a galaxy.
In general, there appears a trend along the Hubble sequence from rapidly decaying SFRs for early type galaxies to constant or even increasing SFRs for dwarf irregular galaxies (see e.g. reviews by Sandage 1986; Kennicutt 1992). On average, the observed trend corresponds to a decrease of the ratio of mean past to present SFR along the Hubble sequence. Most of the LSB galaxies in our sample belong to the group of late-type galaxies for which exponentially decreasing SFR models are in good agreement with the observations. The remaining LSB galaxies, for which slowly decreasing or constant (sporadic) SFR models are more appropriate, belong to a group of galaxies with properties intermediate to those of disk galaxies with weak or absent spiral arms and Sm/Im galaxies. Thus, in general LSB galaxies comply well with the observed trend of SFR variation with Hubble type.
The “burst” scenario discussed in section 6 indicates that current star formation in virtually all the LSB galaxies in our sample is local both in time and space and suggests that sporadic star formation has been a continuous process from the time star formation started in the disks of LSB galaxies.
The low star formation rates of $``$0.1 $`M_{}`$yr<sup>-1</sup> derived for LSB galaxies as well as the local nature of the star formation in these systems are consistent with the idea of a critical threshold for the onset of global star formation in disk galaxies (e.g. Skillman 1987; Kennicutt 1989; Davies 1990).
Even though LSB galaxies contain large amounts of gas, only very limited amounts participate in the process of star formation. If we assume that LSB galaxies maintain their current star formation rate of $``$0.1 $`M_{}`$yr<sup>-1</sup>, their typical present-day amount of gas $``$ M<sub>g</sub> = 2.5 10<sup>9</sup> $`M_{}`$ will be consumed within $`\tau _{\mathrm{gas}}`$30 Gyr (for a recycled fraction of 25%). Such gas consumption times for LSB galaxies are much larger than a Hubble time (e.g. Romanishin 1980). For comparison, $`\tau _{\mathrm{gas}}24`$ Gyr in HSB galaxies, assuming typically M$`{}_{\mathrm{g}}{}^{}10^{10}`$ $`M_{}`$ and SFR$`{}_{1}{}^{}5`$ $`M_{}`$yr<sup>-1</sup>, which implies that HSB galaxies will run out of gas soon (see Kennicutt 1992).
The presence of an old stellar population in many late-type LSB galaxies as indicated by their optical colours (e.g. vdH93; dB95) and as confirmed by the results from the modelling suggests that LSB galaxies roughly follow the same evolutionary history as HSB galaxies, but at a much lower rate.
The models suggest that very slowly decreasing star formation rates with $`e`$-folding times much larger than $``$ 5 Gyr probably can be ruled out for the majority of the LSB galaxies examined here. The observed colours are too blue, provided extinction is indeed negligible. This result is insensitive to the possible occurrence of infall or accretion of gas. Similarly the models indicate that the dominant stellar population in galaxies with $`(BV)0.5`$ mag cannot be as young as 10 or 5 Gyr (cf. Fig. 5). In contrast the models can not rule out a dominant luminosity contribution by stellar populations significantly younger than 5 – 10 Gyr for LSB galaxies with $`(BV)0.4`$.
First this indicates that the mean age of the stellar population in most LSB and HSB galaxies is similar even though the disks of LSB galaxies are in a relative early evolutionary stage. Although we have not explored an entire range of star formation e-folding times, values much below or in excess of $``$ 5 Gyr are not very likely. Smaller values would increase the colour contribution from the underlying older stellar population and hence produce redder colours than observed, while much larger values correspond to slowly decreasing or almost constant star formation rates implying a relatively larger contribution of the younger stellar population.
Secondly, the combined effect of extinction and metallicity on galaxy colours is sufficient to explain the colour differences observed between LSB galaxies and HSB galaxies. Since the amount of extinction depends strongly on the dust content, which in turn is coupled to the heavy element abundances in the ISM (see Sect. 2), metallicity probably is the main quantity determining the colour differences observed between LSB and HSB galaxies. In this manner the much lower rate of star formation in LSB galaxies, implying lower metallicities and dust contents, indirectly determines the blue colours of LSB galaxies compared to HSB galaxies.
## 9 Summary
We have examined the star formation histories of LSB galaxies using models which take into account both the photometric and the chemical evolution of the galaxies. For the majority of the LSB galaxies in our sample, observed $`UBVRI`$ magnitudes, \[O/H\] abundances, gas masses and fractions, and Hi mass-to-light ratios are best explained by an exponentially decreasing global SFR ending at a present-day gas-to-total mass-ratio of $`\mu _10.5`$. When small amplitude bursts are involved to decrease the predicted $`M_\mathrm{g}/L`$ ratios, models ending at $`\mu _10.7`$ may also apply. In addition to exponentially decreasing SFR models, $`15\%`$ of the LSB galaxies require modest amounts of internal extinction $`E_{(BV)}\begin{array}{c}<\\ \end{array}0.1`$ mag to explain the relatively red colours of $`(BV)0.6`$ of these systems.
A substantial fraction ($`35\%`$) of the LSB galaxies in our sample have colours $`(BV)\begin{array}{c}<\\ \end{array}0.45`$ mag and exhibit properties similar to those resulting from exponentially decreasing SFR models at evolution times of $`510`$ Gyr. Alternatively, recent episodes of enhanced star formation superimposed on exponentially decreasing SFR models may provide an adequate explanation for the colours of these systems (see Sect. 6). Recent star formation is observed, at least at low levels, in essentially all the LSB galaxies in our sample. Hence, it seems justified to assume that the disks in LSB galaxies experienced continuous (i.e. frequent small amplitude bursts of) star formation, at least during the last few Gyr.
There is nothing peculiar about the evolution of LSB galaxies. Broadly speaking their evolution proceeds like that of HSB galaxies, but at a much lower rate.
###### Acknowledgements.
We thank Andre Maeder and Georges Meyanet for providing us with the stellar isochrone data and conversion programs. We thank Roelof de Jong for making available to us the data of a large sample of face-on spirals. We are grateful to the referee Uta Fritze-von Alvensleben for constructive comments from which this paper has benefitted.
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# Structure Factor of a Lamellar Smectic Phase with Inclusions
## Abstract
Motivated by numerous X-ray scattering studies of lamellar phases with membrane proteins, amphiphilic peptides, polymers, or other inclusions, we have determined the modifications of the classical Caillé law for a smectic phase as a function of the nature and concentration of inclusions added to it. Besides a fundamental interest on the behavior of fluctuating systems with inclusions, a precise characterization of the action of a given protein on a lipid membrane (anchoring, swelling, stiffening …) is of direct biological interest and could be probed by way of X-ray measurements. As a first step we consider three different couplings involving local pinching (or swelling), stiffening or tilt of the membrane. In the first two cases we predict that independent inclusions induce a simple renormalization of the bending and compression modulii of the smectic phase. The X-ray experiments may also be used to probe correlations between inclusions. Finally we show that asymmetric coupling (such as a local tilt of the membrane) results in a modification of the usual Caillé law.
Lyotropic smectic phases ($`L_\alpha `$ phases) are liquid crystalline systems of parallel, regularly spaced lipid bilayers. Those phases form a quasi-crystalline structure in one (the $`z`$) direction, while retaining their fluid properties in the two others directions ($`xy`$ plane). They have been studied extensively in recent years, as they form one of the most convenient lyotropic structure for the experimental study of fluctuating interfaces. They have been for instance particularly useful in demonstrating the importance of the so-called entropic (Helfrich) repulsion between fluctuating membrane. Furthermore, the degree of dilution (the layer spacing) and other controlled parameters of the lyotropic $`L_\alpha `$ phase can be experimentally varied over a wide range, which allows for a precise study of the scattering intensity of the phase. This scattering intensity is expected to show several peaks corresponding to a stack of regularly spaced membrane. In these layered systems however, the thermal fluctuations destroy the long range order at finite temperature, leading to weak - power law - singularities of the peaks.
Inclusions in membrane have been actively studied for many years. Most of the earlier work was done by the biological community, for which it was a necessary step toward the comprehension of the complexity of biological (cell) membrane. In the last twenty years or so, a more physical description of membrane-inclusion complexes, in term of curvature energy and disruption of the molecular order of the membrane, have led to interesting results on membrane-induced interactions between inclusions. While this “physical” picture of the membrane-inclusion interactions leaves aside many fundamental processes driven by specific interactions, it allows for a description of universal features of the complexes in terms of elastic couplings between membrane and inclusion. This description is applicable to a certain extent to membrane with adsorbed or grafted polymers or colloids, as well as to transmembrane inclusions resembling certain membrane proteins.
In recent years, numerous experimental studies have investigated the role of added membrane proteins in lyotropic smectic phase. This lipid phase seems to be of particular interest for the study “in vitro” of the bilayer-protein association, and may have some practical application for gene therapy. From a physicist point of view, it is interesting to understand the effect of inclusions on the thermal fluctuations of flexible, fluid membranes. From a more “biological” point of view, a precise characterization of the physical action of a given protein onto a lipid membrane (anchoring, swelling, stiffening …) is of direct biological interest, and could be reach via X-ray measurements. This work have some relevance in material science as well, since for practical purposes, a lamellar phase is very rarely composed of two components (solvent+surfactants) only. The adjunction of a third or fourth component, such as cosurfactant, polymers, proteins or colloids is of common practice to tune the properties of the phase. The object of the present work is to determined the modifications of the smectic X-ray scattering due to the presence of inclusions in the lamellar phase. Ultimately, one could expect to map out different membrane-inclusion couplings by their effect on the X-ray structure factor. We consider three model inclusions, leading to different modifications of the membrane elastic properties, and we show that the qualitative analysis of the diffraction spectrum is enough to discriminate between these inclusions.
Several recent theoretical works deal with lamellar systems with flexible polymers or other soluble inclusions (so-called doped solvent lamellar phases). These works are complementary, as we are interested in insoluble inclusions such as the membrane proteins used in , while studied effects such as polymer depletion - induced interactions between lamellae, resulting from an equilibrium between the trapped inclusions and free inclusions in pure solvent.
$`\mathrm{}`$ The standard description of Caillé for pure $`L_\alpha `$ phases has been confirmed experimentally in numerous systems. It is based on the classical smectic hamiltonian:
$$H_0=\frac{1}{2}d^3r\left(B\left(_zu\right)^2+K\left(_{}^2u\right)^2\right)$$
(1)
where $`u(𝐫)`$ is the (continuous) normal displacement of the layer at the point $`𝐫`$, $`z`$ is the coordinate normal to the layers, and $`_{}`$ is the gradient along the layers. The elasticity of the smectic phase is characterized by the two coefficients $`K`$ and $`B`$, which are related to the energy cost of bending and compressing the sample. Typically, $`K=\kappa /d`$, where $`\kappa `$ is the bending modulus of a single membrane and $`d`$ is the average layer spacing, and $`B`$ is function of the intermembrane interactions. Those modulii define the characteristic smectic length $`\lambda =\sqrt{K/B}`$. The scattering intensity $`I(𝐪)`$ is related to the Fourier transform of the density correlation function $`G_n(𝐫)`$, itself related to the exponential of the layer displacement correlation function $`g(𝐫)=q_0^2/2\left(u(𝐫)u(0)\right)^2`$, where $`\mathrm{}`$ represents a thermal average of the fluctuations, and $`q_0=2\pi /d`$ (see , sec 6.3.2). As usual, the Fourier decomposition of the layer displacement is used: $`u(𝐫)=d^3q/(2\pi )^3u_qe^{i\mathrm{𝐪𝐫}}`$. Using Eq.(1):
$$g(𝐫)=\frac{q_0^2}{V}\frac{d^3q}{(2\pi )^3}|u_q^2|\left(1e^{i\mathrm{𝐪𝐫}}\right)\mathrm{with}|u_q^2|_0=\frac{Vk_BT}{Bq_z^2+Kq_{}^4}$$
(2)
where $`V`$ is the volume of the sample. After integration over $`q_z`$, an asymptotic solution of $`g(𝐫)`$ can be found:
$`g^{(0)}(𝐫)={\displaystyle \frac{q_o^2k_BT}{4\pi \sqrt{KB}}}{\displaystyle 𝑑q_{}\frac{\left(1J_0(q_{}r_{})e^{\lambda zq_{}^2}\right)}{q_{}}}\eta \left(2\gamma +2\mathrm{log}{\displaystyle \frac{\pi r_{}}{l}}+E_1\left({\displaystyle \frac{r_{}^2}{4\lambda z}}\right)\right)`$ (3)
$`\mathrm{with}\eta {\displaystyle \frac{q_0^2k_BT}{8\pi \sqrt{KB}}}`$ (4)
where $`J_0(x)`$ is the Bessel function of the first kind, $`\gamma `$ is the Euler constant, and $`E_1(x)`$ is the exponential integral ($`l`$ is a molecular size). From Eq.(4), one can show that the scattering intensity $`I_n(𝐪)`$ near the n<sup>th</sup> peak, defined by
$$I_n(𝐪)=G_n(𝐪+n𝐪_0)+G_n(𝐪n𝐪_0)\mathrm{with}G_n(𝐫)e^{n^2g(𝐫)}$$
(5)
shows power law singularities of the form
$$G_n(𝐪n𝐪_0)\{\begin{array}{cc}(q_znq_0)^{2+\eta }\mathrm{if}q_{}=0\hfill & \\ q_{}^{4+2\eta }\mathrm{if}q_z=0\hfill & \end{array}$$
(6)
The above law gives a satisfactory description of the X-ray structure factor in the vicinity of the diffraction peaks, and has been observed experimentally (see and references therein). A more accurate calculation of $`I(𝐪)`$ over a broad range of wave vector requires a discrete description of the lamellar phase, and the account of the fluctuations of the lipid concentration, and is not tractable analytically. As one can see, this peculiar power law divergence allows for an experimental determination of the product $`KB`$, through the parameter $`\eta `$. Low values of the parameter $`\eta `$, such as those observed for electrostatically stabilized systems ($`0.2`$) correspond to a well defined smectic organization with little fluctuations. Larger values of order unity are observed for weakly interacting neutral or screened systems stabilized by the Helfrich interaction.
$`\mathrm{}`$ Several couplings between the inclusions and the membrane are considered below. They are referred to as “pinch”, “stiff”, and “tilt”, and are schematically depicted in Fig.1 and Fig.2. The “pinch” inclusion can be thought of as exerting a force which tends to pinch ($`\beta _{pinch}>0`$, see Eq.(7) below) or swell ($`\beta _{pinch}<0`$) neighboring membranes. Since we are using a continuous description for the lamellar phase, the inclusion may reside within a membrane or between two membranes. Inclusions which would typically induce such deformations are amphiphilic proteins, with hydrophobic parts which penetrate in the bilayers and hydrophilic parts lying in the water, or simply particles larger than the layer spacing which sterically swell the membrane. A pinched lamellar structure has also been identified in the case of anionic polymers added to a lamellar phase formed by a mixture of cationic and neutral lipids. The hamiltonian corresponding to this perturbation couples the density of inclusion $`\rho (𝐫)`$ with the variation of the layer spacing:
$$\mathrm{\Delta }H_{pinch}=d^3r\beta _{pinch}\rho (𝐫)_zu$$
(7)
The “stiff” inclusion corresponds to a transmembrane protein which locally stiffens (or softens) the membrane and modifies its fluctuations. A stiffening is expected for many integral membrane proteins, or for rigid inclusions laying on the lipid bilayer. On the other hand, pores in the membrane are expected to induce a softening of the membrane (among other, less trivial phenomena). The inclusion density is coupled with a variation ($`\delta K`$) of the bending modulus:
$$\mathrm{\Delta }H_{stiff}=\frac{1}{2}d^3r\delta Kv\rho (𝐫)\left(_{}^2u\right)^2$$
(8)
where $`v`$ is the effective volume of an inclusion in the membrane.
The “tilt” inclusion corresponds to an anisotropic inclusion (such as a grafted polymer) which would locally induce a curvature to a membrane without changing the layer spacing. It has been shown that such a deformation can be described by a modification of the “pinch” hamiltonian Eq.(7); we will get back to it later. It is clear that any real inclusion will, to a certain extent, combine the three effects described above.
A single inclusion of the “pinch” or “tilt” type induces a non-zero average deformation to the membrane. However, this mean deformation evens out during the course of an experiment, and is not detectable on the scattering spectrum (except if it induces a phase separation, such as expulsion of the solvent). The effect of the inclusions on the spectrum is visible via their effect on the membrane fluctuations. In what follows, we do not investigate a possible coupling between the membrane fluctuations and the inclusion correlation function (assumption of a quenched distribution of inclusions). It has been shown for a single membrane system that “quenched” and “annealed” distributions of inclusions give similar corrections to the fluctuation spectrum of the membrane. The scattering experiments show an average of the structure factor over the positions of the inclusions, it is thus necessary to calculate the membrane correlation function after a proper average over the possible positions of the inclusions.
$`\mathrm{}`$ For the case of “pinch” inclusions, the hamiltonian Eqs.(1,7) can be rewritten in the discrete form:
$`H_{pinch}={\displaystyle \frac{1}{V}}{\displaystyle \underset{q>0}{}}\left\{Q_q\left|u_qu_q^{}\right|^2\left(\xi _qu_q^{}+\xi _q^{}u_q\right)\right\}`$ (9)
$`\mathrm{with}Q_q=Bq_z^2+Kq_{}^4\mathrm{and}\xi _q=\beta _{pinch}iq_z\rho _q`$ (10)
The structure factor is then easily obtained:
$$\frac{|u_q|^2}{V}=\frac{T}{Q_q}+\frac{\xi _q\xi _q^{}/V}{Q_q^2}=\frac{k_BT}{Bq_z^2+Kq_{}^4}+\beta _{pinch}^2\frac{q_z^2}{\left(Bq_z^2+Kq_{}^4\right)^2}\frac{|\rho _q|^2_\rho }{V}$$
(11)
where $`\mathrm{}_\rho `$ designates an average over the position of the inclusions. The first part of the left hand side corresponds to the correlation function in a pure $`L_\alpha `$ phase $`g^{(0)}(𝐫)`$ (Eq.(2)), and the second part describes the corrections due to the inclusions $`\mathrm{\Delta }g`$. This correction reflects correlations between inclusions $`|\rho _q|^2_\rho `$.
If the inclusions are independent from each other, $`|\rho _q|^2/V`$ is the average concentration of inclusions $`\overline{\rho }`$. Similarly to Eq.(4), the correction to the spatial structure factor can be written:
$$\mathrm{\Delta }g(𝐫)=2\eta \frac{\beta _{pinch}^2\overline{\rho }}{2k_BTB}\left(𝑑q_{}\frac{\left(1J_0(q_{}r_{})e^{\lambda zq_{}^2}\right)}{q_{}}+\frac{1}{4}e^{\frac{r_{}^2}{4\lambda z}}\right)$$
(12)
This correction leads to the same power law divergence of the scattering intensity as the pure $`L_\alpha `$ (Eq(6)), with the effective modulii:
$$\left(KB\right)_{pinch}=\frac{KB}{\left(1+{\scriptscriptstyle \frac{\beta _{pinch}^2\overline{\rho }}{2k_BTB}}\right)^2}$$
(13)
One notice right away that the “pinch” inclusions tend to soften the lamellar phase. We thus expect broader peaks in the scattering intensity (see Fig.1). This finding is qualitatively consistent with experiments on comparable systems. Note that results reported for the doped solvent situation tend to show qualitatively similar results for physically different reasons. In that case, the inclusions-mediated attraction between lamellae tends to reduce the compression modulus $`B`$.
Possible correlations between inclusions may be induced either by direct interactions between them, or by membrane mediated interactions. In the Latter case, it has been shown that the inclusions interact via the potential (in Fourier space)
$$U_q=\frac{\beta _{pinch}^2}{2}\frac{q_z^2}{Q_q}$$
(14)
In this case, a ”Debye-Huckel like” expansion of the correlation between inclusions (valid if $`\overline{\rho }U_q/(k_BT)1`$), shows that:
$`{\displaystyle \frac{|\rho _q|^2}{V}}{\displaystyle \frac{\overline{\rho }}{1+\overline{\rho }U_q/(k_BT)}}\mathrm{or}{\displaystyle \frac{|\rho _q|^2}{V}}={\displaystyle \frac{\overline{\rho }Q_q}{Q_q^{}}}`$ (15)
$`\mathrm{with}Q_q^{}B^{}q_z^2+Kq_{}^4\mathrm{and}B^{}B\left(1{\displaystyle \frac{\overline{\rho }\beta _{pinch}^2}{2k_BTB}}\right)`$ (16)
From Eq.(11), the corrections to the structure factor is given by:
$$\mathrm{\Delta }|u_q|^2=\frac{\beta _{pinch}^2q_z^2\overline{\rho }}{Q_qQ_q^{}}\eta _{pinch}=\eta \left(1+\frac{\beta _{pinch}^2\overline{\rho }}{k_BT\sqrt{B^{}}(\sqrt{B}+\sqrt{B^{}})}\right)$$
(17)
The effective modulii are now given by:
$$\left(KB\right)_{pinch}=\frac{KB}{\left(1+\frac{\beta _{pinch}^2\overline{\rho }}{2k_BTB}+\frac{3}{16}\left(\frac{\beta _{pinch}^2\overline{\rho }}{2k_BTB}\right)^2\right)}$$
(18)
Note that in the limit of validity of the expansion, the correction due to correlation is quite small.
$`\mathrm{}`$ The investigation of the effect of the “stiff” inclusions is more delicate, because of the non-quadratic nature of the hamiltonian Eq.(8). This leads to a coupling between different Fourier modes of the membrane displacement, in which case Eq.(2) for the real space correlation function has to be replaced by:
$$g(𝐫)=\frac{1}{2}q_0^2\frac{d^3q}{(2\pi )^3}\frac{d^3q^{}}{(2\pi )^3}u_qu_q^{}^{}\left(e^{i\mathrm{𝐪𝐫}}1\right)\left(e^{i𝐪^{}𝐫}1\right)$$
(19)
The discrete version of the hamiltonian reads:
$`{\displaystyle \frac{H_{stiff}}{k_BT}}=`$ $`{\displaystyle \underset{q,q^{}>0}{}}\left(G_{q,q^{}}^0+\mathrm{\Delta }G_{q,q^{}}\right)u_qu_q^{}^{}`$ (20)
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{q,q^{}}{}}\left({\displaystyle \frac{Q_q}{Vk_BT}}\delta _{q,q^{}}+{\displaystyle \frac{\delta Kv}{k_BTV^2}}q_{}^2q_{}^{}{}_{}{}^{2}\rho _{(qq^{})}^{}\right)u_qu_q^{}^{}`$ (21)
from which the mean square displacement can be calculated:
$$u_qu_q^{}^{}=\left[\left(G^0+\mathrm{\Delta }G\right)^1\right]_{q,q^{}}$$
(22)
In the limit where the correction is only a small perturbation, an expansion of the previous equation $`[G^0+\mathrm{\Delta }G]_{q,q^{}}^1=\left[G^0\right)]^1_{q,q^{}}[\mathrm{\Delta }G]_{q,q^{}}\times [G^0]^2_{q,q^{}}`$ leads to:
$$\frac{u_qu_q^{}^{}}{V}\frac{k_BT}{Bq_z^2+Kq_{}^4}\delta _{q,q^{}}\frac{\delta Kvk_BT}{V}\frac{q_{}^2q^{}{}_{}{}^{2}\rho _{qq^{}}^{}}{(Bq_z^2+Kq_{}^4)(Bq_z^{}{}_{}{}^{2}+Kq_{}^{}{}_{}{}^{4})}$$
(23)
For an homogeneous distribution of particle of average concentration $`\overline{\rho }`$, the Fourier modes are decoupled: $`\rho _q=\overline{\rho }\delta _q`$. It is then easy to show that again, the scattering intensity show peaks with power law divergence, from which the product of the effective compression and bending modulii can be extracted. The renormalization introduced by the inclusions is given by:
$$\left(KB\right)_{stiff}=\frac{KB}{\left(1\frac{\delta Kv\overline{\rho }}{2K}\right)^2}$$
(24)
Note that this results compares to the finding of Netz and Pincus, who have shown that for a single membrane with inclusions, the effective bending rigidity of the membrane is given by $`\kappa _{eff}\kappa +\delta \kappa \varphi `$ at the lowest order in the volume fraction $`\varphi `$ of inclusions.
In contrary to the case of “pinch” inclusions, a stiffening of the $`L_\alpha `$ phase is observed for inclusions increasing locally the bending rigidity of the membrane. This should lead to a sharpening of the peaks in the scattering intensity. Hence, a quick look at the intensity scattered by the sample should give us valuable information on the inclusion-membrane coupling. A stiffening of a L<sub>α</sub> phase has been recently observed upon adjunction of polysoap molecules. This is a more complex case, since the polysoaps are fairly long molecules and can be expected to modify (increase) the interactions between membranes in addition to their effect on an isolated bilayers. The reduction of the effective $`\eta `$ has however been observed below the overlap concentration of polysoaps, where the molecules are expected to be spread flat onto the membrane.
$`\mathrm{}`$ The last coupling investigated in this work is the case of anisotropic inclusions which induce a local tilt (spontaneous curvature) on a membrane. It has been shown in a previous work that such a tilt can be conveniently mimic with the aid of an analogy between the “pinch/swell” inclusions and +/- electrical charges in electrostatics (see Fig.2). The deformation of the smectic induced by a “dipole” $`p_z\beta _{tilt}d`$ (where $`\beta _{tilt}`$, the tilt intensity, has the dimension of an energy) normal to the membrane has been calculated:
$$u_{tilt}=\frac{p_z}{16\pi \sqrt{KB}z^2}\left(1\frac{r_{}^2}{4\lambda z}\right)e^{r_{}^2/\left(4\lambda z\right)}$$
(25)
and correspond to imposing a local curvature $`C_0=\beta _{tilt}/(\pi \kappa d)`$ at the position of the inclusion ($`r_{}=0`$, $`z=d/2`$) without changing the layer spacing ($`\kappa =Kd`$ is the rigidity of an isolated membrane). It can be shown easily in the Fourier space that such a deformation can be induced by a coupling between the second order z-derivative of the layer deformation and the inclusion density:
$$\mathrm{\Delta }H_{tilt}=d^3rp_z\rho (𝐫)_z^2u$$
(26)
which, for independent inclusions, leads to the membrane structure factor::
$$\frac{|u_q|^2}{V}=\frac{k_BT}{Bq_z^2+Kq_{}^4}+p_z^2\frac{q_z^4}{\left(Bq_z^2+Kq_{}^4\right)^2}\overline{\rho }$$
(27)
and the correlation function:
$$\mathrm{\Delta }g(𝐫)=\frac{q_o^2p_z^2\overline{\rho }}{2\pi B^2}\frac{d^2q_{}}{2\pi }\frac{dq_z}{2\pi }\left(1e^{i\mathrm{𝐪𝐫}}\right)\frac{q_z^4}{(q_z^2+\lambda ^2q_{}^4)^2}$$
(28)
After integration over $`q_z`$ (the ($``$) indexes are forgotten):
$`\mathrm{\Delta }g(𝐫)={\displaystyle \frac{q_o^2p_z^2\lambda \overline{\rho }}{8\pi B^2}}`$ (29)
$`\mathrm{where}=3{\displaystyle 𝑑qq^3\left(1J_0(qr)e^{\lambda zq^2}\right)}+\lambda z{\displaystyle 𝑑qq^5J_0(qr)e^{\lambda zq^2}}`$ (30)
Both integrals can be computed after a series expansion of the Bessel function. The correction to the correlation function takes the form:
$$\mathrm{\Delta }g(𝐫)=\frac{q_o^2p_z^2\lambda \overline{\rho }}{16\pi B^2}\left(\frac{e^\omega }{(\lambda z)^2}\left(1+\omega \omega ^2\right)+\mathrm{constant}\right)\mathrm{with}\omega \frac{r_{}^2}{4\lambda z}$$
(31)
The X-ray spectrum is obtained by a combination of the above expression with Eqs.(4,5):
$`G_𝐪{\displaystyle d^3re^{i\mathrm{𝐪𝐫}}\left(\frac{l}{r_{}}\right)^{2\eta }\mathrm{exp}\left(\eta \left(E_1(\omega )+\left(\frac{l_{tilt}}{z}\right)^2e^\omega \left(1+\omega \omega ^2\right)\right)\right)}`$ (32)
$`\mathrm{with}l_{tilt}d\sqrt{{\displaystyle \frac{\beta _{tilt}^2\overline{\rho }}{2k_BTB}}}\mathrm{and}\omega {\displaystyle \frac{r_{}^2}{4\lambda z}}`$ (33)
Unlike the case of symmetrical inclusions considered earlier, the asymptotic scaling laws for $`G_𝐪`$ in the limits $`q_z=0`$ or $`q_{}=0`$ cannot be easily extracted from the above expression, since the exponential term in the integral cannot be written as a function of $`r_{}^2/(\lambda z)`$ alone. However, up to the first order in the dipole strength, $`l_{tilt}/z<1`$ for $`z>d/2`$ and we can expand the exponential in the integral Eq.(33). One can see clearly that unlike symmetrical inclusions such as “pinch” and “stiff”, the effect of a “tilt” inclusion cannot be matched by a renormalization of the smectic modulii, as they modify of the power law divergence of the scattering intensity near the Bragg positions. In fact, the inclusions add a term to the intensity, which for $`q_{}=0`$ can be written:
$`I_{tilt}(q_z)q_z^\eta l_{tilt}^2l^{2\eta }\lambda ^{1\eta }{\displaystyle }{\displaystyle \frac{dZ}{Z^{1+\eta }}}e^{iZ}{\displaystyle }{\displaystyle \frac{d\omega }{\omega ^\eta }}e^\omega (1+\omega \omega ^2)e^{\eta E_1(\omega )};Zq_zz`$ (34)
and in the other limit $`q_z=0`$:
$`I_{tilt}(q_{})q_{}^{2\eta }l_{tilt}^2l^{2\eta }\lambda {\displaystyle }{\displaystyle \frac{dR}{R^{1+2\eta }}}J_0(R){\displaystyle }d\omega e^\omega (1+\omega \omega ^2)e^{\eta E_1(\omega )};Rq_{}r_{}`$ (35)
The integrals being prefactors only, the scattering intensity should take the form:
$`I_{tilt}\{\begin{array}{cc}\alpha _1q_z^{2+\eta }\alpha _2q_z^\eta \hfill & \text{for }q_{}=0\hfill \\ \alpha _1q_{}^{4+2\eta }\alpha _3q_{}^{2\eta }\hfill & \text{for }q_z=0\hfill \end{array}`$ (36)
where the $`\alpha _i`$ are constants. The first term is the Caillé law for a pure smectic phase and the second term is the signature of the asymmetric inclusions. Note that the effect of the tilt inclusion as described by Eq.(36) is most probably very small and might be hardly detectable in experiments. Note also that as a “pinch” inclusion lead to a renormalization of the compression modulus $`B`$ (see Eq.(16)), a “tilt” inclusion can be seen as renormalizing the elastic modulus of an higher order term of the smectic hamiltonian, of the form $`K^{}(_z^2u)^2`$, which has been neglected in the reference hamiltonian Eq.(1). This reinforces the fact that their effect on the structure factor must be small. One nevertheless might hope to detect such an effect because of its dependence on the concentration of inclusion.
$`\mathrm{}`$ In conclusion, we have derived the effect of inclusions on the shape of the scattering curve of a smectic $`L_\alpha `$ phase. We have considered the three inclusions depicted in Figs.1 and 2, and have shown that the two symmetrical inclusions (“pinch” and “stiff”) renormalize the elastic modulii of the phase without modifying the power law divergence near the peaks (Caillé law). In agreement with experiments, we have obtained that the “pinch” inclusions soften the lamellar phase, while the “stiff” inclusions makes it more rigid. the “tilt” inclusions are predicted to modify the shape of the scattering curves, but their effect is small and might hardly be detectable.
As a last remark, we note that in view of our results, the study of the X-ray scattering spectrum of a lamellar phase with inclusions should be a powerful tool to investigate some changes in conformation of membrane proteins. One can argue that the bridge (pinch) between neighboring membranes created by some amphiphilic peptides and reported in may be pulled out of one of the two membranes as the intermembrane spacing is increased. If the peptide is flexible enough, it is then expected to lay flat on one membrane. This conformation change represents a transition from a “pinch” to a “stiff” inclusion, and should be directly detectable from the shape of the scattering peak. Moreover, some amphiphilic helical peptides in contact with a lipid bilayer are known to go from an adsorbed state at low peptide concentration, to an inserted state (a certain number of peptide aggregating to form an hydrophilic pore) at higher concentration. This transition from one state to the other could be interpreted in terms of “stiff” and “soft” (in the case of a pore) inclusions; One could investigate X-ray scattering from such a solution with this transition in mind. Nevertheless, one should always bear in mind that the different states of the amphiphilic peptides are most probably a complex mixed between the types of inclusion proposed in our work, and probably others.
We acknowledge stimulating discussions with W. Urbach, N. Taulier and N. Tsapis (E.N.S Paris), and we would like to thank J.-F. Joanny and A. Johner (I.C.S Strasbourg) for fruitful comments.
Fig.1 The two symmetrical inclusions considered in the paper. A “pinch” inclusion modifies
the layer spacing and broadens the Bragg peak - A “stiff” inclusion increase the local rigidity
and makes the Bragg peak narrower
Fig.2 a) The “tilt” inclusion induces a spontaneous curvature
b) How to make a “tilt” from the combination of “pinch” and “swell”
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# Nodal Quasiparticle Lifetime in the Superconducting State of 𝐁𝐢_𝟐𝐒𝐫_𝟐𝐂𝐚𝐂𝐮_𝟐𝐎_{𝟖+𝛿}
## Abstract
We have measured the complex conductivity, $`\sigma `$, of a $`Bi_2Sr_2CaCu_2O_{8+\delta }`$ (BSCCO) thin film between 0.2 and 0.8 THz. We find $`\sigma `$ in the superconducting state to be well described as the sum of contributions from quasiparticles, the condensate, and order parameter fluctuations which draw 30$`\%`$ of the spectral weight from the condensate. An analysis based on this decomposition yields a quasiparticle scattering rate on the order of $`k_BT/\mathrm{}`$ for temperatures below $`T_c`$.
The unusual properties of the quasiparticle (QP) lifetime $`\tau _{qp}`$ in high-$`T_c`$ cuprate superconductors provide evidence for a non-Fermi liquid normal state and a non-BCS mechanism. Because angle-resolved photoemission (ARPES) shows that QP properties are highly anisotropic, it is important to distinguish the behavior of $`\tau _{qp}`$ at different parts of the Fermi surface. Near the maximum of the d-wave gap at the ($`\pi `$, 0) point in momentum space, the ARPES lineshape narrows rapidly below $`T_c`$. The change suggests the onset of a well-defined QP in this antinodal region. Recently, improvements in detectors have made it possible to resolve the lineshape in the nodal or ($`\pi ,\pi `$) direction as well. Valla et al. reported that in the normal state, $`\tau _{qp}`$ can be described by the marginal Fermi liquid phenomenology, $`1/\tau _{qp}`$ max$`(\omega ,T)`$. Surprisingly, $`1/\tau _{qp}`$ of nodal quasiparticles (at the Fermi surface) remains approximately proportional to $`T`$ in the superconducting state as well. Thus the change in the lifetime of such QP’s upon entering the superconducting state is much less than that of the antinodal QPs.
The nodal QP’s apparent insensitivity to superconductivity is surprising for two reasons. First, one might expect the reduction in the phase space for scattering which accompanies the formation of a gap to increase the QP lifetime. Second, $`\tau _{qp}`$, as determined by microwave and thermal conductivity measurements in the closely related compound $`YBa_2Cu_3O_{7\delta }`$ (YBCO) does indeed increase rapidly below $`T_c`$. While such transport experiments are not momentum resolved, there is strong evidence that they probe the nodal QP lifetime. It was reported in Ref. that the microwave conductivity, $`\sigma =\sigma _1+i\sigma _2`$, can be described (at low $`T`$) by a two-fluid model which sums the contributions of condensate and QPs. The spectral weight in the normal fluid, $`\rho _n\sigma _1𝑑\nu `$, grows linearly with $`T`$ at low temperature at a rate which matches the linear decrease in superfluid density $`\rho _s(T)`$, measured independently via the penetration depth. The linear rate of transfer from $`\rho _s`$ to $`\rho _n`$ agrees with the theoretical prediction for nodal QPs in a d-wave superconductor.
Thus, the transport measurements of $`\tau _{qp}`$ in YBCO stand in apparent contradiction to ARPES measurements in BSCCO. It would appear that either the single-particle lifetime probed by ARPES is qualitatively different than the transport lifetime, or that $`\tau _{qp}`$ differs in BSCCO and YBCO despite their similar bilayer structure and $`T_c`$. The resolution of this issue requires interchanging the two probes and the two material systems. Since ARPES lineshapes in YBCO are difficult to obtain, it is important to focus on the transport properties of BSCCO, for example the microwave conductivity.
Unfortunately, it has proved difficult to determine $`\tau _{qp}`$ reliably from microwave data in BSCCO. In contrast to the strong dispersion in YBCO, $`\sigma _1`$ is approximately constant in the range from $`\nu `$=10-50 GHz. Consequently, the transport lifetime is not directly resolved in the conductivity spectra, although they do indicate that $`\tau _{qp}1/2\pi `$50 GHz, which is much smaller than in YBCO. Despite the lack of frequency dependence, in principle $`\tau _{qp}`$ could still be estimated from the magnitude of $`\sigma _1`$, which (in the $`\omega \tau _{qp}1`$ limit) is proportional to $`\rho _n\tau _{qp}`$. Using the two-fluid model, the magnitude of $`\rho _n(T)`$ follows from measurement of $`\rho _s(T)`$. In practice, this analysis yields results which are not self-consistent, for the following reason. Just as in YBCO, the condensate density varies as $`\rho _s(0)\alpha T`$ at low T, implying that the samples are in the clean d-wave regime where $`\rho _n=\alpha T`$. However, $`\sigma _1`$ in BSCCO does not tend to zero as $`T`$ approaches zero. Even at the lowest temperatures measured, about 5 K, $`\sigma _1`$ remains approximately 8 times larger than the normal state value just above $`T_c`$, $`\sigma _n(T_c)`$, and is far larger than the ‘universal conductivity’ proposed by Lee. Thus $`\tau _{qp}`$ calculated from $`\sigma _1/\rho _n`$ appears to diverge, which is inconsistent with the observation that the conductivity is frequency independent in the microwave regime. To extract a finite $`\tau _{qp}`$ one may assume that a large fraction of the quasiparticles remain uncondensed, i.e. $`\rho _n(T)=\rho _n(0)+\alpha T`$, or subtract the zero temperature limit of the dissipation prior to the two-fluid model analysis.
In view of the uncertainties inherent in these assumptions, it is desirable to extend the conductivity measurements to the frequency range probed by ARPES, $`k_BT/\mathrm{}`$. Here we report measurements of $`\sigma `$ in a BSCCO thin film in the frequency range from 0.18 to 1.0 THz, above the domain of microwave and below that of infrared spectroscopies. To cover this region, we used time-domain terahertz spectroscopy, a technique which is based on the generation and detection of single cycle electromagnetic pulses. The transmission coefficient of such pulses yields both the real and imaginary parts of $`\sigma `$ directly and independently. The sample, with $`T_c`$=$`85K`$, was grown on a $`LaAlO_3`$ substrate using atomic layer-by-layer molecular beam epitaxy. The resistance $`R`$ versus $`T`$ is linear and $`R(0)/R(300K)=4\times 10^2`$, where $`R(0)`$ is the extrapolation of the linear resistance to zero temperature.
The terahertz data show a strong frequency dependence in $`\sigma `$ below $`T_c`$, from which $`\tau _{qp}`$ in the superconducting state can be determined directly.
We find that $`1/\tau _{qp}(T)0.8k_B(T+10K)/\mathrm{}`$, consistent with ARPES measurements . Furthermore, we find the total conductivity to be incompatible with the two-fluid model. Instead, our measurements indicate that the conductivity comprises three contributions: QPs, the condensate, and a low frequency collective mode with spectral weight drawn from the condensate.
We now turn to $`\sigma _1(\omega ,T)`$ and discuss its salient features. The left panel of Figure 1 shows $`\sigma _1`$ plotted versus $`T`$ for a representative sample of frequencies in our range. At temperatures above $`T_c`$ there is little frequency dependence, as we would expect since $`1/\tau _{qp}2\pi `$1 THz in this temperature range. Below $`T_c`$, however, $`\sigma _1(T)`$ depends strongly on frequency. At $`\nu `$=0.2 THz the $`T`$ dependence is quite similar to that measured in the microwave regime on BSCCO single crystals . $`\sigma _1`$ rises to a maximum at 25 K, but at 5 K remains much greater than $`\sigma _n(T_c)`$. As the frequency increases from 0.2 to 0.8 THz the maximum shifts from 25 K to $`T_c`$, indicating that $`1/\tau _{qp}`$ is sweeping through our bandwidth. For a given frequency, the QP conductivity will decrease below the temperature where $`1/\tau _{qp}(T)\omega `$. At each measurement frequency, this $`T`$ is marked with an arrow. The measurement frequency (in radians/s) and corresponding $`T`$ for each arrow are shown in the upper right panel of Fig. 1. This admittedly rough estimate clearly suggests a QP lifetime with the form $`1/\tau _{qp}T`$. Finally, the peak in $`\sigma _1`$ near $`T_c`$, most visible at low frequencies, has been attributed to phase fluctuations of the order parameter due to thermal vortex, anti-vortex pairs.
Following its success in YBCO, we attempt to describe the conductivity using a two fluid model, in which $`\sigma `$ below $`T_c`$ is comprised of condensate and QP contributions. The condensate’s real, or dissipative, conductivity is a $`\delta `$-function at $`\nu =0`$. Therefore, for any nonzero frequency $`\sigma _1`$ is due solely to the QPs. Following the analysis of YBCO data, we choose a Drude conductivity for the QPs, $`\sigma _{qp}(\omega ,T)/\sigma _Q(\rho _n\tau _{qp}/\mathrm{})/(1+\omega ^2\tau _{qp}^2)`$. Here $`\sigma _Qe^2/(\mathrm{}d)`$ is the quantum conductivity of a stack of bilayers with spacing $`d=15.4\AA `$. The spectral weight, $`\rho _n`$ is expressed in units of energy as will be described below.
The two $`T`$ dependent parameters, $`\tau _{qp}(T)`$ and $`\rho _n(T)`$ are clearly suggested by the data. Regarding $`\tau _{qp}(T)`$, the upper right panel of Fig. 1 suggests that $`1/\tau _{qp}T`$ over a wide temperature range. The form of $`\rho _n(T)`$ follows from the $`T`$ dependence of the superfluid density. We obtain $`\rho _s(T)`$ from the imaginary part of the conductivity, $`\sigma _2`$, which is measured independently of $`\sigma _1`$ in our experiment. We observe that $`\rho _s(T)=\rho _0\alpha T`$ below about 30 K, suggesting that $`\rho _n=\alpha T`$ at low $`T`$. For purposes of modeling, we have found that the most accurate description of the terahertz conductivity results if we assume that $`\rho _n=\alpha T`$ up to $`T_c`$. Of course the total normal fluid weight must increase faster than this in order that $`\rho _n+\rho _s`$ remains constant as $`T`$ changes. However, we are modeling only the spectral weight contained in a Lorentzian component centered at zero frequency.
In the right-hand panel of Fig. 1, we plot $`\sigma _{qp}(T)`$, the Drude conductivity with $`\rho _n(T)`$ and $`\tau _{qp}(T)`$ as described above. Focusing first on the highest frequencies in our range, we see that $`\sigma _{qp}(T)`$ agrees rather well with the observed $`\sigma _1(T)`$. On the other hand, $`\sigma _1`$ at low frequency is much larger than predicted by a Drude model for the QP conductivity. In fact, we have found it impossible to describe the conductivity in the range from 0.2-1 THz with a two-fluid model ($`\rho _n(T)+\rho _s(T)`$ is constant) regardless of the choice of $`\tau _{qp}(T)`$. The difficulty is that the spectral weight in the low-frequency end of our range increases with decreasing $`T`$, rather than decreasing as expected for the normal fluid component in a two-fluid model. The conductivity spectra suggest the presence of an additional component, whose spectral weight increases with decreasing $`T`$. (Simulations based on show that the spectral weight at low frequency cannot be caused by inadvertent coupling to the c-axis plasmon due to off-normal incidence).
To help identify the additional component of $`\sigma `$, we consider the difference between the measured spectral weight and the amount ascribed to QPs using the parameters given above. Fig. 2 shows the spectral weight measured in our bandwidth, $`\mathrm{\Sigma }(T)`$, and that due to the normal fluid, $`\mathrm{\Sigma }_{qp}(T)`$, plotted versus $`T`$. Because of the extreme anisotropy of BSCCO we express spectral weight in terms of an areal density per CuO<sub>2</sub> bilayer. The areal density is proportional to the phase or charge stiffness per bilayer, and can be conveniently expressed in thermal units of energy, or Kelvins. The spectral weights denoted by $`\mathrm{\Sigma }`$ are distinct from those discussed earlier labeled with $`\rho `$’s. Whereas $`\rho (T)`$ represents an integration of $`\sigma _1`$ over all frequencies, $`\mathrm{\Sigma }(T)`$ is confined to an integration over our experimental bandwidth. The difference between $`\mathrm{\Sigma }(T)`$ and $`\mathrm{\Sigma }_{qp}(T)`$ is plotted as triangles. It is apparent that $`\mathrm{\Sigma }_{qp}(T)`$ is less than the total spectral weight observed. Notice that at low temperatures the difference, $`\mathrm{\Sigma }(T)\mathrm{\Sigma }_{qp}(T)`$, is proportional to $`\rho _s(T)`$, shown here multiplied by 0.18. The spectral weight in excess of $`\mathrm{\Sigma }_{qp}`$ is proportional to $`\rho _s`$ only if we set $`1/\tau _{qp}(T)T`$. If we choose $`1/\tau _{qp}(T)T^\beta `$ with any $`\beta `$ other than 1, the unaccounted for spectral weight does not have such a reasonable and recognizable $`T`$ dependence.
The proportionality of the residual spectral weight and $`\rho _s(T)`$, suggests that the excess conductivity arises from fluctuations of the condensate order parameter. The fluctuations are not thermally generated because their spectral weight increases as $`T0`$. Although the temperature dependence is consistent with intrinsic quantum fluctuations, the associated conductivity is expected to be of order $`\sigma _Q`$ or $`1.6\times 10^5(\mathrm{\Omega }^1m^1)`$, which is far too small to explain the data. However, the fluctuation conductivity can be much larger in the presence of static (or quasistatic) spatial variation in the superfluid density. Two examples have been discussed recently in the literature. The first is a one-dimensional periodic modulation of the phase-stiffness parameter. The second example is a ”granular” superconductor in which the Josephson coupling between grains is randomly distributed about its mean value. Both of these models predict a contribution to $`\sigma _1`$ above zero frequency which is not present if the phase stiffness is homogeneous throughout the material. Moreover, the spectral weight in this new contribution to $`\sigma _1`$ is proportional to the mean value of the phase stiffness, or superfluid density, just as is observed in our experiments. In both cases the proportionality constant is $`(\mathrm{\Delta }\rho _s)^2/\overline{\rho _s}^2`$, the fractional mean square variation in $`\rho _s`$.
Motivated by these models, we attempt a decomposition of $`\sigma _1`$ into QP and collective mode (CM) contributions, $`\sigma _1\sigma _{qp}+\sigma _{cm}`$. The difference between $`\sigma _1`$ and $`\sigma _{qp}`$ shows the collective mode contribution to be a low frequency peak contained, predominately, within our frequency range. We therefore test this decomposition with $`\sigma _{cm}`$ chosen to be a Lorenztian centered at $`\nu =0`$. This introduces two collective mode parameters: width, $`\mathrm{\Gamma }_{cm}(T)`$, and spectral weight, $`\rho _{cm}(T)`$. As is suggested by Fig. 2 we set the spectral weight of the collective mode, $`\rho _{cm}(T)`$, equal to a fixed fraction, $`\kappa `$, of $`\rho _s(T)`$. For the normal fluid spectral weight we choose $`\rho _n=\alpha T`$ as before, and we vary $`\mathrm{\Gamma }_{cm}`$ and $`\tau _{qp}`$ in order to find the best fit to the data. The best fit to $`1/\tau _{qp}(T)`$ is shown in the upper right panel of Fig. 3. (The dashed line indicates the range where the large thermal fluctuation conductivity dominates the QP contribution, making it difficult to extract $`\tau _{qp}`$). The width of the collective mode is found to be $`T`$ independent within the noise level, with a value 0.24 THz, and $`\kappa =0.30`$.
The left panel of Fig. 3 compares $`\sigma _1(T)\sigma _{qp}(T)`$ with $`\sigma _{cm}(T)`$, obtained using the parameter values given above. At each frequency $`\sigma _1(T)\sigma _{qp}(T)`$ is plotted as symbols and $`\sigma _{cm}(T)`$ as a line. The difference between the total conductivity and that due to QPs is obviously well described by the collective mode. What deviation there is between the model and the data is shown in the bottom right panel of Fig. 3. We see that difference corresponds to the phase fluctuations near $`T_c`$.
At this point we comment on the collective mode conductivity spectrum. In both models described previously, the frequency of the collective mode is related to the screened plasma frequency of the condensate, or $`\mathrm{}\omega _p=\sqrt{\rho _se^2/ϵd}`$. With $`ϵ10`$, this yields $`\mathrm{}\omega _p1400`$ K, much higher than the energies probed in our experiment. However, the plasma mode can be strongly perturbed by thermal QPs. For example, the damping rate of the plasmon, $`\mathrm{\Gamma }`$, is given by $`\mathrm{}\mathrm{\Gamma }=\rho _s\sigma _Q/\sigma _{qp}`$. For the values of $`\sigma _{qp}`$ seen in our experiment, this is of order 100 K, indicating that the plasma mode is highly overdamped. The corresponding conductivity will be centered at $`\nu =0`$ and its width will depend on the damping and screening effects of the QP background.
One final observation concerning our data is that causality dictates that the collective mode feature in $`\sigma _1(\nu )`$ must also be clearly manifest in $`\sigma _2(\nu )`$. Since both $`\sigma _1`$ and $`\sigma _2`$ are measured directly and independently by our coherent spectroscopy this is an important check of our model’s description of the conductivity. We find the total $`\sigma _2(\nu )`$ to be a sum of contributions from the QPs, the condensate and the collective mode. $`\sigma _2(\nu )`$, therefore, is not described by a two-fluid model but rather by the model set forth above.
In conclusion, we have shown that the dissipative conductivity of BSCCO in the superconducting state is not due solely to a normal fluid of QPs. There is an additional contribution whose spectral weight increases with decreasing temperature. We describe this contribution as a Lorentzian whose width is $`T`$ independent and whose spectral weight is a constant fraction of $`\rho _s(T)`$. This additional dissipation could result from phase fluctuations in the presence of static or quasistatic spatial variations of the local superfluid density. From the spectral weight of this dissipation we infer a fractional mean square variation in $`\rho _s`$ of $``$ 0.30. Adding this contribution to that of QPs with $`1/\tau _{qp}T`$ successfully describes the conductivity over the entire experimental range of frequency. Thus the transport lifetime in the superconducting state increases with decreasing $`T`$ far more slowly in BSCCO than in YBCO. The difference in lifetime may arise from the presence of same inhomogeneities which generate the collective mode dissipation.
We thank D.H. Lee, A. Maeda and D. Stroud for helpful discussions. This work was supported by NSF Grant No. 9870258, DOE Contract No. DE-AC03-76SF00098, and ONR Contract No. N00014-94-C-001.
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# 1 Introduction
## 1 Introduction
The string theory and the noncommutative geometry are fundamental theories possessing common goals: the elimination of the problems appearing in the standard field theory, like ultra-violet divergences, and, perspectively, a quantization of the gravity. In the former theory the elementary objects are strings, i.e. they are not point-like, whereas from the latter the notion of a point, as the elementary geometrical or physical entity, is eliminated from the very beginning (for detailed discussion see books , and ). However, it appears that there are deep relations between both approaches: the Yang-Mills theory on a noncommutative tori can be re-interpreted as a limiting model of the $`M`$-theory .
In the idea is discussed that the noncommutative space-time can appear naturally in the particular low energy limit of the string theory, as a direct consequence of a (constant) non-vanishing $`B`$-field. This fact is closely related to deformation quantization . Conversely, the deformation quantization can be interpreted in terms of a topological string theory . Here the noncommutative geometry appears in the space-time (the target space of string coordinates).
In this article we shall analyze the different link between the string theory and the noncommutative geometry relating the known deformation of the Virasoro algebra with the noncommutative geometry on a string world-sheet.
The Virasoro algebra is an infinite Lie algebra with generators $`L_m`$, $`m𝐙`$, satisfying commutation relations
$$[L_m,L_m^{}]=(mm^{})L_{m+m^{}}+\frac{c}{12}(m^3m)\delta _{m+m^{},0},$$
(1)
where $`c`$ is a central element commuting with all $`L_m`$. In an unitary irreducible representation $`L_m^+=L_m`$ and $`c`$ is a real constant. The Virasoro algebra is usually realized in terms of an infinite set of (bosonic or fermionic) oscillators, and is closely related to the symmetry properties of the 2D (conformal) field theory in question. The simplest such representation, with $`c=1`$, is given by the Sugawara construction of $`L_m`$ in terms of an infinite set of bosonic oscillators $`a_n`$, $`a_n^+=a_n`$, $`n=1,2,\mathrm{}`$, satisfying the commutation relations
$$[a_n,a_m]=n\delta _{n+m,0},n,m0.$$
(2)
The Sugawara formula for $`L_m`$ reads
$$L_m=\frac{1}{2}\underset{k,n0}{}:a_ka_n:\delta _{k+n,m},m𝐙.$$
(3)
The deformations of the Virasoro algebra are related to deformations of the Sugawara construction (3). The oscillator realizations (with a nontrivial central element) of these deformations of Virasoro algebra were constructed in . They have been intensively studied in -, mainly in connection with a formal deformations in the conformal and/or string field theories. Later a developed mathematical structure was found, the Zamolodchikov-Faddeev algebras, which proved to be a natural framework for the deformed Virasoro algebras . Our construction in represents a paricular realization of a bosonization of a ZF algebra. For a recent review see .
Our strategy is as follows. First we describe a model on a standard (commutative) cylinder, in which framework the Virasoro algebra appears naturally. Then we generalize the model to the noncommutative analog of the cylinder. This leads to the discrete time evolution. In Sec. 2 we summarize briefly free scalar field theory on a commutative and noncommutative cylinder. In Sec. 3 we describe the Euclidean version of the model and analyze its symmetry properties. We show that there is a particular model on a noncommutative cylinder possessing, as a symmetry algebra, the deformed Virasoro algebra proposed earlier. Thus, the Virasoro algebra on a complex plane $`𝐂`$ is replaced by the deformed Virasoro algebra on a $`q`$-deformed complex plane $`𝐂_q`$: there is a direct link between the noncommutative geometry on a world-sheet and the deformation of the Virasoro algebra. This link allows us to ascribe a clear physical meaning to the second index of the deformed Virasoro generators which necessarily appears in all abstract generalizations of the Virasoro algebra. It turns out that appearance of the additional index is a straightforward consequence of the noncommutativity of coordinates of the space-time ($`q`$-complex plane or, equivalently, noncommutative cylinder) on which the underlying field theory is defined. The field theoretical origin of the deformed Virosoro algebra can serve for the better (physical) motivation and understanding of its role as well as understanding of all related constructions ($`q`$-strings, $`q`$-vertex operators and Zamolodchikov-Faddeev algebras). Now all physical intuition (based on a notions like field action, quantization, Fock space, etc) can be used for further extensions, new constructions, etc originated and related to the q-deformed Virasoro algebras. In the last Sec. 4 we formulate the problem of a supersymmetric extension of the model and add concluding remarks.
## 2 Scalar field on a cylinder
### 2.1 Commutative case
In this section, we first discuss a free scalar field on a standard (commutative) cylinder $`C`$ (see, e.g., , and refs therein), and then we generalize the model to the noncommutative case. More detailed description of field-theoretical models on a noncommutative cylinder can be found in -.
The cylinder $`C`$ which we identify with the set of points $`C=𝐑\times S^1=\{x=(\rho \mathrm{cos}\varphi ,\rho \mathrm{sin}\varphi ,\tau )𝐑^3,\rho =\text{const}\}`$. If the function $`f(x)`$ on $`C`$ can be expanded as
$$f(\tau ,\phi )=\underset{k𝐙}{}c_k(\tau )e^{ik\phi },$$
(4)
we introduce the standard integral on $`C`$ by putting
$$I_0[f]=\frac{1}{2\pi }_C𝑑\tau 𝑑\phi f(\tau ,\phi )=_𝐑𝑑\tau c_0(\tau ).$$
(5)
The field action for a free massless real scalar field on a space-time modeled as the cylinder is defined by
$$S[\mathrm{\Phi }]=\frac{1}{2}I_0[\mathrm{\Phi }_\tau ^2\mathrm{\Phi }\mathrm{\Phi }_\phi ^2\mathrm{\Phi }].$$
(6)
It can be interpreted as the action describing (one coordinate of the) free closed bosonic string. We shall have in mind this interpretation in what follows. We can expand the fields $`\mathrm{\Phi }`$ into Fourier modes
$$\mathrm{\Phi }(x)=\underset{k0}{}c_k(\tau )e^{ik\phi },c_k^{}(\tau )=c_k(\tau ),$$
(7)
(we eliminated the zero mode which is inessential for us). Inserting (7) into the action we obtain
$$S[\mathrm{\Phi }]=\frac{1}{2}_𝐑𝑑\tau \underset{k0}{}[c_k(\tau )\ddot{c}_k(\tau )k^2c_k(\tau )c_k(\tau )].$$
(8)
The canonically conjugated momentum to the mode $`c_k(\tau )`$ is $`\pi _k(\tau )=\dot{c}_k(\tau )`$. We assume standard equal-time canonical Poisson brackets between modes and their conjugate momenta
$$\{c_k(\tau ),c_k^{}(\tau )\}=\{\pi _k(\tau ),\pi _k^{}(\tau )\}=0,$$
$$\{c_k(\tau ),\pi _k^{}(\tau )\}=\delta _{kk^{}}.$$
(9)
The quantization means an operator realization of the corresponding equal-time canonical commutation relations
$$[c_k(\tau ),c_k^{}(\tau )]=[\pi _k(\tau ),\pi _k^{}(\tau )]=0,$$
$$[c_k(\tau ),\pi _k^{}(\tau )]=i\delta _{kk^{}}.$$
(10)
This can be performed by solving corresponding classical Euler-Lagrange equations of motion
$$\ddot{c}_k(\tau )=k^2c_k(\tau ).$$
(11)
They possess positive frequency oscillating solutions
$$c_k(\tau )=\frac{i}{k}[a_ke^{ik\tau }b_ke^{ik\tau }],k>0.$$
(12)
The solutions for negative $`k`$ are fixed by the reality condition: $`c_k(\tau )=c^{}c_k(\tau )`$, i.e. $`a_k=a_k^{}`$, $`b_k^{}=b_k`$. Explicitly, the formulas for the field $`\mathrm{\Phi }(\tau ,\phi )`$ and the conjugated momentum $`\mathrm{\Pi }(\tau ,\phi )=_\tau \mathrm{\Phi }(\tau ,\phi )`$ read
$$\mathrm{\Phi }(\tau ,\phi )=\underset{k0}{}\frac{i}{k}[a_ke^{ik\tau +ik\phi }b_ke^{ik\tau ik\phi }],$$
$$_\tau \mathrm{\Phi }(\tau ,\phi )=\underset{k0}{}[a_ke^{ik\tau +ik\phi }+b_ke^{ik\tau ik\phi }].$$
(13)
The terms with expansion coefficients $`a_k`$ are interpreted as the right-movers on a closed bosonic string, whereas those solutions with $`b_k`$ as the left-movers. They are independent, and we can treat them separately. The canonical commutations relations are indeed satisfied if we replace the complex coefficients $`a_k`$ and $`b_k`$, $`k0`$, by two independent infinite set of bosonic oscillators satisfying commutation relations
$$[a_k,b_k^{}]=0,[a_k,a_k^{}]=[b_k,b_k^{}]=k\delta _{k+k^{},0}$$
(14)
(we are using the same notation for the classical expansion parameters and annihilation and creation operators, this should not lead to any confusion).
### 2.2 Noncommutative case
The noncommutative cylinder we realize by ”quantizing” the Poisson structure on $`C`$, see -. This is defined by
$$\{f,g\}=\frac{f}{\phi }\frac{g}{\tau }\frac{f}{\tau }\frac{g}{\phi },$$
(15)
for any pair of functions $`f`$ and $`g`$ on $`C`$. Using the Leibniz rule, it can be generated from the elementary brackets
$$\{\tau ,x_\pm \}=\pm ix_\pm ,\{x_+,x_{}\}=0,$$
(16)
where, $`x_\pm =\rho e^{\pm i\varphi }`$. Eqs. (17) are just $`e(2)`$ Lie algebra defining relations. The function $`x_+x_{}`$ is central: $`\{x_0,x_+x_{}\}=\{x_\pm ,x_+x_{}\}=0`$, i.e. the restriction $`x_+x_{}=\rho ^2`$ is consistent with the Poisson bracket structure. The operators $`_\phi ^2`$ and $`_\tau ^2`$ entering the action can be expressed in terms of Poisson brackets:
$$_\phi ^2f=\{\tau ,\{\tau ,f\}\},_\tau ^2f=\frac{1}{\rho ^2}\{x_{},\{\{x_+,f\}\}.$$
(17)
In the noncommutative case we replace the commuting variables $`\tau ,x_\pm `$ by the $`e(2)`$ Lie algebra generators $`\widehat{\tau },\widehat{x}_\pm `$ satisfying the relations
$$[\widehat{\tau },\widehat{x}_\pm ]=\pm \lambda \widehat{x}_\pm ,[\widehat{x}_+,\widehat{x}_{}]=0,$$
$$\widehat{x}_+\widehat{x}_{}=\rho ^2.$$
(18)
The algebra $`e(2)`$ possesses one series of infinite dimensional unitary representations (parameterized by one real parameter $`\rho >0`$). These representations can be realized in the Hilbert space $`L^2(S^1,d\phi )`$ as follows
$$\widehat{\tau }=i\lambda _\phi ,\widehat{x}_\pm =\rho e^{\pm i\phi }.$$
(19)
The Casimir operator takes the value $`\widehat{x}_+\widehat{x}_{}=\rho ^2`$. In what follows we put $`\rho =1`$.
For any operator $`\widehat{f}`$ on $`C`$ possessing the expansion
$$\widehat{f}=\underset{k𝐙}{}c_k(\widehat{\tau })e^{ik\phi },$$
(20)
we introduce the noncommutative analog of integral (5) by
$$I_\lambda [\widehat{f}]=\lambda \mathrm{Tr}[\widehat{f}]=\lambda \underset{n𝐙}{}c_0(n\lambda ).$$
(21)
Here $`c_0(n\lambda )`$ are spectral coefficients of the operator $`c_0(\widehat{\tau })`$: $`c_0(\widehat{\tau })e^{in\phi }=c_0(n\lambda )`$ $`e^{in\phi }`$ ($`e^{in\phi }`$ is an eigenfunction of $`\widehat{\tau }`$ with the eigenvalue $`n\lambda `$). There is a straightforward generalization of (22) to an integral over finite discrete time interval $`\alpha =a\lambda \tau b\lambda =\beta `$:
$$I_{\lambda \alpha }^\beta [\widehat{f}]:=\lambda \mathrm{Tr}_\alpha ^\beta [\widehat{f}]=\lambda \underset{a}{\overset{b}{}}c_0(n\lambda ).$$
(22)
Integrals of this type appear, e.g., if one calculates the field action for a finite time interval.
The operators $`\widehat{}_\phi ^2`$ and $`\widehat{}_\tau ^2`$, the noncommutative analogs of (18), are obtained replacing the Poisson brackets by the commutators: $`\{.,.\}(i/\lambda )[.,.]`$. Performing the Fourier expansion (21) we obtain
$$\widehat{}_\phi ^2\widehat{f}=\frac{1}{\lambda ^2}[\widehat{\tau },[\widehat{\tau },\widehat{f}]]=\underset{k𝐙}{}k^2c_k(\widehat{\tau })e^{ik\phi },$$
$$\widehat{}_\tau ^2\widehat{f}=\frac{1}{\lambda ^2\rho ^2}[\widehat{x}_{},[\widehat{x}_+,\widehat{f}]]=\underset{k𝐙}{}\delta _\lambda ^2c_k(\widehat{\tau })e^{ik\phi }.$$
(23)
Here $`\delta _\lambda ^2`$ is the second order difference operator
$$\delta _\lambda ^2c_k(\widehat{\tau })=\frac{1}{\lambda ^2}[c_k(\widehat{\tau }+\lambda )2c_k(\widehat{\tau })+c_k(\widehat{\tau }\lambda )].$$
The free hermitian scalar field action we take in the form
$$S[\widehat{\mathrm{\Phi }}]=\frac{1}{2}I_\lambda [\mathrm{\Phi }_\tau ^2\widehat{\mathrm{\Phi }}+K_\lambda ^2(i_\phi )\widehat{\mathrm{\Phi }}],$$
(24)
where $`K_\lambda (i_\phi )=\frac{2}{\lambda }\mathrm{sin}(i\frac{\lambda }{2}_\phi )`$. This specific form of the operator $`K_\lambda `$ has been chosen for the later convenience: as we shall see in section 3.2, the Euclidean version of (24) corresponds precisely to the action which can be reinterpreted as a theory on the noncommutative $`q`$-plane with a most simple and natural Lagrangian. The operator $`K_\lambda (i_\phi )`$ is defined by the Fourier expansion (21): $`K_\lambda (i_\phi )\widehat{f}=c_k(\widehat{\tau })K_\lambda (k)e^{ik\phi }`$$`K_\lambda (k)=\frac{2}{\lambda }\mathrm{sin}(\frac{k\lambda }{2})`$.
The field $`\widehat{\mathrm{\Phi }}`$ is a function in the noncommutative variables $`\widehat{\tau }`$ and $`\widehat{x}_\pm `$. It possesses, by assumption, the Fourier expansion
$$\widehat{\mathrm{\Phi }}=\underset{k0}{}c_k(\widehat{\tau })e^{ik\phi },c_k^{}(\widehat{\tau })=c_k(\widehat{\tau }k\lambda ).$$
(25)
The latter relation guarantees the hermiticity of the field:
$$\widehat{\mathrm{\Phi }}^{}=\underset{k0}{}e^{ik\phi }c_k^{}(\widehat{\tau })=\underset{k0}{}c_k^{}(\widehat{\tau }+k\lambda )e^{ik\phi }$$
$$=\underset{k0}{}c_k^{}(\widehat{\tau }k\lambda )e^{ik\phi }=\underset{k0}{}c_k(\widehat{\tau })e^{ik\phi }=\widehat{\mathrm{\Phi }}.$$
Inserting the mode expansion (25) into $`S[\widehat{\mathrm{\Phi }}]`$ and using the relation
$$c_k(\widehat{\tau })e^{ik\phi }c_k^{}(\widehat{\tau })e^{ik^{}\phi }=c_k(\widehat{\tau })c_k^{}(\widehat{\tau }k\lambda )e^{i(k+k^{})\phi },$$
(26)
we obtain
$$S[\widehat{\mathrm{\Phi }}]=\frac{\lambda }{2}\underset{n,k}{}[c_k(n\lambda k\lambda )\delta _\lambda ^2c_k(n\lambda )K_\lambda ^2(k)c_k(n\lambda k\lambda )c_k(n\lambda )].$$
(27)
Its extremalization leads to the discrete-time Euler-Lagrange equations
$$\delta _\lambda ^2c_k(n\lambda )=K_\lambda ^2(k)c_k(n\lambda ).$$
(28)
We see that the noncommutativity demonstrates itself dominantly as a discreteness of time: the action (27) does not contain explicitely noncommutative quantities, it depends on spectral modes at various discrete time slices. Moreover, in the equations of motion the time derivatives are replaced by the time differences.
The momentum, conjugated to the field mode $`c_k(n\lambda )`$, is
$$\pi _k(n\lambda )=\frac{1}{2\lambda }[c_k(n\lambda k\lambda +\lambda )c_k(n\lambda k\lambda \lambda ).$$
We postulate the standard equal-time Poisson brackets among modes and conjugated momenta
$$\{c_k(n\lambda ),c_k^{}(n\lambda )\}=\{\pi _k(n\lambda ),\pi _k^{}(n\lambda )\}=0,$$
$$\{c_k(n\lambda ),\pi _k^{}(n\lambda )\}=\delta _{kk^{}}.$$
(29)
The quantization means an operator realization of the corresponding equal-time canonical commutation relations
$$[c_k(n\lambda ),c_k^{}(n\lambda )]=[\pi _k(n\lambda ),\pi _k^{}(n\lambda )]=0,$$
$$[c_k(n\lambda ),\pi _k^{}(n\lambda )]=i\delta _{kk^{}},.$$
(30)
This can be performed similarly as in the commutative case, in terms of suitable sets of annihilation and creation operators. We shall not discuss this problem here since in the next Section, we analyze in detail the analogous problem within the Euclidean version.
## 3 Euclidean model on a cylinder
### 3.1 Commutative case
In order to analyze the Euclidean case (see, e.g., , ), we have to continue the time $`\tau `$ to $`i\tau `$, as a result in the action the kinetic term changes sign. The Euclidean action is usually defined by
$$S[\mathrm{\Phi }]=\frac{1}{2}I_0[\mathrm{\Phi }_\tau ^2\mathrm{\Phi }\mathrm{\Phi }_\phi ^2\mathrm{\Phi })].$$
(31)
The corresponding Euler-Lagrange equations for modes
$$\ddot{c}_k(\tau )=k^2c_k(\tau ).$$
(32)
can be solved similarly as in the real-time case. The formula for the field $`\mathrm{\Phi }(\tau ,\phi )`$ and the conjugate momentum $`\mathrm{\Pi }(\tau ,\phi )=_\tau \mathrm{\Phi }(\tau ,\phi )`$ reads
$$\mathrm{\Phi }(\tau ,\phi )=\underset{k0}{}\frac{i}{k}[a_ke^{k\tau +ik\phi }+b_ke^{k\tau ik\phi }],$$
$$_\tau \mathrm{\Phi }(\tau ,\phi )=i\underset{k0}{}[a_ke^{k\tau +ik\phi }+b_ke^{k\tau ik\phi }].$$
(33)
The equal-time canonical commutation relations among $`\mathrm{\Phi }(\tau ,\phi )`$ and $`\mathrm{\Pi }(\tau ,\phi ^{})`$ are satisfied provided that $`a_k`$ and $`b_k`$, $`k0`$, satisfy commutation relations (14).
It is useful to introduce the new complex variables
$$z=e^{\tau i\phi },\overline{z}=e^{\tau +i\phi }.$$
(34)
In this variables the action reads
$$S[\mathrm{\Phi }]=\frac{1}{4\pi }𝑑z𝑑\overline{z}\mathrm{\Phi }(z,\overline{z})\overline{}\mathrm{\Phi }(z,\overline{z}).$$
(35)
The corresponding Euler-Lagrange equations
$$\overline{}\mathrm{\Phi }(z,\overline{z})=0$$
(36)
have a general solution
$$\mathrm{\Phi }(z,\overline{z})=i\underset{k0}{}\frac{1}{k}[a_kz^k+b_k\overline{z}^k]=:\mathrm{\Phi }(z)+\overline{\mathrm{\Phi }}(\overline{z}).$$
(37)
Comparing with (33), we see that the fields $`\mathrm{\Phi }(z)`$ and $`\overline{\mathrm{\Phi }}(\overline{z})`$ correspond to the right- and left-movers, respectively, of a closed bosonic string.
The energy momentum tensor $`T_{ij}`$ is traceless: $`T_{z\overline{z}}=0`$ in complex notation with $`i,j=z,\overline{z}`$. This, together with the energy-momentum conservation,
$$\overline{}T_{zz}+T_{\overline{z}z}=0,T_{\overline{z}\overline{z}}+\overline{}T_{\overline{z}z}=0,$$
gives $`T(z,\overline{z})=T(z)+\overline{T}(\overline{z})`$ with
$$T(z)=\frac{1}{2}:\mathrm{\Phi }(z)\mathrm{\Phi }(z):,\overline{T}(\overline{z})=\frac{1}{2}:\overline{}\overline{\mathrm{\Phi }}(\overline{z})\overline{}\overline{\mathrm{\Phi }}(\overline{z}):.$$
(38)
Here there appears the normal product $`:\mathrm{}:`$ defined by
$$:\mathrm{\Phi }(z,\overline{z})\mathrm{\Phi }(w,\overline{w}):=\mathrm{\Phi }(z,\overline{z})\mathrm{\Phi }(w,\overline{w})+\mathrm{log}|zw|^2.$$
(39)
The last term is just the 2-point correlator (Green function) $`\mathrm{\Phi }(z,\overline{z})\mathrm{\Phi }(w,\overline{w})`$. It can be found, e.g. in the framework of Euclidean field theory, with the quantum expectation value of a field functional $`F[\mathrm{\Phi }]`$ being given as the path integral
$$F[\mathrm{\Phi }]=D\mathrm{\Phi }F[\mathrm{\Phi }]e^{S[\mathrm{\Phi }]}.$$
(40)
Let us now consider the infinitesimal transformation of fields $`\delta _\xi \mathrm{\Phi }=\xi \mathrm{\Phi }`$ with some given function $`\xi =\xi (z,\overline{z})`$. This induces the following variation of the action
$$\delta _\xi S=\frac{1}{4\pi }𝑑z𝑑\overline{z}[(\overline{}\xi )(\mathrm{\Phi })^2+(\xi \mathrm{\Phi }\overline{}\mathrm{\Phi })]$$
$$=\frac{1}{2\pi }𝑑z𝑑\overline{z}(\overline{}\xi )T,$$
(41)
with $`T(z)`$ given in (38) (here we used the formula $`𝑑z𝑑\overline{z}(\mathrm{})`$ $`=0`$). The analogous formula for $`\overline{T}(\overline{z})`$ can be obtained by considering the variations $`\delta _{\overline{\xi }}\mathrm{\Phi }=\overline{\xi }\overline{}\mathrm{\Phi }`$.
Inserting into $`T(z)`$ the mode expansion (33), using (38) and replacing the expansion coefficient by bosonic oscillators, we recover the Sugawara formula:
$$T(z)=L_mz^{m2},$$
(42)
where
$$L_m=\frac{1}{2\pi i}z^{n+1}T(z)=\frac{1}{2}:a_ka_n:\delta _{k+n,m}.$$
(43)
We note that here, on r.h.s., there appears the standard normal ordering with respect to the annihilation and creation operators. It can be shown that it is equivalent to the normal ordering introduced in (39), therefore, we have not introduced a new specific notation.
For variations $`\delta _\xi \mathrm{\Phi }=\xi \mathrm{\Phi }`$ with $`\xi =\xi (z)`$ it holds $`\delta _\xi S=0`$ (see (41)). Thus, they are symmetry transformations of the model. The variation $`\delta _\xi \mathrm{\Phi }`$ in question corresponds to an infinitesimal conformal transformation $`zz+\xi (z)`$ of the conformal plane:
$$\mathrm{\Phi }(z)\mathrm{\Phi }(z+\xi (z))=\mathrm{\Phi }(z)+\xi (z)\mathrm{\Phi }(z).$$
The infinitesimal conformal transformations of a plane form a Lie algebra, and the same is true for the variations: $`[\delta _\xi ,\delta _\xi ^{}]\mathrm{\Phi }=\delta _{\xi _\xi ^{}_\xi \xi ^{}}\mathrm{\Phi }`$. Consequently, the Virasoro generators $`L_m`$, $`m𝐙`$, close to a Lie algebra too. Wee see that the conformal mappings of a complex plane are behind the symmetry of the action in question.
### 3.2 Noncommutative case
Now we shall investigate along the similar lines the noncommutative Euclidean version of the model. It is obtained by replacing $`\lambda i\lambda `$ in (27). The corresponding Euclidean action is
$$S[\widehat{\mathrm{\Phi }}]=\frac{1}{2}I_\lambda [\widehat{\mathrm{\Phi }}\widehat{}_\tau ^2\widehat{\mathrm{\Phi }}\widehat{\mathrm{\Phi }}K_\lambda ^2(i_\phi )\widehat{\mathrm{\Phi }}],$$
(44)
where now $`K_\lambda (i_\phi )=\frac{2}{\lambda }\mathrm{sinh}(i\frac{\lambda }{2}_\phi )`$. Inserting into (44) the field mode expansion, we obtain
$$S[\widehat{\mathrm{\Phi }}]=\frac{\lambda }{2}\underset{n,k}{}[c_k(n\lambda k\lambda )\delta _\lambda ^2c_k(n\lambda )+K_\lambda ^2(k)c_k(n\lambda k\lambda )c_k(n\lambda )],$$
(45)
with $`K_\lambda (k)=\frac{2}{\lambda }\mathrm{sinh}(\frac{k\lambda }{2})`$. The corresponding discrete Euler-Lagrange equations
$$\delta _\lambda ^2c_k(n\lambda )=K_\lambda ^2(k)c_k(n\lambda ),k0,$$
(46)
have the solution
$$c_k(n\lambda )=\frac{i\lambda }{\mathrm{sinh}(k\lambda )}[a_ke^{\lambda k^2/2}e^{kn\lambda }b_ke^{\lambda k^2/2}e^{kn\lambda }],k0,$$
(47)
where $`a_k`$ and $`b_k`$ are independent constants. The conjugate momentum to the mode $`c_k(n\lambda )`$ is
$$\pi _k(n\lambda )=\frac{1}{2\lambda }[c_k(n\lambda k\lambda \lambda )c_k(n\lambda k\lambda +\lambda )]$$
$$=[a_ke^{\lambda k^2/2}e^{kn\lambda }+b_ke^{\lambda k^2/2}e^{kn\lambda }],k0.$$
(48)
The Euclidean reality condition, $`\widehat{\mathrm{\Phi }}^{}(\widehat{\tau },\phi )\stackrel{\mathrm{def}}{=}\widehat{\mathrm{\Phi }}^{}(\widehat{\tau },\phi )=\widehat{\mathrm{\Phi }}(\widehat{\tau },\phi )`$, requires $`c_k(n\lambda )=c_k^{}(n\lambda +k\lambda )`$. Taking this into account, the canonical commutation relations (30) are satisfied provided that the coefficients $`a_k`$ and $`b_k`$, $`k0`$, are replaced by bosonic operators satisfying the commutation relations
$$[a_k,b_k^{}]=0,[a_k,a_k^{}]=[b_k,b_k^{}]=\frac{\mathrm{sinh}(k\lambda )}{\lambda }\delta _{k+k^{},0}.$$
(49)
Inserting solution (47) into the field mode expansion, we obtain the field configuration, minimizing the action, in the form
$$\widehat{\mathrm{\Phi }}(\widehat{z},\widehat{\overline{z}})=\underset{k0}{}\frac{i\lambda }{\mathrm{sinh}(k\lambda )}[a_k\widehat{z}^kb_k\widehat{\overline{z}}^k]=:\widehat{\mathrm{\Phi }}(\widehat{z})+\widehat{\overline{\mathrm{\Phi }}}(\widehat{\overline{z}}).$$
(50)
Here we introduced the new noncommutative variables
$$\widehat{z}=e^{\widehat{\tau }i\phi }=e^{i\lambda _\phi i\phi },\widehat{\overline{z}}=e^{\widehat{\tau }+i\phi }=e^{i\lambda _\phi +i\phi }.$$
(51)
The operators $`\widehat{z}`$ and $`\widehat{\overline{z}}`$ satisfy the commutation relation
$$\widehat{z}\widehat{\overline{z}}=q^2\widehat{\overline{z}}\overline{z},q=e^\lambda .$$
(52)
which are typical for the $`q`$-deformed plane $`𝐂_q`$ possessing an $`E_q(2)`$ symmetry. Bellow we summarize necessary formulas following the conventions of .
The $`q`$-deformed Euclidean group $`E_q(2)`$ is a Hopf algebra generated by $`v`$, $`\overline{v}`$, $`n`$, $`\overline{n}`$ with relations
$$vn=q^2nv,v\overline{n}=q^2\overline{n}v,n\overline{n}=q^2\overline{n}n,$$
$$\overline{n}\overline{v}=q^2\overline{v}\overline{n},n\overline{v}=q^2\overline{v}n=v\overline{v}=\overline{v}v=1,$$
with comultiplication, counit and antipode given by
$$\mathrm{\Delta }v=vv,\mathrm{\Delta }\overline{v}=\overline{v}\overline{v},\mathrm{\Delta }n=n1+vn,\mathrm{\Delta }\overline{n}=\overline{n}1+\overline{v}\overline{n},$$
$$\epsilon (v)=\epsilon (\overline{v})=1,\epsilon (n)=\epsilon (\overline{n})=0,$$
$$S(v)=\overline{v},S(\overline{v})=v,S(n)=\overline{v}n,S(\overline{n})=v\overline{n}.$$
For a real $`q`$ the compatible involution is $`v^{}=\overline{v}`$, $`n^{}=\overline{n}`$.
The dual enveloping algebra $`𝒰_q(e(2))`$ is generated by the elements $`P_\pm `$ and $`J`$ such that
$$[P_+,P_{}]=0,[J,P_\pm ]=\pm P_\pm ,$$
with vanishing counit and
$$\mathrm{\Delta }J=J1+1J,\mathrm{\Delta }P_\pm =q^JP_\pm +P_\pm q^J,$$
$$S(J)=J,S(P_\pm )=q^{\pm 1}P_\pm ,J^{}=J,P_\pm ^{}=P_{}.$$
The plane $`𝐂_q`$ is associated with functions analytic in $`\widehat{z}`$, $`\widehat{\overline{z}}`$. It is invariant with respect to the $`𝒰_q(e(2))`$ coaction $`\mathrm{\Lambda }`$ given by (see ):
$$\mathrm{\Lambda }(q^{\pm J})\widehat{z}^m\widehat{\overline{z}}^n=q^{\pm mn}\widehat{z}^m\widehat{\overline{z}}^n,$$
$$\mathrm{\Lambda }(P_{})\widehat{z}^m\widehat{\overline{z}}^n_q\widehat{z}^m\widehat{\overline{z}}^n=[m]_qq^{n2}\widehat{z}^{m1}\widehat{\overline{z}}^n,$$
$$\mathrm{\Lambda }(P_+)\widehat{z}^m\widehat{\overline{z}}^n\overline{}_q\widehat{z}^m\widehat{\overline{z}}^n=[n]_qq^{m+1}\widehat{z}^m\widehat{\overline{z}}^{n1},$$
where $`[m]_q=(q^mq^m)/(qq^1)`$. The corresponding Casimir operator (the Laplacian on $`𝐂_q`$) is $`\mathrm{\Delta }_q=\mathrm{\Lambda }(P_+P_{})`$. Its action on $`𝐂_q`$ is
$$\mathrm{\Delta }_q\widehat{z}^m\widehat{\overline{z}}^n=\mathrm{\Lambda }(P_+P_{})\widehat{z}^m\widehat{\overline{z}}^n=q^{m+n2}[m]_q[n]_q\widehat{z}^{m1}\widehat{\overline{z}}^{n1}.$$
(53)
The operator $`\mathrm{\Delta }_\lambda =\widehat{}_{\overline{\tau }}^2+\frac{4}{\lambda }^2\mathrm{sinh}^2(i\frac{\lambda }{2}_\phi )`$ entering (44) acts on $`\widehat{z}^m\widehat{\overline{z}}^n=e^{(m+n)\widehat{\tau }i(mn)\phi mn\lambda }`$ as follows
$$\mathrm{\Delta }_\lambda \widehat{z}^m\widehat{\overline{z}}^n=[\widehat{}_\tau ^2+\frac{4}{\lambda }^2\mathrm{sinh}^2(i\frac{\lambda }{2}_\phi )]\widehat{z}^m\widehat{\overline{z}}^n$$
$$=\frac{(qq^1)^2}{\lambda ^2}[m]_q[n]_q\widehat{z}^m\widehat{\overline{z}}^n,q=e^\lambda .$$
(54)
Comparing (53) and (54) it can be shown straightforwardly that
$$\mathrm{\Delta }_\lambda \widehat{z}^m\widehat{\overline{z}}^n=\frac{(qq^1)^2}{\lambda ^2}\widehat{r}(\mathrm{\Delta }_q\widehat{z}^m\widehat{\overline{z}}^n)\widehat{r},\widehat{r}=e^{\widehat{\tau }}.$$
(55)
This is the noncommutative analog of the known link between Laplacian on a plane and those on a cylinder.
To any function $`f(\widehat{z},\widehat{\overline{z}})=_{n,m0}C_{m,n}\widehat{z}^m\widehat{\overline{z}}^n`$ on $`𝐂_q`$ we assign the Jackson-type integral by
$$_q𝑑\widehat{z}𝑑\widehat{\overline{z}}f(\widehat{z},\widehat{\overline{z}})=\lambda \mathrm{Tr}[\widehat{r}^2f_0(\widehat{r}^2)]=\lambda \underset{n𝐙}{}q^{2n}f_0(q^{2n}),$$
(56)
where $`\mathrm{Tr}`$ denotes the trace over the spectrum of $`\widehat{r}^2`$ is $`q^{2n}`$, $`n𝐙`$, and
$$f_0(\widehat{r})=\underset{n0}{}C_{n,n}q^{n^2}\widehat{r}^{2n}.$$
(57)
We can extend the definition (56) by taking a partial trace over the spectrum: putting $`\alpha =q^{2a}`$ and $`q^b=\beta `$, we define
$$_{q\alpha }^\beta 𝑑\widehat{z}𝑑\widehat{\overline{z}}f(\widehat{z},\widehat{\overline{z}})=\lambda \underset{n=a}{\overset{b}{}}q^{2n}f_0(q^{2n}).$$
(58)
Taking $`\widehat{f}=\widehat{r}^2f(\widehat{z},\widehat{\overline{z}})`$ it can be seen easily that the integrals $`I_\lambda [\widehat{f}]`$ and $`_q𝑑\widehat{z}𝑑\widehat{\overline{z}}f(\widehat{z},\widehat{\overline{z}})`$ are equal. The factor $`\widehat{r}^2`$ represents a Jacobian of the transformation (51). We see that in the noncommutative case there are two equivalent approaches, linked by the transformation (51):
(i) The first one corresponds to the model on a noncommutative sphere described by the noncommutative variables $`\widehat{\tau }`$ and $`\phi `$. The field action (44) leads to the equations of motion (46).
(ii) The second one is a model on a $`q`$-plane described by the variables $`\widehat{z}`$, $`\widehat{\overline{z}}`$. The action
$$S[\widehat{\mathrm{\Phi }}]=_q𝑑\widehat{z}𝑑\widehat{\overline{z}}\widehat{\mathrm{\Phi }}(\widehat{z}q,\widehat{\overline{z}}q^1)(\mathrm{\Delta }_q)\widehat{\mathrm{\Phi }}(\widehat{z},\widehat{\overline{z}}).$$
(59)
The shifts of arguments in the first $`\widehat{\mathrm{\Phi }}`$ guarantee that (59) is equivalent to (44) (use the link (53) between Laplacians and the relation $`\widehat{r}\widehat{\mathrm{\Phi }}(\widehat{z}q,\widehat{\overline{z}}q^1)=\widehat{\mathrm{\Phi }}(\widehat{z},\widehat{\overline{z}})\widehat{r}`$).
The action (59) depends on general hermitian field configurations
$$\widehat{\mathrm{\Phi }}(\widehat{z},\widehat{\overline{z}})=\underset{k0}{}\frac{i\lambda }{\mathrm{sinh}(k\lambda )}[a_k(q^k\widehat{r}^2)\widehat{z}^kb_k(q^k\widehat{r}^2)\widehat{\overline{z}}^k].$$
(60)
The field is hermitian provided that $`a_k^{}(q^k\widehat{r}^2)=a_k(q^1\widehat{r}^2)`$. The Euler-Lagrange equation for $`S[\widehat{\mathrm{\Phi }}]`$ is just the $`q`$-harmonicity condition
$$\mathrm{\Delta }_q\widehat{\mathrm{\Phi }}(\widehat{z},\widehat{\overline{z}})=_q\overline{}_q\widehat{\mathrm{\Phi }}(\widehat{z},\widehat{\overline{z}})=0.$$
(61)
Obviously, it is equivalent to the equation of motion (46).
Let us now investigate the variations of the action under infinitesimal field transformations
$$\delta _{\widehat{\xi }}\widehat{\mathrm{\Phi }}=\widehat{\xi }_q\widehat{\mathrm{\Phi }},\widehat{\xi }=\xi (\widehat{z},\widehat{\overline{z}}).$$
(62)
Using the Leibniz rule
$$_q[f(\widehat{z}q,\widehat{\overline{z}}q)g(\widehat{z},\widehat{\overline{z}})]=(_qf)(\widehat{z}q,\widehat{\overline{z}}q)g(\widehat{z}q,\widehat{\overline{z}}q)+f(\widehat{z}q,\widehat{\overline{z}}q)(_qg)(\widehat{z}q,\widehat{\overline{z}}q),$$
we can rewrite the action in the form
$$S[\widehat{\mathrm{\Phi }}]=q^2_q𝑑\widehat{z}𝑑\widehat{\overline{z}}(_q\widehat{\mathrm{\Phi }})(\widehat{z}q,\widehat{\overline{z}}q^1)(\overline{}_q\widehat{\mathrm{\Phi }})(\widehat{z},\widehat{\overline{z}}).$$
It follows straightforwardly,
$$\delta _{\widehat{\xi }}S[\widehat{\mathrm{\Phi }}]=q^2_qd\widehat{z}d\widehat{\overline{z}}[_q(\widehat{\xi }_q\widehat{\mathrm{\Phi }})(\widehat{z}q,\widehat{\overline{z}}q^1)(\overline{}_q\widehat{\mathrm{\Phi }})(\widehat{z},\widehat{\overline{z}}),$$
$$+(_q\widehat{\mathrm{\Phi }})(\widehat{z}q,\widehat{\overline{z}}q^1)\overline{}_q(\widehat{\xi }_q\widehat{\mathrm{\Phi }})(\widehat{z},\widehat{\overline{z}})].$$
(63)
We shall now restrict ourselves to the vicinity of solutions of the equation of motion, i.e. we take
$$\widehat{\mathrm{\Phi }}(\widehat{z},\widehat{\overline{z}})=\mathrm{\Phi }(\widehat{z})+\overline{\mathrm{\Phi }}(\widehat{\overline{z}}).$$
The first term in (63) does not contribute since, the integrand can be written as $`_q(\widehat{\xi }_q\mathrm{\Phi }\overline{}_q\overline{\mathrm{\Phi }})`$ and $`_q_q(\mathrm{})=0`$. Thus,
$$\delta _{\widehat{\xi }}S[\widehat{\mathrm{\Phi }}]=q^2_qd\widehat{z}d\widehat{\overline{z}}(_q\mathrm{\Phi })(\widehat{z}q,\widehat{\overline{z}}q^1)(\overline{}_q)\xi (\widehat{z},\widehat{\overline{z}})(_q\mathrm{\Phi }(\widehat{z},\widehat{\overline{z}}).$$
(64)
Any function $`\widehat{\xi }=\xi (\widehat{z},\widehat{\overline{z}})`$ can be written as a linear combination of functions $`\widehat{\xi }_k=\eta (\widehat{z})\widehat{\overline{z}}^k`$. It holds
$$\delta _{\widehat{\xi }_k}S[\widehat{\mathrm{\Phi }}]=_q𝑑\widehat{z}𝑑\widehat{\overline{z}}(_q\mathrm{\Phi })(\widehat{z}q)(\overline{}_q\widehat{\xi }_k)(\widehat{z},\widehat{\overline{z}})(_q\mathrm{\Phi })(\widehat{z})$$
$$=q^2_q𝑑\widehat{z}𝑑\widehat{\overline{z}}(\overline{}_q\widehat{\xi }_k)(\widehat{z},\widehat{\overline{z}})(_q\mathrm{\Phi })(\widehat{z}q^{2k+1})(_q\mathrm{\Phi })(\widehat{z}).$$
Shifting $`\widehat{z}\widehat{z}q^{k\frac{1}{2}}`$ we obtain for any $`k`$ the generator of field transformation in the splitted form
$$T_k(\widehat{z})=q^2:(_q\mathrm{\Phi })(\widehat{z}q^{k+\frac{1}{2}})(_q\mathrm{\Phi })(\widehat{z}q^{k\frac{1}{2}}).$$
(65)
This is exactly the formula for the splitted Virasoro generators, -. Inserting here the mode expansion (see (50))
$$(_q\mathrm{\Phi })(\widehat{z})=\frac{i\lambda q^2}{(qq^1)}\underset{k0}{}a_k\widehat{z}^{k1},$$
(66)
we obtain
$$T_k(\widehat{z})=\frac{\lambda ^2}{(q^21)^2}\underset{n𝐙}{}L_n^k\widehat{z}^{n2}.$$
(67)
Here
$$L_n^k=\underset{l,l^{}0}{}q^{(k1)(ll^{})}:a_la_l^{}:\delta _{l+l^{},n},$$
(68)
are generators of the (double indexed) deformed Virasoro algebra proposed in :
$$[L_n^k,L_n^{}^k^{}]=\frac{1}{4}\underset{\sigma \sigma ^{}}{}[\frac{nn^{}}{2}+n\sigma ^{}k^{}+n^{}\sigma k]_{}L_{n+n^{}}^{\sigma k^{}\sigma ^{}k}+C_n^{kk^{}}\delta _{n+n^{},0},$$
(69)
where
$$C_n^{kk^{}}=\frac{1}{2}\underset{m=1}{\overset{n}{}}[(n2m)k]_+[(n2m)k^{}]_+[m]_{}[nm]_{}.$$
(70)
Here we introduced the notation $`[x]_{}=(q^xq^x)/(qq^1)`$ and $`[x]_+=(q^x+q^x)/2`$.
Some comments are in order:
(i) In the commutative case the formula (42) for the transformation generators $`T(z)`$ follows for a general $`\xi (z,\overline{z})`$ and a general field configuration. In addition, in the commutative version, the model has the conformal symmetry.
(ii) In the noncommutative case the situation is different: we have different expressions for $`T_k(\widehat{z})`$ for integers $`k`$ (depending on the point splitting of the arguments in (65)) which generate transformations of $`q`$-harmonic functions among each other. They form the deformed Virasoro algebra introduced in (a particular realization of the Zamolodchikov-Faddeev algebra ). Notice that the appearance of the additional index $`k`$ in the deformed Virasoro algebra is directly related to the noncommutativity of the underlying two-dimensional space-time.
(iii) There are various reasons for which these symmetries of the model can not have as a background some ”conformal symmetry” of a $`q`$-plane: 1) transformations of the type $`zz+\xi (z)`$ do not spoil the $`𝐂_q`$ structure only for a very limited set of $`\xi (z)`$, and 2) the simple $`q`$-Taylor expansion formula $`\mathrm{\Phi }(z+\xi (z))=\mathrm{\Phi }(z)+\xi (z)_q\mathrm{\Phi }(z)`$ $`+\mathrm{}`$, is not valid even for an infinitesimal $`\xi (z)`$. The meaning of symmetry transformations, generated by the deformed Virasoro operators, requires further investigations. This problem is currently under study.
## 4 Concluding remarks
Recently have been found deep relations between string theory and noncommutative geometry in the space-time , . We have analyzed an alternative possibility introducing the noncommutative geometry on a string world-sheet. To achieve this aim we have investigated a free bosonic string on a noncommutative cylinder. Our results can be summarized as follows:
\- The field theory on a noncommutative cylinder leads in a natural way to the discrete time evolution -. We started with a suitable model for a free scalar field (the one component of bosonic string) on a noncommutative cylinder with a suitable particular symmetry of the field action.
\- The model in the Euclidean version can be equivalently formulated as a model on a $`q`$-deformed complex plane $`𝐂_q`$. Its symmetry is described by the deformed Virasoro algebra suggested earlier , appearing in the context of Zamolodchikov-Faddeev algebras .
The field theoretical origin of the deformed Virosoro algebra can serve for a better (physical) motivation and understanding of its role in all related constructions ($`q`$-strings, $`q`$-vertex operators and Zamolodchikov-Faddeev algebras). We see that the suggested formal deformation of the Virasoro algebra appears not only within the developed mathematical structure of Zamolodchikov-Faddeev algebras but also that there exists a physical theory, namely, the free bosonic string on a noncommutative cylinder, in which framework the deformed Virasoro algebra emerges as the symmetry of the theory. This fact allows one to justify the appearance and to clarify the meaning of the second index which labels (in addition to the usual one) the deformed Virasoro generators. As we have shown, this is related to the noncommutativity of the underlying two-dimensional space (world-sheet).
In this context it would be of great interest to extend our model to the supersymmetric case. We have strong indications that this can be achieved along the same lines as in the bosonic case:
(i) One can start from the fermionic realization of the deformed Virasoro algebra (69) proposed in . Introducing the Dirac operator on a noncommutative cylinder this can be interpreted in terms of a fermionic spinor field on a noncommutative cylinder.
(ii) Such spinor model can be reformulated as a theory on a noncommutative supercylinder. Consequently, the bosonic and fermionic realizations can be joined to a superfield theory on a noncommutative supercylinder. In the Euclidean version the resulting theory is formulated on a $`q`$-deformed superplane (with symmetries described by the quantum supergroup $`s`$-$`E_q(2))`$.
Both indicated steps require careful constructions of all noncommutative analogs of objects in question (the noncommutative (super)cylinder, Dirac operator, operator orderings, etc). Investigations in this direction are under current study.
Acknowledgments
The financial support of the Academy of Finland under the Projects No. 163394 is greatly acknowledged. A.D.’s work was partially supported also by RFBR-00-02-17679 grant and P.P.’s work by VEGA project 1/4305/97.
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# On the genealogy of a population of biparental individuals
### SUMMARY
If one goes backward in time, the number of ancestors of an individual doubles at each generation. This exponential growth very quickly exceeds the population size, when this size is finite. As a consequence, the ancestors of a given individual cannot be all different and most remote ancestors are repeated many times in any genealogical tree. The statistical properties of these repetitions in genealogical trees of individuals for a panmictic closed population of constant size $`N`$ can be calculated. We show that the distribution of the repetitions of ancestors reaches a stationary shape after a small number $`G_c\mathrm{log}N`$ of generations in the past, that only about $`80\%`$ of the ancestral population belongs to the tree (due to coalescence of branches), and that two trees for individuals in the same population become identical after $`G_c`$ generations have elapsed. Our analysis is easy to extend to the case of exponentially growing population.
## I Introduction
In the case of sexual reproduction, the ancestry of an individual is formed by 2 parents, 4 grandparents two generations ago, and in general $`2^G`$ individuals $`G`$ generations back into the past. The explosive growth of the number of ancestors belonging to the genealogical tree of, say, a present human should stop at some point due, at least, to the finite size of previous populations. For instance, only $`G33`$ generations ago (spanning a period of less than thousand years), the number of potential ancestors in the tree of any of us is about $`8.5\times 10^9`$, more than the present population of the Earth, and of course much larger than the population living about year 1000. The answer to this apparent paradox is simple: The branches of a typical genealogical tree often coalesce, indicating that many of the ancestors were in fact relatives and appear repeatedly in the tree (Ohno, 1996; Derrida et al., 1999; Gouyon, 1999). It might be difficult to test the statistical properties of such repetitions for an actual large, randomly mating population. Nevertheless, some exceptions can be found in royal genealogy. Since nobles usually married within their own castes, the presence of repeated ancestors in royal genealogical trees is far from rare. The example of the English king Edward III, where some ancestors appear up to six times, has been analysed in our previous work (Derrida et al., 1999).<sup>*</sup><sup>*</sup>*We used the tree of Edward III which can be found at http://uts.cc.utexas.edu/$``$churchh/edw3chrt.html.
Much attention has been paid in the past to a related problem, namely the statistical properties of branching processes (Harris, 1963) and its applications to the characteristics of the successive descendants of a single ancestor (Kingman, 1993). Actually, first applications of the branchig processes technique go back to the twenties. J.B.S. Haldane (Haldane, 1927) calculated the probability that a mutant allele be fixed in a population through a method developed previously by R.A. Fisher (Fisher, 1922). There, the relevant quantity was the survival probability of the descendants of the first individual carrying the mutation. All these studies apply to the vertical transmission of names, to the inheritance of characters coming only from one of the parents, like mithochondrial DNA or the Y chromosome, or to the fate of a mutant gene, for example, and correspond to an effective monoparental population. The heart of our problem is to take into consideration that reproduction is biparental. The distribution of repetitions of ancestors described below does however satisfy an equation similar to those which appear in branching processes (Harris, 1963).
Our problem of repetitions of ancestors in genealogical trees is much closer to the counting of the descendants of an individual in a sexual population. For example, in the case of a population of constant size, the average number of offspring is two per couple. Therefore after $`G`$ generations each individual has on average $`2^G`$ descendants. What prevents the number of descendants from growing exponentially with $`G`$ and to exceed the population size is interbreeding: When $`2^G`$ becomes comparable to the population size, interbreeding happens between the descendants and different lines of descent coalesce. The problem of the statistical properties of these coalescences is very similar to our present study of genealogical trees. None of them has –to our knowledge– yet been analysed.
In the present work, we study theoretically the problem of repetitions in the genealogical trees in the case of a closed, panmictic population. The study of the properties of a single tree with coalescent branches and the comparison of the genealogical trees of two contemporary individuals allows us to show that
1. There is a finite fraction (about 20% for a population of constant size $`N`$) of the initial population whose descendants becomes extinct after a number of generations $`G_c\mathrm{log}N`$. All the rest of the initial population (about 80%) belongs to all genealogical trees,
2. The distribution of the repetitions of ancestors living more than $`G`$ generations ago reaches a stationary shape after about $`G_c`$ generations,
3. The genealogical trees of two individuals in the same population become identical after a small number of generations $`G_c`$ back into the past,
4. The similarity between two genealogical trees changes from 1% (almost all ancestors in the two trees are different) to 99% (the repetitions of the ancestors in the two trees are almost identical) within 14 generations around $`G_c`$, independently of the population size $`N`$.
Our work can be generalized (see section IV) to describe coalescent processes, understood as the study of the gene tree originated when looking for the ancestry of a random sample of sequences (Kingman, 1982; Hudson, 1991; Donnelly & Tavaré, 1995). In the absence of recombination, each sequence has a single ancestor. The topology of thus reconstructed trees is equivalent to that generated through branching processes. Next in complexity, one can consider a two-loci sequence and assume that recombination can occur only between the two loci and with a small probability (meaning correlated genealogiesIn this paper, we use the term genealogy to refer to the ancestry of a single gene or of a whole set of sequences. In all cases, the genealogy is the complete set of ancestors contributing to the present object, this object being an individual (as in section II), a group of individuals (as in section III), a sequence (section IV), or a single locus (as quoted here). In this case, correlated genealogies simply means that the different sets of ancestors for the two-loci are not independent. for the two loci). The statistical properties of such process can be estimated until the most recent common ancestor (MRCA) is reached (Hudson, 1991). Instead, if one faces the study of a chromosome (Wiuf & Hein, 1997; Derrida & Jung-Muller, 1999) or of the whole genome, the number of ancestors grows as one proceeds back in time, since each individual has two parents and, apart from coalescence, also recombination (meaning splitting of the branches in the tree) is frequent.
If one considers a population or a sample of individuals within a population, there are relevant differences between the genealogy of a single gene and the genealogy of a chromosome or of the whole genome (which we study here). While in the first case, in fact, there exists a MRCA for the sample (where the gene tree ends), the genealogical tree of a chromosome or of the genome with two parents proceeds backwards in time and never reduces to a single ancestor. The genealogical tree representing the pedigree of a diploid organism contains a large fraction of the ancestral population. In this case, one may then talk about the most recent common set of ancestors, and study the similarities among different individuals now within the same population.
## II Statistical properties of an individual tree
Here we consider a simple neutral model of a closed population evolving under sexual reproduction and with non-overlapping generations.The Wright-Fisher model for allele frequencies works in the same set of hypothesis (Wright, 1931; Fisher, 1930). More recently, Serva and Peliti (Serva & Peliti, 1991) obtained a number of statistical results for the genetic distance between individuals in a sexual population evolving in the absence of natural selection. If the population size is $`N(G)`$ at generation $`G`$ in the past, we form couples at random (by randomly choosing $`N(G)/2`$ pairs of individuals) and assign each couple a random number $`k`$ of descendants. The probability $`p_k`$ of the number $`k`$ of offspring is given and if the population size is $`N`$ at present, its size $`N(G)`$ at generation $`G`$ in the past is given by
$$N(G)=\left(\frac{2}{m}\right)^GN$$
(1)
where the factor $`m`$ is obtained from
$$m=\underset{k}{}kp_k.$$
(2)
For $`m=2`$, the population size remains constant in time, whereas for $`m2`$ the number of individuals in the next generation is multiplied by a factor $`m/2`$. After a number of generations, the tree of each of the individuals in the youngest generation is reconstructed. To quantify the contribution of each of the ancestors to the genealogical tree of an individual, we define the weight $`w_\gamma ^{(\alpha )}(G)`$ of an ancestor $`\gamma `$ in the tree of individual $`\alpha `$ at generation $`G`$ in the past as
$$w_\gamma ^{(\alpha )}(G+1)=\frac{1}{2}\underset{\gamma ^{}\text{children of}\gamma }{}w_\gamma ^{}^{(\alpha )}(G)$$
(3)
We take $`w_\gamma ^{(\alpha )}(0)=\delta _{\alpha ,\gamma }`$, as this ensures that at generation $`G=0`$ all the weight is carried by the individual $`\alpha `$ itself. The factor $`1/2`$ in (3) keeps the sum of the weights normalized $`_{\gamma =1}^{N(G)}w_\gamma ^{(\alpha )}(G)=1`$, for any past generation $`G`$. The weight $`w_\gamma ^{(\alpha )}(G)`$ can be thought of as the probability of reaching ancestor $`\gamma `$ if one climbs up the reconstructed genealogical tree of individual $`\alpha `$ by choosing at each generation one of the two parents at random. The weights essentially measure the repetitions (see figure 1) in the genealogical tree. Without repetitions, $`w_\gamma ^{(\alpha )}(G)`$ would simply be $`2^G`$ for each ancestor $`\gamma `$ in the tree.
As an illustration of the previous quantities, we represent in Fig. 1 the result of random matings inside a small closed population of constant size $`N=14`$ (thus $`m=2`$) during 7 generations. The lines link progenitors with their offspring. The grey scale gives the weight $`w_\gamma (G)`$ of each of the individuals in the tree. The numbers on the left, all of them of the form $`r/2^G`$, give the weight of the leftmost individual in each generation. The denominators simply indicate the potential maximum number of ancestors at each generation. As counted by the numerator, each of them would appear repeated $`r`$ times in this tree if all the branches were explicitly shown.
We further assume that the probability $`p_k`$ of having $`k`$ children per couple follows a Poisson distribution, $`p_k=m^ke^m/k!`$ (most of what follows could be easily extended to other choices of $`p_k`$). We represent in Fig. 2 the probability for an English couple to have $`k`$ marrying sons during the period 1350-1986 (Dewdney, 1986). The solid line corresponds to a Poisson distribution with average $`1.15`$ (i.e., the average number of offspring per individual in that period, which corresponds to $`m=2.3`$ in our analysis), and implies that the total population is growing. These data spanning six centuries and taken over an homogeneous population support the hypothesis that the number of offspring is indeed Poisson distributed.<sup>§</sup><sup>§</sup>§Nonetheless, deviations from this distribution induced by a social transmission of the reproductive behaviour have been reported (Austerlitz & Heyer, 1998).
If we define $`S^{(\alpha )}(G)`$, the fraction of the population (at a generation $`G`$ in the past) which does not belong to the genealogical tree of individual $`\alpha `$, (i.e. such that $`w_\gamma ^{(\alpha )}(G)=0`$) one can show (see the appendix) that
$$S^{(\alpha )}(G+1)=\mathrm{exp}\left[m+mS^{(\alpha )}(G)\right].$$
(4)
This recursion, together with the initial condition $`S^{(\alpha )}(0)=11/N`$, determines this quantity for any $`G`$ (Derrida et al., 1999).
For large $`G`$ and for any individual $`\alpha `$, this fraction $`S^{(\alpha )}(G)`$ converges to the fixed point $`S(\mathrm{})`$ of (4). This gives for $`m=2`$ (i.e. for a population of constant size) a fraction $`S(\mathrm{})0.2031878..`$ which becomes extinct, so that the remaining fraction $`1S(\mathrm{})80\%`$ of the population belongs to the genealogical tree of any individual $`\alpha `$. A similar calculation shows that this 80% of the population which is not extinct after a large number of generations appears in the genealogical trees of all individuals: If $`S^{(\alpha ,\beta )}(G)`$ is the fraction of the population which does not belong to any of the two trees of two distinct individuals $`\alpha `$ and $`\beta `$, $`S^{(\alpha ,\beta )}(G)`$ satisfies the same recursion (4) as $`S^{(\alpha )}(G)`$, and converges to the same fixed value $`S(\mathrm{})`$. Thus, within this neutral model, an individual either becomes extinct (with a probability of 20%) or becomes an ancestor of the whole population after a large number of generations (with a probability of 80%). For an exponentially growing population with $`m=2.3`$ as in figure 2, the results are the same except for the precise value of $`S(\mathrm{})`$ (for $`m=2.3`$, one finds $`S(\mathrm{})14\%`$).
When $`G`$ is large enough, as shown in the appendix, the whole distribution $`P(w)`$ of the weights $`w_\gamma ^{(\alpha )}(G)`$ reaches a stationary shape, the properties of which can be calculated (Derrida et al., 1999). We show in Fig. 3 the distribution $`P(w/w)`$ for different values of $`m`$. As can be seen, it has a power-law dependence, $`P(w)w^\xi `$ for small values of the ratio $`w/w`$, with an exponent given by
$$\xi =\frac{\mathrm{log}S(\mathrm{})}{\mathrm{log}m}2,$$
(5)
and achieves a maximum value for $`w/w1`$.
## III Similarity between two trees
We would like to know how similar are the genealogical trees of two contemporary individuals and how they evolve in time within the same population. We have seen that a large fraction $`1S(\mathrm{})80\%`$ of the ancestral population constitutes the pedigree of every present individual. As a next step, one can compare two individuals and compute the degree of similarity between their trees, that is, the set of ancestors appearing at each generation in both trees simultaneously. We will see in particular that the two trees become identical after a number $`G_c`$ of generations.
We start with the definition of the overlap between the genealogical trees of two different individuals, $`\alpha `$ and $`\beta `$. Let $`w_\gamma ^{(\alpha )}(G)`$ be the weight of the ancestor $`\gamma `$ in the tree of $`\alpha `$ at generation $`G`$ in the past, and similarly let $`w_\gamma ^{(\beta )}(G)`$ be the weight of the same ancestor $`\gamma `$ at generation $`G`$ for $`\beta `$. These weights evolve according to (3) with $`w_\gamma ^{(\alpha )}(0)=\delta _{\gamma ,\alpha }`$ and $`w_\gamma ^{(\beta )}(0)=\delta _{\gamma ,\beta }`$ at generation $`G=0`$. In order to quantify the similarity between the two trees, we introduce the quantities
$$X^{(\alpha )}(G)=\underset{\gamma =1}{\overset{N(G)}{}}\left[w_\gamma ^{(\alpha )}(G)\right]^2$$
and
$$Y^{(\alpha ,\beta )}(G)=\underset{\gamma =1}{\overset{N(G)}{}}w_\gamma ^{(\alpha )}(G)w_\gamma ^{(\beta )}(G).$$
$`Y^{(\alpha ,\beta )}(G)`$ measures the correlation between the two trees at generation $`G`$ in the past and $`X^{(\alpha )}(G)`$ acts as a normalization factor. We then define the overlap $`q^{(\alpha ,\beta )}(G)`$ between the two trees at that generation by
$$q^{(\alpha ,\beta )}(G)=\frac{Y^{(\alpha ,\beta )}(G)}{\left[X^{(\alpha )}(G)X^{(\beta )}(G)\right]^{1/2}}$$
This overlap is a measure of the (cosine of the) angle between the two $`N`$dimensional vectors $`w_\gamma ^{(\alpha )}(G)`$ and $`w_\gamma ^{(\beta )}(G)`$.Similar quantities have been proposed as an indicator of the amount of evolutionary divergence between populations (Kimura, 1983). The quantity analogous to our weight $`w_\gamma ^{(\alpha )}`$ in the population genetics approach is the frequency of the sampled alleles, the number of ancestors $`\gamma `$ corresponds to the number of genes (that is the dimension of the space in which the vector $`w_\gamma ^{(\alpha )}`$ is embedded), and our individuals $`\alpha `$ and $`\beta `$ correspond to the compared populations (Cavalli-Sforza & Conterio, 1960). When $`q^{(\alpha ,\beta )}(G)0`$, the two vectors are essentially orthogonal and the ancestors of $`\alpha `$ and $`\beta `$ are all different. On the other hand, when $`q^{(\alpha ,\beta )}(G)1`$, the vectors are almost identical (as for brothers).
For a large enough population, the fluctuations of $`X^{(\alpha )}(G)`$ and $`Y^{(\alpha ,\beta )}(G)`$ are small around the population averaged values $`X(G)`$ and $`Y(G)`$ for almost all choices of $`\alpha `$ and $`\beta `$. Of course, if $`\alpha `$ and $`\beta `$ are brothers, $`Y^{(\alpha ,\beta )}(G)=X^{(\alpha )}(G)`$, a value very different from its average $`Y(G)`$; it is however very unlikely to get brothers, sisters or even cousins if one picks up two individuals at random from a large population.
The averages $`X(G)`$ and $`Y(G)`$ can be calculated from the evolution of the weights (3). Initially, $`X(0)=1`$ and $`Y(0)=0`$ since the individuals $`\alpha `$ and $`\beta `$ in any pair are different. Using the fact that for large $`N`$ the fluctuations of $`X^{(\alpha )}(G)`$ and $`Y^{(\alpha ,\beta )}(G)`$ are small, the expected value of the overlap $`q(G)`$ between two randomly chosen individuals is given by
$$q(G)\frac{Y(G)}{X(G)}=\frac{1}{1+m^{G_cG}}$$
(6)
where
$$G_c=\frac{\mathrm{log}((m1)N)}{\mathrm{log}m}1.$$
(7)
This expression is derived in the appendix. Of course Eq (6) is only valid with probability one with respect to the random choice of $`\alpha `$ and $`\beta `$ and with respect to the dynamics. We see that for large $`N`$, the overlap $`q(G)`$ is essentially zero for a number of generations of order $`G_c\mathrm{log}N/\mathrm{log}m`$ and then within a number of generations $`\mathrm{\Delta }G`$ which does not depend on $`N`$, it becomes equal to unity. Fig. 4 displays the averaged overlap $`q(G)`$ as a function of the number of generations $`G`$ for different values of $`N`$. We have chosen $`m=2`$ so that the population remains constant in size. We see that changing $`N`$ does not change the $`G`$ dependence except for a translation of the curve. In particular the range $`\mathrm{\Delta }G`$ on which the overlap changes from $`1\%`$ to $`99\%`$ does not depend on $`N`$. It is easy to check from (6) that for $`m=2`$, the overlap should satisfy
$$q(G+1)=\frac{2q(G)}{1+q(G)}$$
(8)
(plain line in the insert). The fixed point $`q(G)=0`$ is unstable for this map. All the trajectories finally converge to the stable fixed point $`q(G)=1`$ for large $`G`$. Also the quantity $`\mathrm{\Delta }G`$ can be estimated by counting how many generations are required for the overlap to change from 1 % to 99 % and this gives from (6)
$$\mathrm{\Delta }G\mathrm{log}(10^4)/\mathrm{log}m,$$
that is $`\mathrm{\Delta }G14`$ for $`m=2`$ and $`\mathrm{\Delta }G11`$ for $`m=2.3`$ as in figure 2. Typical values of $`G_c`$ are $`G_c20`$ for a population of constant size $`N=10^6`$. For a population increasing with $`m=2.3`$ as in figure 2, one gets $`G_c=21`$ if the size in the last generation is $`N=N(0)=75`$ millions.
The previous analysis can be easily extended to the hypothetical case of having an arbitrary number $`p`$ of parents instead of 2. As is shown in the appendix, the statistical properties of genealogical trees in a population of constant size but arbitrary $`p`$ are the same as for a population with only two parents and an expanding or shrinking size according to Eq. (1). The described statistical properties are thus equivalent in (i) a system with sexual reproduction and a growth rate $`m=p`$ and (ii) a system with constant population size but a number $`m`$ of genders.
The existence of a generation $`G_c`$ around which the genealogical similarity among individuals changes from 0 to 1 and which grows logarithmically with the size of the population is one of our main results. This has to be compared with the number of generations required for the population to become genetically homogeneous (Donnelly & Tavaré, 1991; Harpending et al., 1998), which grows proportionally to $`N`$. The difference is that when $`G_cGN`$, all the overlaps are 1, i.e all the genealogical trees in the population have the same ancestors with the same weights, but the genomes are still very different: This is just an extension of the situation of brothers who have exactly the same genealogical tree but different genomes.
## IV Simple model for the contribution of the ancestors to the genome
The evolution of a set of sequences subject to coalescence and recombination was first described by Hudson (1983). In this case, evolution proceeds until the most recent common ancestor for each set of homologous sites has been found. The set of MRCA sites does not necessarily belong to the genome of a single ancestor, on the contrary, it is in general spread on a finite fraction of the original population (Wiuf & Hein, 1997; 1999). In this section, we focus our attention on the statistical properties of the ancestry of a single extant genome. In particular, we calculate the equilibrium distribution for the fraction of material contributed by each ancestor.
Consider the whole set of genes that a present diploid organism has inherited from its parents. Although both parents contributed $`50\%`$ each, it is no longer true that grandparents contributed $`25\%`$ each, since independent assortment of chromosomes plus crossing over mixed in each of the parental gametes the material inherited from the previous generation. As a rough approximation to the output of genetic recombination, one might consider that each sequence is obtained as the addition of a fraction $`f`$ of the genetic material of one parent and a fraction $`1f`$ of the genetic material of the other parent with $`f(0,1)`$. This would be true if the length $`L`$ of the sequence was long enough (or infinitely long), so that there would be no restriction on the number of times it could be divided, and if one could forget the linear structure of the sequence. The process of coalescence and recombination (for small $`N`$) is schematically represented in Fig. 5.
We can now repeat the analysis done previously to the present extension. We will discard the correlations between the values of $`f`$ coming from a couple. This is equivalent to our assumption that fixing the pairs for $`k`$ offspring or choosing the parents of each individual at random only has effects of order $`O(N^1)`$ (see the appendix), and we can therefore work in the simplest realization of the process. Hence, we assume that the fraction $`f`$ takes independent values for each parent. The recursive equations (3) for the weights become
$$w_\gamma ^{(\alpha )}(G+1)=\underset{\gamma ^{}\text{children of}\gamma }{}f_\gamma ^{}w_\gamma ^{}^{(\alpha )}(G),$$
(9)
where the weight $`w_\gamma ^{(\alpha )}(G)`$ means now the fraction of the genetic material of individual $`\alpha `$ inherited from ancestor $`\gamma `$ at generation $`G`$. The random fraction $`f`$ is chosen anew for each offspring from a distribution $`\rho (f)`$ (with average value $`f=1/2`$). This implies that now even brothers would have different weights for their ancestors, and hence brings us slightly closer to the real genetic process.
Following the procedure described in the appendix, one can calculate the fraction $`S`$ of ancestors without lines of descent in the present (as we also show in Sec. II) and the exponent $`\xi `$ for the distribution $`P(w)`$. In general, given the distribution $`\rho (f)`$ for the contributions of the parents, we get
$$S(\mathrm{})=e^{mS(\mathrm{})m}$$
(10)
$$1=S(\mathrm{})m^{2+\xi }f^{1+\xi }f^{\xi 1}\rho (f)\text{d}f.$$
(11)
as one can easily show from (9) that the generating function $`h_G(\lambda )`$ defined by $`h_G(\lambda )=\mathrm{exp}[\lambda w(G)/w(G)]`$ has a limit $`h_{\mathrm{}}(\lambda )`$ for large $`G`$ which satisfies
$$h_{\mathrm{}}(\lambda )=\mathrm{exp}\left[m+m\rho (f)\text{d}fh\left(\frac{\lambda f}{mf}\right)\right]$$
Fig. 6 summarizes the changes in the distribution $`P(w)`$ for different distributions $`\rho (f)`$ of the random variable $`f`$. We have considered a simple case of a population of constant size (i.e. $`m=2`$) and with $`\rho (f)=1/(2\delta )`$ uniform in the interval $`(1/2\delta ,1/2+\delta )`$. In this particular case, an implicit relation between $`\delta `$ and the exponent $`\xi `$ can be obtained,
$$\delta \xi =S\left[\left(\frac{1}{2}\delta \right)^\xi \left(\frac{1}{2}+\delta \right)^\xi \right].$$
(12)
As $`\delta `$ varies, $`P(w)`$ remains a power law at small $`w`$ (i.e. $`P(w)w^\xi `$), and the exponent $`\xi `$ monotonously decreases with $`\delta `$. In particular, for $`\delta 0.35`$, $`\xi `$ changes sign: The maximum of $`P(w)`$ moves discontinuously from $`w/w1`$ to $`w/w0`$. The exponents obtained through simulations of the process are represented in Fig. 7 together with the numerical solution of Eq. (12), showing a good agreement.
## V Discussion
We have analysed the statistical properties of genealogical trees generated inside a closed sexual population. We focused our interest on the distribution of the repetitions of ancestors in the trees and on the amount of genetic material contributing to an extant genome. The precise values of $`\xi ,S(\mathrm{}),G_c`$ and $`\mathrm{\Delta }G`$ depend only weakly on the details of the model and do not change qualitatively if for instance a non Poissonian distribution of offspring is used. Moreover, we have shown how our results can be extended to the hypothetical case of having an arbitrary number $`p`$ of parents: Indeed, this case proves to be equivalent to a biparental population with a growth rate $`m/2=p/2`$.
The problem analysed here presents a number of connections to other fields. Equations similar to (3) appear also in the distribution of constraints in granular media where the variables $`w`$ represent the force acting on each grain and the recursion (3) expresses the way in which constraints are transmitted from one layer to the next (Coppersmith et al., 1996). In this case, $`p2`$ and even fluctuating $`p`$ would be perfectly realistic. The fact that the overlap changes from 0 to 1 within a small number of generations $`\mathrm{\Delta }G`$ independent of the size of the population and after $`G_c\mathrm{log}N`$ generations is also very reminiscent of the sharp cutoff phenomenon characteristic of some natural mixing processes modelled by Markov chains. One example of such systems is the shuffling of cards, where the stationary state in which the system has lost almost all information about the initial ordering of the $`n`$ cards is reached through a sharp cutoff after about $`\mathrm{log}n`$ riffle shuffles (Diaconis, 1996).
It is clear that the study of the interplay between the weights calculated in our generalized model and the structure of the genome would require more sophisticated approaches (Derrida & Jung-Muller, 1999; Wiuf & Hein, 1997; 1999). We have discarded the correlations between the history of neighboring sites in a sequence and assumed the independence of the factors $`f`$. Actually, the closer in the sequence two positions are, the more correlated are their genealogical histories (Kaplan & Hudson, 1985). This fact constrains the possible breaking points for our simulated sequences, implying that the random factors $`f`$ in (9) are a crude approximation to reality.
Since we have faced the problem from a statistical perspective, our results represent the average, typical behaviour, and are only valid with probability one when the population size is large. We did not study fluctuations due to the finite size of the population. Nonetheless, we hope that our results contribute to a better understanding of the role of genealogy in the degree of diversity of finite-size interbreeding populations.
## Appendix
In this appendix we have regrouped the technical aspects of the derivations of the main equations (4,5,6,10,11) presented in the body of the paper.
One may consider several variants of the model which all give a Poisson distribution for the number of offspring when the size of the population is large. For instance, the population size could be strictly multiplied by a factor $`m/2`$ at each generation or it could fluctuate (if we take the number of offspring from the Poisson distribution). One might decide that each individual has two parents chosen at random in the previous generation or form fixed couples and assign each couple some children. All these variants do not change the results when the population size is large, but might affect some finite size corrections that we compute in this appendix.
We will choose the following version of the model, which makes the calculation of the finite size corrections not too difficult. Our population has a given size $`N(G)`$ at each generation $`G`$ in the past, and we will assume that all the $`N(G)`$ are very large, at least in the range of generations $`G`$ that we will consider. Now, to construct the ancestors of all the $`N(G)`$ individuals at generation $`G`$ in the past, we choose for each of them a pair of parents at random among the $`N(G+1)`$ individuals at the previous generation (to facilitate the calculation, we do not even exclude that the two parents might be equal). Within this model, the number $`k`$ of children of an individual at generation $`G+1`$ is random and can be written as
$$k=\underset{i=1}{\overset{2N(G)}{}}z_i$$
where $`z_i=1`$ with probability $`1/N(G+1)`$ and $`z_i=0`$ otherwise. It follows that the whole distribution of $`k`$ can be calculated. The probability $`p_k`$ that an individual at generation $`(G+1)`$ has exactly $`k`$ children is given by the binomial distribution
$$p_k=\frac{(2N(G))!}{k!(2N(G)k)!}\left(\frac{1}{N(G+1)}\right)^k\left(1\frac{1}{N(G+1)}\right)^{2N(G)k}.$$
(13)
In particular,
$`k={\displaystyle \frac{2N(G)}{N(G+1)}}`$ (14)
$`k(k1)={\displaystyle \frac{2N(G)[2N(G)1]}{N(G+1)^2}}`$ (15)
$`k(k1)(k2)={\displaystyle \frac{2N(G)[2N(G)1][2N(G)2]}{N(G+1)^3}}.`$ (16)
If the population size is multiplied by a factor $`m/2`$ at each generation, i.e. if $`N(G)=N(G+1)m/2`$ (as $`G`$ counts the number of generations in the past), one recovers from (13) the Poisson distribution $`p_k=m^ke^m/k!`$ for large $`N(G)`$.
### A Calculation of the density of individuals without long term descendants and derivation of (4)
To establish (4), one simply needs to notice that for an individual to have no descendants after $`G+1`$ generations, all his children should have no descendants after $`G`$ generations. Let $`M(G)`$ be the number of individuals with no descendants at generation $`G`$ in the past. Given $`M(G)`$, one can write $`M(G+1)`$ as
$$M(G+1)=\underset{\gamma =1}{\overset{N(G+1)}{}}y_\gamma $$
where $`y_\gamma =1`$ if all the children of $`\gamma `$ are among the $`M(G)`$ and $`y_\gamma =0`$ otherwise. It can be shown that
$$y_\gamma =\left(1\frac{1}{N(G+1)}\right)^{2N(G)2M(G)}$$
and
$$y_\gamma y_\gamma ^{}=\left(1\frac{2}{N(G+1)}\right)^{2N(G)2M(G)}$$
for $`\gamma \gamma ^{}`$. This gives
$$M(G+1)=N(G+1)\left(1\frac{1}{N(G+1)}\right)^{2N(G)2M(G)}$$
(17)
$$M^2(G+1)=M(G+1)+N(G+1)[N(G+1)1]\left(1\frac{2}{N(G+1)}\right)^{2N(G)2M(G)}$$
(18)
When all the $`M`$’s and $`N`$’s are large, we see from (17,18) that the fluctuations of $`M(G+1)`$ are small (as $`M^2(G+1)M(G+1)^2M(G+1)^2`$), and one finds from (17) that the ratio $`M(G)/N(G)S^{(\alpha )}(G)`$ satisfies
$$S^{(\alpha )}(G+1)=\mathrm{exp}\left[\frac{2N(G)}{N(G+1)}(S^{(\alpha )}(G)1)\right]$$
which is identical to (4) for $`N(G)=N(G+1)m/2`$.
### B Time evolution of the distribution of the weights
From the recursion (3) and from the known distribution (13) of $`k`$ one can write recursions for the moments of the weights
$`w_\gamma ^{(\alpha )}(G+1)={\displaystyle \frac{k}{2}}w_\gamma ^{(\alpha )}(G)`$ (19)
$`[w_\gamma ^{(\alpha )}(G+1)]^2={\displaystyle \frac{k}{4}}[w_\gamma ^{(\alpha )}(G)]^2+{\displaystyle \frac{k(k1)}{4}}w_\gamma ^{(\alpha )}(G)w_\gamma ^{}^{(\alpha )}(G),`$ (20)
where $`\gamma \gamma ^{}`$. The normalization $`_\gamma w_\gamma ^{(\alpha )}=1`$ allows one to rewrite
$$w_\gamma ^{(\alpha )}(G)w_\gamma ^{}^{(\alpha )}(G)=\frac{1}{N(G)1}\left[w_\gamma ^{(\alpha )}(G)[w_\gamma ^{(\alpha )}(G)]^2\right]$$
and together with the known moments (16) gives that
$`w_\gamma ^{(\alpha )}(G+1)={\displaystyle \frac{N(G)}{N(G+1)}}w_\gamma ^{(\alpha )}(G)={\displaystyle \frac{1}{N(G+1)}}`$ (21)
$`[w_\gamma ^{(\alpha )}(G+1)]^2=\left[{\displaystyle \frac{N(G)}{2N(G+1)}}{\displaystyle \frac{N(G)[2N(G)1]}{2N(G+1)^2[N(G)1]}}\right][w_\gamma ^{(\alpha )}(G)]^2`$ (22)
$`+{\displaystyle \frac{2N(G)1}{2N(G+1)^2[N(G)1]}}`$ (23)
where $`\gamma \gamma ^{}`$.
For large $`N(G)`$, if the ratio $`N(G+1)/N(G)=2/m`$, as in the case of a population increasing by a factor $`m/2`$ at each new generation, expression (23) becomes simpler and one gets
$`[w_\gamma ^{(\alpha )}(G+1)]^2={\displaystyle \frac{m}{4}}[w_\gamma ^{(\alpha )}(G)]^2+{\displaystyle \frac{m^2}{4}}\left({\displaystyle \frac{1}{N(G)}}\right)^2`$ (24)
In this limit, we have from (16) that $`k=m`$ and $`k(k1)=m^2`$, and we see that (24) means that in (19) the weights $`w_\gamma ^{(\alpha )}`$ and $`w_\gamma ^{}^{(\alpha )}`$ are, for large $`N(G)`$, uncorrelated. The calculation of higher moments of the weights can be done in the same manner and for large $`N(G)`$ the weights of different ancestors become again uncorrelated.
If the population size changes in time, the distribution of the weights cannot be stationary. This is already visible in the expression (19) which shows that even the first moment of the weights changes with $`G`$. One can however check from (19) and (24) that the ratio $`[w_\gamma ^{(\alpha )}(G)]^2/w_\gamma ^{(\alpha )}(G)^2`$ which satisfies
$`{\displaystyle \frac{[w_\gamma ^{(\alpha )}(G+1)]^2}{w_\gamma ^{(\alpha )}(G+1)^2}}={\displaystyle \frac{1}{m}}{\displaystyle \frac{[w_\gamma ^{(\alpha )}(G)]^2}{w_\gamma ^{(\alpha )}(G)^2}}+1`$ (25)
has a limit $`m/(m1)`$ as $`G`$ increases. Moreover, as the initial value of this ratio is $`N(0)`$, the number of generations $`G_c`$ to converge to this limit is $`G_c\mathrm{log}N(0)/\mathrm{log}m`$. Higher moments of the weights behave in a similar way and one can write recursions for ratios which generalize (25) and which show that all the ratios have limits.
This indicates that the distribution of the ratio $`w/w`$ becomes stationary. In the limit of large $`N(G)`$ (considering that the weights of the different children $`\gamma ^{}`$ in (3) can be taken as independent and that the distribution of $`k`$ becomes Poissonian), one finds that the generating function $`h_G(\lambda )`$ defined by
$$h_G(\lambda )=\mathrm{exp}\left[\lambda \frac{w_\gamma ^{(\alpha )}(G)}{w(G)}\right]$$
(26)
satisfies
$$h_{G+1}(\lambda )=\underset{k}{}p_k\left[h_G\left(\frac{\lambda w(G)}{2w(G+1)}\right)\right]^k=\mathrm{exp}\left[m+mh_G(\lambda /m)\right].$$
(27)
Recursion (27) generalizes to the case $`m2`$ (i.e. the case of an exponentially increasing population) the result of our previous work obtained for a population of constant size ($`m=2`$). Similar recursions have been studied in the theory of branching processes (Harris, 1963). The use of generating functions in population genetics is well illustrated in the book by Gale (1990), where this method is for example applied to the calculation of the probability of fixation of a mutant allele.
It is remarkable, that if one considers an imaginary world where each individual would have $`p`$ parents (instead of $`2`$), the generating function (26), in the case of a population of constant size, would satisfy the recursion (27) with $`m=p`$. This means that as long as the distribution of weights is concerned, the problem of a large population of constant size with $`m`$ parents per individual is identical to the problem of a population of size increasing at each generation by a factor $`m/2`$ with two parents per individual.
### C Stationary distribution
For large $`G`$, if we fix the ratio $`N(G)/N(G+1)=m/2`$, the generating function $`h_G(\lambda )`$ converges to $`h_{\mathrm{}}(\lambda )`$ solution of
$$h_{\mathrm{}}(\lambda )=\mathrm{exp}\left[m+mh_{\mathrm{}}(\lambda /m)\right]$$
(28)
If one expands the solution around $`\lambda =0`$, one finds that
$$h_{\mathrm{}}(\lambda )=1+\lambda +\frac{1}{2}\frac{m}{m1}\lambda ^2+\frac{1}{6}\frac{m^2(m+2)}{(m^21)(m1)}\lambda ^3+\mathrm{}$$
and the comparison with (26) gives for large $`G`$
$$\frac{w^2}{w^2}\frac{m}{m1};\frac{w^3}{w^3}\frac{m^2(m+2)}{(m^21)(m1)};$$
which means that in principle the whole shape of $`P(w)`$ can be extracted from (28). In particular, one can predict the power law of $`P(w)`$ at small $`w`$. Trying to solve (28) for large negative $`\lambda `$, if one writes
$$h_{\mathrm{}}(\lambda )S(\mathrm{})\frac{1}{|\lambda |^{\xi +1}}$$
(29)
one finds, as expected, that $`S(\mathrm{})`$ is the fixed point of (4). Eq. (29) is equivalent to the asumption that $`P(w)w^\xi `$ at small $`w`$, where the exponent $`\xi `$ should satisfy
$$1=S(\mathrm{})m^{\xi +2}.$$
This completes the derivation of (5) which was already discussed in our previous work (Derrida et al., 1999).
### D Overlap between two trees
Let us now show how (6) can be derived. Starting from recursion (3), one obtains by averaging over all the links relating generation $`G`$ to generation $`G+1`$
$$w_\gamma ^{(\alpha )}(G+1)w_\gamma ^{(\beta )}(G+1)=\frac{k}{4}w_\gamma ^{(\alpha )}(G)w_\gamma ^{(\beta )}(G)+\frac{k(k1)}{4}w_\gamma ^{(\alpha )}(G)w_\gamma ^{}^{(\beta )}(G),$$
(30)
where $`\gamma \gamma ^{}`$ and the averages over $`k`$ are carried out with respect to (13). This gives
$$w_\gamma ^{(\alpha )}(G+1)w_\gamma ^{(\beta )}(G+1)=\frac{m}{4}w_\gamma ^{(\alpha )}(G)w_\gamma ^{(\beta )}(G)+\frac{1}{4}\left(m^2\frac{m}{N(G+1)}\right)w_\gamma ^{(\alpha )}(G)w_\gamma ^{}^{(\beta )}(G).$$
(31)
Using the fact that the sum $`_\gamma ^{}w_\gamma ^{}^{(\beta )}(G)=1`$, so that $`w_\gamma ^{(\alpha )}(G)=1/N(G)`$ at all generations, one gets that
$`w_\gamma ^{(\alpha )}(G+1)w_\gamma ^{(\beta )}(G+1)={\displaystyle \frac{m}{4}}w_\gamma ^{(\alpha )}(G)w_\gamma ^{(\beta )}(G)`$ (32)
$`+{\displaystyle \frac{1}{4}}\left(m^2{\displaystyle \frac{m}{N(G+1)}}\right){\displaystyle \frac{\frac{1}{N(G)}w_\gamma ^{(\alpha )}(G)w_\gamma ^{(\beta )}(G)}{N(G)1}}.`$ (33)
Keeping only the dominant contributions for large $`N`$’s we arrive at
$$w_\gamma ^{(\alpha )}(G+1)w_\gamma ^{(\beta )}(G+1)=\frac{m}{4}w_\gamma ^{(\alpha )}(G)w_\gamma ^{(\beta )}(G)+\frac{m^2}{4}\frac{1}{N(G)^2}.$$
Comparing this expression with (30), one sees that for large $`N`$, one could have simply neglected the correlations between the weights of different individuals, (i.e. directly replaced $`w_\gamma ^{(\alpha )}(G)w_\gamma ^{}^{(\beta )}(G)`$ by $`w_\gamma ^{(\alpha )}(G)w_\gamma ^{}^{(\beta )}(G)`$) and used the Poisson distribution instead of (13)). The previous recursion can be integrated
$$w_\gamma ^{(\alpha )}(G)w_\gamma ^{(\beta )}(G)=\left[w_\gamma ^{(\alpha )}(0)w_\gamma ^{(\beta )}(0)+\frac{1}{N^2}\frac{m}{m1}(m^G1)\right]\left(\frac{m}{4}\right)^G,$$
(34)
and using the fact that $`w_\gamma ^{(\alpha )}(G)w_\gamma ^{(\beta )}(G)`$ is equal to $`Y(G)/N(G)`$ when $`\alpha \beta `$ and to $`X(G)/N(G)`$ when $`\alpha =\beta `$, one finds (with $`X(0)=1`$ and $`Y(0)=0`$)
$$\frac{Y(G)}{X(G)}=\frac{(m^G1)m^{G_c}}{1+(m^G1)m^{G_c}}$$
where $`G_c`$ is given by (7). For large $`N`$, that is for large $`G_c`$ this reduces to (6) in the whole range where the expression departs from $`0`$ or $`1`$, that is for $`G`$ of order $`G_c`$. Finally, one can check that with the value of $`G_c`$ given by (7), $`N(G)`$ is always large, as long as $`N`$ is large, so that our assumption that all the $`N`$’s were large was legitimate.
ACKNOWLEDGEMENTS. The authors acknowledge discussions with Jordi Bascompte, Ugo Bastolla and Julio Rozas. SCM thanks the Alexander von Humboldt Foundation for support.
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# Long range correlations in the non-equilibrium quantum relaxation of a spin chain
\[
## Abstract
We consider the non-stationary quantum relaxation of the Ising spin chain in a transverse field of strength $`h`$. Starting from a homogeneously magnetized initial state the system approaches a stationary state by a process possessing quasi long range correlations in time and space, independent of the value of $`h`$. In particular the system exhibits aging (or lack of time translational invariance on intermediate time scales) although no indications of coarsening are present.
\]
Non-equilibrium dynamical properties of quantum systems have been of interest recently, experimentally and theoretically. Measurements on magnetic relaxation at low-temperatures show deviations from the classical exponential decay, which was explained by the effect of quantum tunneling. On the theoretical side, among others, integrable and non-integrable models were studied in the presence of energy or magnetic currents, as well as the phenomena of quantum aging in systems with long-range and short-range interactions.
Here we pose a different question: Consider a quantum mechanical interacting many body system described by a Hamilton operator $`\widehat{H}`$ without any coupling to an external bath, which means that the system is closed. Suppose the system is prepared in a specific state $`|\psi _0`$ at time $`t=0`$, which is not an eigenstate of the Hamiltonian $`\widehat{H}`$. Then we are interested in the natural quantum dynamical evolution of this state which is described by the Schrödinger equation and is formally given by
$$|\psi (t)=\mathrm{exp}\left(\frac{i}{\mathrm{}}\widehat{H}t\right)|\psi _0.$$
(1)
Obviously the energy $`E=\psi _0|\widehat{H}|\psi _0`$ is conserved. In particular we want to study the time evolution of the expectation value $`A(t)`$ of an observable $`\widehat{A}`$ or the two-time correlation function $`C_{AB}(t_1,t_2)`$ of two observables $`\widehat{A}`$ and $`\widehat{B}`$, defined by
$`A(t)`$ $`=`$ $`\psi _0|\widehat{A}_H(t)|\psi _0`$ (2)
$`C_{AB}(t_1,t_2)`$ $`=`$ $`\psi _0|\{\widehat{A}_H(t_1)\widehat{B}_H(t_2)\}_S|\psi _0,`$ (3)
where $`\widehat{A}_H(t)=\mathrm{exp}(+i\widehat{H}t)\widehat{A}\mathrm{exp}(i\widehat{H}t)`$ is the operator $`\widehat{A}`$ in the Heisenberg picture (with $`\mathrm{}`$ set to unity) and $`\{\widehat{A}\widehat{B}\}_S=1/2(\widehat{A}\widehat{B}+\widehat{B}\widehat{A})`$ the symmetric product of two operators.
One should emphasize that in such a situation one does not expect time translational invariance to hold, which would manifest itself in, for instance, $`A(t)=A_0=\mathrm{const}.`$ and $`C_{AB}(t_1,t_2)=C_{AB}(t_1t_2)`$. There will be a transient regime in which these relations are violated and depending on the complexity of the system this non-equilibrium regime will extend over the whole time axis, in which we would denote it as quantum aging, as it can be observed for instance for the universe, which is (most probably) a closed system.
To be concrete we consider the prototype of an interacting quantum systems, the Ising model in a transverse field (TIM) in one dimension defined by the Hamiltonian:
$$H=\frac{1}{2}\left[\underset{l=1}{\overset{L1}{}}\sigma _l^x\sigma _{l+1}^x+h\underset{l=1}{\overset{L}{}}\sigma _l^z\right],$$
(4)
where $`\sigma _l^{x,z}`$ are spin-$`1/2`$ operators on site $`l`$. We consider initial many body states that are eigenstates either of all local $`\sigma _l^x`$ or of all local $`\sigma _l^z`$ operators. We will mostly be concerned with fully magnetized initial states, either in the $`x`$ or the $`z`$ direction, which we denote with $`|x`$ and $`|z`$, respectively, and which obey $`\sigma _l^x|x=+|x`$ and $`\sigma _l^z|z=+|z`$, respectively.
In passing we note that one obtains the zero temperature equilibrium situation by choosing the ground state of the Hamiltonian (4) as the initial state. This ground state has a quantum phase transition at $`h=1`$ from a paramagnetic ($`h>1`$) to a ferromagnetic ($`h<1`$) phase, the latter being indicated by long range order in the $`x`$-component, i.e. a non-vanishing expectation value of $`\sigma ^x`$. Moreover, non-zero temperature equilibrium relaxation has been considered in .
The expectation values and correlation functions we are interested in are those that originate from these spin operators $`\sigma _l^x`$ and $`\sigma _l^z`$. In order to compute them, we have to express the Hamiltonian (4) in terms of fermion creation (annihilation) operators $`\eta _q^+`$ ($`\eta _q`$)
$$H=\underset{q}{}ϵ_q\left(\eta _q^+\eta _q\frac{1}{2}\right).$$
(5)
The energy of modes, $`ϵ_q`$, $`q=1,2,\mathrm{},L`$ are given by the solution of the following set of equations
$`ϵ_q\mathrm{\Psi }_q(l)`$ $`=`$ $`h\mathrm{\Phi }_q(l)\mathrm{\Phi }_q(l+1),`$ (6)
$`ϵ_q\mathrm{\Phi }_q(l)`$ $`=`$ $`\mathrm{\Psi }_q(l1)h\mathrm{\Psi }_q(l),`$ (7)
and we use free boundary conditions which implies for the components $`\mathrm{\Phi }_q(L+1)=\mathrm{\Psi }_q(0)=0`$. The spin-operators can then be expressed by the fermion operators as
$`\sigma _l^x`$ $`=`$ $`A_1B_1A_2B_2\mathrm{}A_{l1}B_{l1}A_l,`$ (8)
$`\sigma _l^z`$ $`=`$ $`A_lB_l,`$ (9)
with
$$\begin{array}{ccc}A_i\hfill & =& _{q=1}^L\mathrm{\Phi }_q(i)(\eta _q^++\eta _q),\hfill \\ B_i\hfill & =& _{q=1}^L\mathrm{\Psi }_q(i)(\eta _q^+\eta _q),\hfill \end{array}$$
(10)
and the time-evaluation of the spin operators follows from that of the fermion operators: $`\eta _q^+(t)=e^{itϵ_q}\eta _q^+`$ and $`\eta _q(t)=e^{itϵ_q}\eta _q`$.
To calculate different non-equilibrium correlation functions we have developed a systematic method in which the time-dependent contractions are defined by:
$`A_lA_k_t`$ $`=`$ $`{\displaystyle \underset{q}{}}\mathrm{cos}(ϵ_qt)\mathrm{\Phi }_q(l)\mathrm{\Phi }_q(k),`$ (11)
$`A_lB_k_t`$ $`=`$ $`B_kA_l_t=i{\displaystyle \underset{q}{}}\mathrm{sin}(ϵ_qt)\mathrm{\Phi }_q(l)\mathrm{\Psi }_q(k),`$ (12)
$`B_lB_k_t`$ $`=`$ $`{\displaystyle \underset{q}{}}\mathrm{cos}(ϵ_qt)\mathrm{\Psi }_q(l)\mathrm{\Psi }_q(k).`$ (13)
play a central role. For general position of the spin, $`l=O(L/2)`$, one finds simple formulas for the expectation values and correlation functions involving $`\sigma _l^z`$ operators, whereas the calculation of those involving $`\sigma _l^x`$ operators is a difficult task and the final result is complicated . However, both the surface-spin auto-correlations and the end-to-end correlations are given in quite simple form, both for the equilibrium and for the non-equilibrium case.
First we study the $`x`$-end-to-end correlations defined by
$$C_L^{x,\psi }(t)=\psi _0|\{\sigma _1^x(t)\sigma _L^x(t)\}_S|\psi _0,$$
(14)
which contain information about the existence or absence of magnetic order in the $`x`$-direction. The single time $`t`$ at which the correlations between the two spins are measured indicates the age of the system after preparation. For the fully ordered initial state $`|\psi _0=|x`$ we obtain
$$C_L^{x,x}(t)=A_1A_1_tB_LB_L_t+|A_1B_L_t|^2,$$
(15)
The first term in the r.h.s. of Eq.(15) is the product of surface magnetizations at the two ends of the chain. Therefore $`lim_{L,t\mathrm{}}C_L^{x,x}(t)=\overline{m_1}^2`$ and the stationary state, starting with $`|x`$, has long-range order for $`h<1`$ as $`\overline{m_1}=1h^2`$. Thus the surface order-parameter, $`\overline{m_1}`$, vanishes continuously at the transition point, $`h_c=1`$, with a non-equilibrium exponent, $`\beta _1^{ne}=1`$.
The time dependence of the connected correlations (generally defined via 2 as $`\stackrel{~}{C}_{AB}(t_1,t_2)=C_{AB}(t_1,t_2)A(t_1)B(t_2)`$) $`\stackrel{~}{C}_L^{x,x}(t)=|A_1B_L_t|^2`$ shows the following features which can be read from fig. 1: 1) They are zero for times smaller than a time $`\tau _h(L)`$ which is equal to the system size $`L`$ for $`h1`$ and increases monotonically with decreasing $`h`$ for $`h<1`$. 2) At $`t=\tau _h(L)`$ a jump occurs to a value that decreases algebraically with the system size $`L`$:
$$\stackrel{~}{C}_{\mathrm{max}}^{x,x}(L)=\stackrel{~}{C}_L^{x,x}(t=\tau _h(L))L^a,$$
(16)
with $`a=2/3`$ for $`h=1`$ and $`a=1/2`$ for $`h>1`$. 3) For $`t\tau _h(L)`$ the correlations decay slower than exponentially, roughly with a stretched exponential. 4) For $`t=3\tau _h(L)`$ again a sudden jump occurs as for $`t=\tau _h(L)`$ followed by a slightly slower oscillatory decay. 5) This pattern is repeated for time $`t=5\tau _h(L),7\tau _h(L),\mathrm{}`$, but gets progressively smeared out by oscillations.
These features can be interpreted as follows: the elementary (tunnel) processes of the quantum dynamics of the Hamiltonian (4) are spin flips induced by the transverse field operator $`\sigma _l^z`$. In this picture two spins can only act coherently and thus give a contribution to the connected correlation function if the information about such a spin flip processes reaches the two spins simultaneously. Feature 1 tells us that signals generated in the center of the system travel with a speed of proportional to $`L/\tau _h(L)`$ to the boundary spins and reaches both simultaneously. At this moment $`\stackrel{~}{C}_L^{x,x}(t)`$ jumps to its maximum (see 2). After this, this signal is superposed by other more incoherent signals (see 3). However, the strongest initial signal is reflected at both boundaries and reaches the opposite boundary spins simultaneously again at time $`t=3\tau _h(L)`$ (see 4), and so on. Of course more and more incoherent processes occur in the meantime, giving rise to feature 5.
A similar behavior can be observed for the end-to-end correlations when starting with the state $`|z`$, which is
$$C_L^{x,z}(t)=\underset{k}{}\left(A_1B_k_tB_LA_k_tA_1A_k_tB_LB_k_t\right).$$
(17)
The only difference to the behavior of $`C_L^{x,x}(t)`$ reported above is a) its long time limit vanishes for all values of $`h`$ and b) $`\tau _h(L)`$, i.e. the earliest time at which the two boundary spins are correlated, is only half as big as in the previous case. Obviously it is easier to generate and to propagate spin flip signals when starting with a $`z`$-state.
Next we study the bulk behavior of the expectation values and correlations involving $`\sigma _l^z`$ operators. We start its non-equilibrium expectation value
$`e_l^\psi (t)`$ $`=`$ $`\psi _0|\sigma _l^z(t)|\psi _0`$ (18)
$`=`$ $`{\displaystyle \underset{k}{}}\left(A_lB_k_tB_lA_{i(k)}_tA_lA_{i(k)}_tB_lB_k_t\right),`$ (19)
with $`i(k)=k,(k+1)`$ for $`\psi =z,(x)`$. We note that the equilibrium (i.e. ground state) expectation value, $`e_l^0`$, corresponds to the energy-density in the two-dimensional classical Ising model and we use this terminology also in this non-equilibrium situation. For long times the non-equilibrium energy-density approaches a finite limit, $`\overline{e}_l^\psi `$, which for a bulk spin is a) for the initial state $`|\psi _0=|x`$ given by $`\overline{e^x}=h/2`$ for $`h1`$ and $`\overline{e^x}=1/(2h)`$ for $`h>1`$. and b) for the initial state $`|\psi _0=|z`$ by $`\overline{e^z}=1/2`$ for $`h1`$ and $`\overline{e^z}=11/2h^2`$ for $`h>1`$. Therefore the analogue to the specific heat $`c_v\overline{e^z}/h`$ is discontinuous at the transition point. The relaxation of the energy-density to its stationary value is algebraic and follows a $`t^{3/2}`$ low for any value of the transverse field, $`h`$. At the transition point, $`h=1`$, we have the analytical results in terms of the Bessel-function, $`J_\nu (x)`$: $`e_l^\psi (t)=1/2\pm J_1(4t)/4t`$, where the + (–) sign refer to $`\psi =z(x)`$.
The two-spin non-equilibrium dynamical and spatial correlations involve contributions from different processes described by the contractions (13) and the corresponding formulas are complicated, therefore they will be presented elsewhere. Here we report on the basic features of the asymptotic behavior of correlations. The two-time correlation function ($`t_1t_2`$),
$$G_l^{z,\psi }(t_1,t_2)=\psi _0|\{\sigma _l^z(t_1)\sigma _l^z(t_2)\}_S|\psi _0,$$
(20)
is non-stationary for intermediate times, $`t_1/(t_2t_1)=O(1)`$, which can be read off from our analytical result for the connected bulk auto-correlations at $`h=1`$:
$$\stackrel{~}{G}_l^{z,\psi }(t_1,t_2)=J_0^2(2t_22t_1)\frac{1}{4}[f(t_2+t_1)\pm g(t_2t_1)]$$
(21)
where $`f(x)=J_2(2x)+J_0(2x)`$, $`g(x)=J_2(2x)J_0(2x)`$ and the + (–) sign refer to $`\psi =x(z)`$. Thus we conclude that for intermediate times there is aging in the $`z`$-component auto-correlation function, contrary to what is reported in . Asymptotically we have $`lim_{t_1\mathrm{}}G_l^{z,\psi }(t_1,t_2)=(\overline{e}^\psi )^2`$, and the connected two-time correlations depends only on the time difference, e.g. for $`h=1`$ via (21) $`lim_{t_1\mathrm{}}\stackrel{~}{G}_l^{z,\psi }(t_1,t_2)=J_0^2(2[t_2t_1])\{J_1^{}(2[t_2t_1])\}^2`$. For bulk spins this stationary correlation function decays algebraically as $`(t_2t_1)^2`$ for any value of $`h`$.
Next we consider the spatial equal-time correlations
$$C^{z,\psi }(r,t)=\psi _0|\{\sigma _{ir/2}^z(t)\sigma _{i+r/2}^z(t)\}_S|\psi _0,$$
(22)
where $`i=L/2`$ in a finite system. For long times they approach the stationary limit, $`lim_t\mathrm{}C^{z,\psi }(r,t)=(\overline{e}^\psi )^2`$, the same as for the auto-correlation function. For the connected correlation function $`\stackrel{~}{C}^{z,\psi }(r,t)=C^{z,\psi }(r,t)e_l^\psi (t)e_{l+r}^\psi (t)`$ we can derive an analytic expression at the transition point, $`h=1`$ in the limit $`L\mathrm{}`$
$`\stackrel{~}{C}^{z,\psi }(r,t)`$ $`=`$ $`\left[{\displaystyle \frac{r}{2t}}J_{2r}(4t)\right]^2`$ (23)
$``$ $`{\displaystyle \frac{r^21}{4t^2}}J_{2r+1}(4t)J_{2r1}(4t),`$ (24)
which is valid both for $`|\psi _0=|x`$ and $`|\psi _0=|z`$. In fig. 2 we show the $`r`$-dependence of $`\stackrel{~}{C}^{z,z}(r,t)`$ for various times $`t`$. We see that for fixed time $`t`$ the correlations increase proportional to $`r^2`$ for distances $`rt`$ to a maximum value $`\stackrel{~}{C}_{\mathrm{max}}^{z,z}(t)`$ at $`r=2t`$, which decreases with time proportional to $`t^1`$. For distances larger than $`r=2t`$ they drop rapidly, faster than exponentially, to zero.
The latter two features correspond perfectly to what we observed also for the $`z`$-end-to-end correlations, see eq(17): spins that are separated by a distance $`r`$ can only be correlated after the first signal from spin flip processes in between them reach simultaneously the two spins, i.e. for times $`t`$ larger than $`r/2`$ (for $`h=1`$ and $`|\psi _0=|z`$. The first feature, that correlations for distances smaller than $`2t`$ are diminished only algebraically rather than via a stretched exponential in the case of end-to-end correlations, is new and characteristic for bulk spins. For $`r2t`$ the correlation function $`\stackrel{~}{C}^{z,\psi }(r,t)`$ obeys the characteristic scaling form
$$\stackrel{~}{C}^{z,\psi }(r,t)=t^1g(r/t)$$
(25)
with $`g(x)x^2`$ for $`x1`$. The scaling parameter $`r/t`$ appearing in the scaling function $`g(x)`$ is reminiscent of the fact that space and time scales are connected linearly at the critical point in the transverse Ising chain since the dynamical exponent is $`z=1`$. Away from the critical point we have to evaluate our expressions for $`\stackrel{~}{C}^{z,\psi }(r,t)`$ numerically for finite but large ($`L=512`$) system sizes. Essentially we observe the same scenario as at the critical point, the only difference being that the general functional dependency is superposed by strong oscillations. Moreover, starting with $`|\psi _0=|x`$ instead of $`|z`$ changes the correlations only by a constant factor.
We collect now our results for the maximum value for connected spin-spin correlations since they decay algebraically with various new exponents. we define ourselves to $`h1`$ since here the time $`\tau _h`$ of maximum correlation is fixed, whereas for $`h<1`$ the value of $`\tau _h`$ depends on $`h`$ and has to be determined numerically that renders the precise determination of the decay exponents difficult. We define the ratio $`\alpha =t/L`$ and $`\alpha _{\mathrm{max}}=\tau _h(L)/L`$ and consider equal time correlations for fixed values of $`\alpha `$. In the picture of a propagating front, that separates a region in the space-time diagram for the chain in which spins are uncorrelated from a region in which they are correlated, one observes quasi long range correlations on the front, the latter being defined by the ratio $`t/L=\alpha _{\mathrm{max}}`$. For distances smaller than the distance of maximum correlation or times larger than $`\tau _h`$ the correlations decay slower than exponential in time, e.g. algebraically for bulk spins ($`\stackrel{~}{C}^{z\psi }(t,r=\mathrm{fixed})t^2`$)). When we vary both space and time with fixed ratio $`t/L`$ or $`t/r`$ we get power laws, as long as we stay behind the front (i.e. $`t\tau _h`$). For $`\alpha >\alpha _{\mathrm{max}}`$ we observe again power laws, but with different exponents; they are listed in table 1.
To conclude we studied a novel type of dynamically produced long range correlations in a quantum relaxation process in a quantum spin chain. Starting with a homogeneous initial state the quantum mechanical time evolution according to the Schrödinger equation drives the system into a stationary state, which has algebraically decaying time-dependent autocorrelations but no critical fluctuations. However, during the relaxation process spin-spin correlation build up upon arrival of a front of coherent signals, which afterwards decay algebraically in the bulk. On the front and behind it for fixed ratio of space and time scales one observes quasi long range order. This does not depend on any external parameter like the transverse field. This type of algebraic correlation needs not to be triggered by some tuning parameter and is therefore reminiscent of phenomena in self organized criticality . The scenario we have reported here is a result of quantum interference and one may expect that a similar one holds for other quantum systems, too. At this point one should mention the possibility of coarsening in quantum systems as for instance reported in , which is different from the scenario we have reported here.
Acknowledgment: This study has been partially performed during our visits in Saarbrücken and Budapest, respectively. F. I.’s work has been supported by the Hungarian National Research Fund under grant No OTKA TO23642, TO25139, MO28418 and by the Ministery of Education under grant No. FKFP 0596/1999.
| | | $`h=1`$ | | $`h>1`$ | |
| --- | --- | --- | --- | --- | --- |
| | $`\alpha _{\mathrm{max}}`$ | $`\alpha `$=$`\alpha _{\mathrm{max}}`$ | $`\alpha >\alpha _{\mathrm{max}}`$ | $`\alpha `$=$`\alpha _{\mathrm{max}}`$ | $`\alpha >\alpha _{\mathrm{max}}`$ |
| $`\stackrel{~}{C}_L^{xx}(t`$=$`\alpha L)`$ | $`1`$ | $`L^{2/3}`$ | $`L^1`$ | $`L^{1/2}`$ | $`L^1`$ |
| $`\stackrel{~}{C}_L^{xz}(t`$=$`\alpha L)`$ | $`1/2`$ | $`L^{1/4}`$ | $`L^{1/2}`$ | $``$ | $``$ |
| $`C_L^{z\psi }(t`$=$`\alpha L)`$ | $`1/2`$ | $`L^{5/4}`$ | $`L^1`$ | $`L^{5/8}`$ | $`L^1`$ |
| $`\stackrel{~}{C}^{z\psi }(t`$=$`\alpha r)`$ | $`1/2`$ | $`r^{4/3}`$ | $`r^1`$ | $`r^{2/3}`$ | $`r^1`$ |
TABLE I
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# Renormalization group study of the 1d polaron problem
\[
## Abstract
We investigate the one-dimensional polaron problem and calculate the ground state energy and the effective mass. We use a real-time renormalization group method and compare our results with first and second order perturbation theory, with Feynman’s variational principle and with the method of Lee, Low and Pines.
\]
Introduction. The polaron has often been examined since Fröhlich proposed the corresponding Hamiltonian . It serves as a standard model for various problems involving a nonrelativistic particle moving in a scalar field, e.g. the interaction between nucleons and scalar mesons or a single electron in a solid interacting with longitudinal optical phonons. The physical picture is that the particle polarizes the environment and must drag this polarization with it, which affects its energy and effective mass. To describe the polaron not only in the weak coupling limit one has to consider it beyond perturbation theory. The most successful methods for determining the ground state energy and effective mass of the polaron are those of Lee, Low and Pines , Pekar et al. , and Feynman . Since polaron effects have been observed in low-dimensional systems, the problem has also been studied in two dimensions .
In this paper we will examine the one-dimensional polaron problem. It can be realized e.g. for a Bloch-electron in a one-dimensional wire or macromolecular structure. The excitation of an electron will be strongly influenced by the interaction with optical phonons . If the conduction band is partially filled, one can linearize the electronic spectrum and the model is exactly solvable by using bosonization techniques . However, if the conduction band is empty, it is necessary to consider a quadratic spectrum for the electron with a bare mass $`m_0`$. This leads to the one-dimensional Fröhlich Hamiltonian with a constant coupling to the phonons. It is our purpose to examine the ground state energy and the effective mass of the electron for this problem. Thereby we will use a recently developed real-time renormalization group (RG) technique . While this method actually allows nonequilibrium descriptions , in the present paper we have only used it to determine spectral properties. Furthermore, we will compare our results with first and second order perturbation theory, with Feynman’s variational principle, and with the method of Lee, Low and Pines generalized to the one-dimensional case with finite band width. For not too large coupling constants we find a ground state energy near Feynman’s method and a value for the effective mass between the result of Feynman and the one of Lee, Low and Pines.
Model. The polaron problem is modelled by the Fröhlich Hamiltonian $`H=H_0+H_1`$ with
$`H_0`$ $`=`$ $`{\displaystyle \underset{k}{}}ϵ_kc_k^{}c_k+{\displaystyle \underset{q}{}}\mathrm{}\omega _qa_q^{}a_q`$ (1)
$`H_1`$ $`=`$ $`{\displaystyle \underset{k,q}{}}M_k^q(a_q^{}+a_q)c_{k+q}^{}c_k.`$ (2)
$`c_k^{}(c_k)`$ is the creation (annihilation) operator of the electron, its energy $`ϵ_k`$ is given by the electron dispersion $`(\mathrm{}k)^2/2m_0`$, $`m_0`$ being the bare mass of the electron in the conduction band and $`a_q^{}(a_q)`$ creates (annihilates) a phonon with frequency $`\omega _q`$. Since the interaction is dominated by the longitudinal optical branch, we presume dispersionless phonons, i.e. $`\omega _q=\omega `$. While the electron-phonon interaction coefficients $`M_k^qM^q`$ are proportional to $`1/q`$ in case of the bulk polaron, the one-dimensional situation involves a $`q`$-independent coefficient $`M`$ . We define
$$M=(\frac{4m_0\omega }{\mathrm{}})^{\frac{1}{4}}\frac{\sqrt{\alpha }}{\sqrt{L}},$$
where $`L`$ is the one-dimensional normalization volume. In analogy to the three-dimensional case $`\alpha `$ is a dimensionless coupling constant. In the following we choose units such that $`\mathrm{}=m_0=\omega =1`$. With these assumptions the perturbation theory produces the ground state energy $`E_g=\alpha `$ and the inverse effective mass $`1/m=1\alpha /2`$.
RG method. In the RG we consider the $`S`$-matrix
$`S`$ $`=`$ $`Te^{i{\scriptscriptstyle 𝑑tH(t)}}`$ (3)
$`=`$ $`e^{i{\scriptscriptstyle 𝑑tH_0}}Te^{i{\scriptscriptstyle 𝑑tH_1(t)_I}},`$ (4)
where $`H_1(t)_I`$ is the interaction part of the Hamiltonian taken in the interaction picture with respect to $`H_0`$. The idea of the RG is to leave the $`S`$-matrix invariant while successively integrating out diagrams of different time scales. The procedure is schematically shown in Fig. 1. For a given cutoff $`t_c`$ in time space, we allow only for correlation functions of the phonons with a time scale $`t>t_c`$. At zero temperature, the latter are given by
$`(a_q^{}+a_q)(t)(a_q^{}+a_q)=e^{i\omega t}.`$ (5)
All correlation functions with time scales shorter than $`t_c`$ are accounted for by renormalized energies and coupling constants. A change of $`t_c`$ to $`t_c+dt_c`$ is made by applying three steps : (i) expanding the second exponential in (4) and introducing normal ordering for the phonon operators using Wick’s theorem, (ii) integrating over the contractions with a time scale between $`t_c`$ and $`t_c+dt_c`$ and (iii) resumming the operators in an exponential form. Consequently these operators will not be limited to zero-phonon and one-phonon operators any more. But we shall see that a good approximation is achieved for not too large coupling constants if we neglect double or higher order vertex operators. One advantage of this method is that we only need to consider the operators of the electronic system, as the phonon degrees of freedom are integrated out in each $`t_c`$-step. For the polaron problem we start with $`t_c=0`$ involving the operators as in (1), (2) and end up with an effective Hamiltonian for $`t_c\mathrm{}`$. We obtain the following RG-equations
$`{\displaystyle \frac{dϵ_k}{dt_c}}`$ $`=`$ $`i{\displaystyle \underset{q}{}}e^{i\mathrm{\Delta }_k^qt_c}M_{k+q}^qM_k^q`$ (6)
$`{\displaystyle \frac{dM_k^q}{dt_c}}`$ $`=`$ $`{\displaystyle \underset{q^{}}{}}(M_{k+q+q^{}}^q^{}M_{k+q^{}}^qM_k^q^{}{\displaystyle \frac{e^{i\mathrm{\Delta }_{k+q}^q^{}t_c}e^{i\mathrm{\Delta }_k^q^{}t_c}}{i(\mathrm{\Delta }_{k+q}^q^{}\mathrm{\Delta }_k^q^{})}}.`$ (8)
$`.+M_k^qM_{k+q^{}}^q^{}M_k^q^{}t_ce^{i\mathrm{\Delta }_k^q^{}t_c}).`$
Here we introduced
$`\mathrm{\Delta }_k^q=ϵ_{k+q}ϵ_k+\omega .`$
The $`(k,q)`$-dependence of the interaction coefficients is generated during the RG-flow. The second term in (8) is a correction term. It is due to the fact that a time interval connected with a contraction becomes a single point in time at one RG step. In the next step this leads to the generation of new terms which were previously not present. The level broadening is also included in (6) since all energies become complex. The terms that generate the double vertex operators are of fourth order in $`M`$. Therefore the equations (6) and (8) contain the order $`M^4`$ exactly.
Feynman’s method. To compare the results of the real-time RG we applied Feyman’s method to the one-dimensional polaron with a finite band-width. Feynman used a variational principle in the path-integral formalism. Within his two-particle approximation one obtains for the ground state energy
$`E_g`$ $`=`$ $`{\displaystyle \frac{(vw)^2}{4v}}\alpha {\displaystyle \frac{v}{\sqrt{\pi }}}{\displaystyle _0^{\mathrm{}}}𝑑\tau e^\tau `$ (10)
$`\times [g(\tau )]^{\frac{1}{2}}[\mathrm{erf}({\displaystyle \frac{Dg(\tau )}{v^2}})^{\frac{1}{2}}]`$
with the band-width $`D`$, the error-function
$`\mathrm{erf}(\mathrm{x})={\displaystyle \frac{2}{\sqrt{\pi }}}{\displaystyle _0^x}𝑑te^{t^2}`$
and
$`g(\tau )=\left({\displaystyle \frac{v^2w^2}{v}}\left(1e^{v\tau }\right)+w^2\tau \right).`$
$`v`$ and $`w`$ are chosen such that $`E_g`$ is minimum. For $`v=w`$ the result of the perturbation theory of first order in $`\alpha `$ is reproduced. Following Feynman we treat small couplings by setting $`v=(1+\delta )w`$ and $`w=1`$ . Considering $`\delta `$ small one can now expand the right-hand-side in (10) and minimize it with respect to $`\delta `$. The resulting energy is
$`E_g`$ $`=`$ $`\alpha {\displaystyle \frac{2}{\pi }}\mathrm{arctan}\sqrt{D}\alpha ^2{\displaystyle \frac{4}{\pi ^2}}(3\mathrm{arctan}\sqrt{D}`$ (12)
$`2\sqrt{2}\mathrm{arctan}\sqrt{{\displaystyle \frac{D}{2}}}{\displaystyle \frac{\sqrt{D}}{D+1}})^2.`$
Within Feynman’s approximation we obtain for the effective mass:
$`m`$ $`=`$ $`1+{\displaystyle \frac{\alpha }{\pi }}\left(\mathrm{arctan}\sqrt{D}+\sqrt{D}{\displaystyle \frac{D1}{(D+1)^2}}\right)`$ (16)
$`+\alpha ^2{\displaystyle \frac{4}{\pi ^2}}\left(3\mathrm{arctan}\sqrt{D}2\sqrt{2}\mathrm{arctan}\sqrt{{\displaystyle \frac{D}{2}}}{\displaystyle \frac{\sqrt{D}}{D+1}}\right)`$
$`\times (3\mathrm{arctan}\sqrt{D}+3\sqrt{2}\mathrm{arctan}\sqrt{{\displaystyle \frac{D}{2}}}`$
$`+\sqrt{D}^3{\displaystyle \frac{3D^3+18D^2+21D2}{(D+1)^3(D+2)^2}}).`$
Since (10) follows from a mininum principle Feynman’s method leads to more accurate results for the energy than for the mass (for quantitative studies concerning the accuracy in three dimensions see ).
Our results for the energy using Feynman’s approach as well as our values following from perturbation theory of second order are in good agreement with .
Results. The differential equations (6) and (8) are solved numerically. Regarding the oscillating terms in these equations one recognizes that given a certain discretization in $`q`$-space one obtains large errors with increasing frequency $`t_c`$. To avoid this the phase $`\mathrm{\Delta }_k^qt_c`$ has been interpolated in $`q`$-space.
Unfortunately (6) and (8) do not show a convergent behaviour for the ground state energy for $`t_c\mathrm{}`$. One reason is that there are undamped modes corresponding to high excitations in the $`q`$-sums leading to increasing effects on the ground state energy. In this context another problem arises from the fact that the correlation function in (5) is not decaying. As a consequence, oscillations decay as a function of $`t_c`$ but reoccur for sufficiently large $`t_c`$ . The idea of our solution is to neglect further renormalizaton effects of both $`\mathrm{\Delta }_k^q`$ and $`M_k^q`$ for $`t_c`$ larger than a certain point $`t_f`$. By doing this we obtained a damped oscillation of the ground state energy $`ϵ_0`$ in $`t_c`$-space for $`t_c>t_f`$. Therefore (6) can be integrated analytically which leads to
$`E_g`$ $``$ $`\underset{t_c\mathrm{}}{lim}ϵ_0(t_c)`$ (17)
$`=`$ $`ϵ_0(t_f){\displaystyle \underset{q}{}}{\displaystyle \frac{M_q^q(t_f)M_0^q(t_f)}{\mathrm{\Delta }_0^q(t_f)}}e^{i\mathrm{\Delta }_0^q(t_f)t_f}.`$ (18)
In Fig. 2 the solution of (18) is shown for different $`t_f`$. The problems mentioned above make it necessary to choose a finite $`t_f`$ where the renormalization effects beyond perturbation theory are contained but the numerical instabilities do not yet occur. We choose $`t_f=2.5`$ for all values of $`\alpha `$. At this point only low excitations ($`\mathrm{\Delta }<0.4`$) are not integrated out yet. Since $`\mathrm{\Delta }1`$ sets the scale for the first excited state, it is reasonable to assume that excited states do not have further important renormalization effects. The change of $`E_g`$ between $`t_f=2`$ and $`t_f=3`$ is approximately $`2\%`$ for $`\alpha =0.5`$.
The $`\alpha `$-dependence of $`E_g`$ is shown in Fig. 3. One notices that for $`\alpha <0.7`$ we obtain lower values for the ground state energy than those of Feynman’s method. For larger couplings our method is no longer reliable and yields worse results (see Fig. 3) which is not surprising since we neglected double and higher order vertex operators.
We also calculated $`E_g`$ for different band-widths $`D`$. As one can see in Fig. 4 for larger band-widths the deviation to the results of Feynman’s method grows.
Another quantity of interest is the mass of the polaron. If one differentiates (6) twice with respect to $`k`$ one obtains a differential equation for $`1/m`$. However, its solution oscillates with increasing amplitude. Therefore for determining the mass we calculate both the ground state energy and low excited energies using the same procedure as above. The resulting dispersion gives us a value for the mass. The accuracy for the mass is worse though. For $`\alpha =0.3`$ the change of $`1/m`$ is approximately $`6\%`$ between $`t_f=2`$ and $`t_f=3`$.
The results for different $`\alpha `$ are shown in Fig. 5. For small couplings ($`\alpha <0.4`$), we find a value for the mass between the variational principle of Feynman and the one of Lee, Low and Pines . For larger couplings the numerical solution is too unstable to make definite statements from the RG-approach. From Fig. 6 we see that our mass depends only slightly on the band-width.
In summary, we used a new renormalization group method to study the polaron problem. Although there arise some difficulties from the structure of the flow-equations, we were able to investigate the energy beyond perturbation theory. The mass was calculated as well, but its accuracy is worse than that of the energy. The results were compared to the values following from perturbation theory, Feynman’s method and the one of Lee, Low and Pines. The performance of the RG is good for couplings $`\alpha 0.5`$. This restriction is due to the neglecting of double and higher order vertex operators and the undamped oscillating correlation function of the phonons.
We acknowledge useful discussions with A. Mielke and K. Schönhammer. This work was supported by the ”Deutsche Forschungsgemeinschaft” as part of ”SFB 345” (M.K.) and “SFB 195” (H.S.)
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# Connected components of real double Bruhat cells
## 1. Introduction
The main geometric objects of study in this paper are *double Bruhat cells* $`G^{u,v}=BuBB_{}vB_{}`$ in a simply-connected connected complex semisimple group $`G`$; here $`B`$ and $`B_{}`$ are two opposite Borel subgroups in $`G`$, and $`u`$ and $`v`$ any two elements of the Weyl group $`W`$. Double Bruhat cells were introduced and studied in as a geometric framework for the study of total positivity in semisimple groups; they are also closely related to symplectic leaves in the corresponding Poisson-Lie groups (see ). It will be convenient for us to replace $`G^{u,v}`$ with a reduced double Bruhat cell $`L^{u,v}`$ introduced in . The variety $`L^{u,v}`$ can be identified with the quotient of $`G^{u,v}`$ modulo the left (or right) action of the maximal torus $`H=BB_{}`$.
As shown in , an algebraic variety $`L^{u,v}`$ is biregularly isomorphic to a Zariski open subset of an affine space of dimension $`m=\mathrm{}(u)+\mathrm{}(v)`$, where $`\mathrm{}(u)`$ is the length of $`u`$ in the Coxeter group $`W`$. However, the smooth topology of $`L^{u,v}`$ can be quite complicated. A first step towards understanding this topology is enumerating the connected components of the real part $`L^{u,v}()`$. In the case when $`G`$ is simply-laced, a conjectural answer was given in \[15, Conjecture 4.1\]. Here we prove this conjecture and extend the result to an arbitrary semisimple group $`G`$. The answer is given in the following terms: as shown in for $`G`$ simply-laced, every reduced word $`𝐢`$ of $`(u,v)W\times W`$ gives rise to a subgroup $`\mathrm{\Gamma }_𝐢(𝔽_2)GL_m(𝔽_2)`$ generated by symplectic transvections (here $`𝔽_2`$ is the $`2`$-element field). We extend the construction of $`\mathrm{\Gamma }_𝐢(𝔽_2)`$ to an arbitrary $`G`$ (it is still generated by transvections but not necessarily by symplectic ones). Extending \[15, Conjecture 4.1\], we show that the connected components of $`L^{u,v}()`$ are in a natural bijection with the $`\mathrm{\Gamma }_𝐢(𝔽_2)`$-orbits in $`𝔽_2^m`$. As explained in , this provides a far-reaching generalization of the results in ; this also refines and generalizes results in .
Our proof uses methods and results developed in . First, it was shown there that every reduced word $`𝐢`$ of $`(u,v)W\times W`$ gives rise to a biregular isomorphism between the complex torus $`_0^m`$ and a Zariski open subset $`U_𝐢L^{u,v}`$. We refine this result by showing that the complement $`L^{u,v}U_𝐢`$ is the union of $`m`$ divisors $`\{M_{k,𝐢}=0\}`$, where $`M_{1,𝐢},\mathrm{},M_{m,𝐢}`$ are some irreducible regular functions on $`L^{u,v}`$. We further show that every $`𝐢`$-bounded index $`n[1,m]`$ (see Section 2 for the definition) gives rise to a regular function $`M_{n,𝐢}^{}`$ on $`L^{u,v}`$ such that replacing the divisor $`\{M_{n,𝐢}=0\}`$ with $`\{M_{n,𝐢}^{}=0\}`$ leads to another “toric chart” $`U_{n,𝐢}`$ in $`L^{u,v}`$. Then we prove that the connected components of the real part of the union of charts $`U_𝐢_nU_{n,𝐢}`$ are in a natural bijection with the $`\mathrm{\Gamma }_𝐢(𝔽_2)`$-orbits in $`𝔽_2^m`$. Finally, we show that the complement in $`L^{u,v}`$ of this union of charts has codimension $`2`$, so the connected components of $`L^{u,v}()`$ are enumerated in the same way as those of the real part of $`U_𝐢_nU_{n,𝐢}`$.
According to , each $`M_{k,𝐢}`$ is a “twisted (generalized) minor” on $`G`$. We show that each $`M_{n,𝐢}^{}`$ is obtained by the same twist from a regular function on $`G`$ which is no longer a minor but can be expressed as a sum of two Laurent monomials in minors. These regular functions are of independent interest for the study of the dual canonical bases in the ring of regular functions $`[G]`$ and its $`q`$-deformation.
The paper is organized as follows. After recalling the necessary background, we formulate our main result (Theorem 2.2) in Section 2. In Section 3, we formulate a lemma (Lemma 3.1) that plays the crucial role in our proof of Theorem 2.2, and then show how this lemma implies the theorem. The proof of Lemma 3.1 is given in Section 4. Finally, Section 5 discusses some examples and applications of the results in Sections 3 and 4.
## 2. Main theorem
To formulate our main result, let us recall the necessary background from . Let $`G`$ be a simply connected semisimple algebraic group with the Dynkin graph $`\mathrm{\Pi }`$. Let $`B`$ and $`B_{}`$ be two $``$-split opposite Borel subgroups, $`N`$ and $`N_{}`$ their unipotent radicals, $`H=BB_{}`$ an $``$-split maximal torus of $`G`$, and $`W=\mathrm{Norm}_G(H)/H`$ the Weyl group of $`G`$. Let $`𝔤=\mathrm{Lie}(G)`$ be the Lie algebra of $`G`$, and $`𝔥=\mathrm{Lie}(H)`$ the Cartan subalgebra of $`𝔤`$. Let $`\{\alpha _i:i\mathrm{\Pi }\}`$ be the system of simple roots in $`𝔥^{}`$ for which the corresponding root subgroups are contained in $`N`$. Let $`\{\alpha _i^{}:i\mathrm{\Pi }\}`$ be the corresponding system of simple coroots in $`𝔥`$, and $`A=(a_{ij}=\alpha _j(\alpha _i^{}))`$ be the *Cartan matrix*. Thus, for $`ij`$ the indices $`i`$ and $`j`$ are adjacent in $`\mathrm{\Pi }`$ if and only if $`a_{ij}a_{ji}0`$; we shall denote this by $`\{i,j\}\mathrm{\Pi }`$. For every $`i\mathrm{\Pi }`$, let $`\phi _i:SL_2G`$ denote the corresponding canonical $`SL_2`$-embedding.
The Weyl group $`W`$ is canonically identified with the Coxeter group $`W(A)`$ generated by the involutions $`s_i`$ for $`i\mathrm{\Pi }`$ subject to the relations $`(s_is_j)^{d_{ij}}=e`$ for all $`ij`$, where $`d_{ij}=2`$ (resp. $`3,4,6`$) if $`a_{ij}a_{ji}=0`$ (resp. $`1,2,3`$). A word $`𝐢=(i_1,\mathrm{},i_m)`$ in the alphabet $`\mathrm{\Pi }`$ is a *reduced word* for $`wW`$ if $`w=s_{i_1}\mathrm{}s_{i_m}`$, and $`m`$ is the smallest length of such a factorization. The length $`m`$ of any reduced word for $`w`$ is called the *length* of $`w`$ and denoted by $`m=\mathrm{}(w)`$. Let $`R(w)`$ denote the set of all reduced words for $`w`$. The identification $`W=W(A)`$ is given by $`s_i=\overline{s_i}H`$, where
$$\overline{s_i}=\phi _i(\begin{array}{cc}0& 1\\ 1& 0\end{array})\mathrm{Norm}_G(H).$$
The representatives $`\overline{s_i}G`$ satisfy the braid relations $`\overline{s_i}\overline{s_j}\overline{s_i}\mathrm{}=\overline{s_j}\overline{s_i}\overline{s_j}\mathrm{}`$ (with $`d_{ij}`$ factors on each side); thus, the representative $`\overline{w}`$ can be unambiguously defined for any $`wW`$ by requiring that $`\overline{uv}=\overline{u}\overline{v}`$ whenever $`\mathrm{}(uv)=\mathrm{}(u)+\mathrm{}(v)`$.
The “double” group $`W\times W`$ is also a Coxeter group. The corresponding graph $`\stackrel{~}{\mathrm{\Pi }}`$ is the union of two disconnected copies of $`\mathrm{\Pi }`$. We identify the vertex set of $`\stackrel{~}{\mathrm{\Pi }}`$ with $`\{+1,1\}\times \mathrm{\Pi }`$, and write a vertex $`(\pm 1,i)\stackrel{~}{\mathrm{\Pi }}`$ simply as $`\pm i`$. For each $`i\mathrm{\Pi }`$, we set $`\epsilon (\pm i)=\pm 1`$ and $`|\pm i|=i`$. Thus, two vertices $`i`$ and $`j`$ of $`\stackrel{~}{\mathrm{\Pi }}`$ are adjacent if and only if $`\epsilon (i)=\epsilon (j)`$ and $`\{|i|,|j|\}\mathrm{\Pi }`$. In this notation, a reduced word for a pair $`(u,v)W\times W`$ is an arbitrary shuffle of a reduced word for $`u`$ written in the alphabet $`\mathrm{\Pi }`$ and a reduced word for $`v`$ written in the alphabet $`\mathrm{\Pi }`$.
The group $`G`$ has two *Bruhat decompositions*, with respect to $`B`$ and $`B_{}`$:
$$G=\underset{uW}{}BuB=\underset{vW}{}B_{}vB_{}.$$
The *double Bruhat cells* $`G^{u,v}`$ are defined by $`G^{u,v}=BuBB_{}vB_{}`$.
Following , we define the *reduced double Bruhat cell* $`L^{u,v}G^{u,v}`$ as follows:
(2.1)
$$L^{u,v}=N\overline{u}NB_{}vB_{}.$$
The maximal torus $`H`$ acts freely on $`G^{u,v}`$ by left (or right) translations, and $`L^{u,v}`$ is a section of this action. Thus, $`G^{u,v}`$ is biregularly isomorphic to $`H\times L^{u,v}`$, and all properties of $`G^{u,v}`$ can be translated in a straightforward way into the corresponding properties of $`L^{u,v}`$ (and vice versa). In particular, Theorem 1.1 in implies that $`L^{u,v}`$ is biregularly isomorphic to a Zariski open subset of an affine space of dimension $`\mathrm{}(u)+\mathrm{}(v)`$.
The *real part* of $`G`$ is the subgroup $`G()`$ of $`G`$ generated by all the subgroups $`\phi _i(SL_2())`$. For any subset $`LG`$, we define its real part by $`L()=LG()`$.
Now let us fix a pair $`(u,v)W\times W`$, and let $`m=\mathrm{}(u)+\mathrm{}(v)`$. Let $`𝐢=(i_1,\mathrm{},i_m)R(u,v)`$ be any reduced word for $`(u,v)`$. We associate to $`𝐢`$ an $`m\times m`$ matrix $`(C_{kl})`$ in the following way: set $`C_{kl}=1`$ if $`|i_k|=|i_l|`$ and $`C_{kl}=a_{|i_k|,|i_l|}`$ if $`|i_k||i_l|`$.
Following , we associate with $`𝐢`$ a directed graph $`\mathrm{\Sigma }(𝐢)`$ on the set of vertices $`[1,m]=\{1,2,\mathrm{},m\}`$. For $`l[1,m]`$, we denote by $`l^{}=l_𝐢^{}`$ the maximal index $`k`$ such that $`1k<l`$ and $`|i_k|=|i_l|`$; if $`|i_k||i_l|`$ for $`1k<l`$ then we set $`l^{}=0`$. The edges of $`\mathrm{\Sigma }(𝐢)`$ are now defined as follows.
###### Definition 2.1.
A pair $`\{k,l\}[1,m]`$ with $`k<l`$ is an edge of $`\mathrm{\Sigma }(𝐢)`$ if it satisfies one of the following three conditions:
(i) $`k=l^{}`$;
(ii) $`k^{}<l^{}<k`$, $`\{|i_k|,|i_l|\}\mathrm{\Pi }`$, and $`\epsilon (i_l^{})=\epsilon (i_k)`$;
(iii) $`l^{}<k^{}<k`$, $`\{|i_k|,|i_l|\}\mathrm{\Pi }`$, and $`\epsilon (i_k^{})=\epsilon (i_k)`$.
The edges of type (i) are called *horizontal*, and those of types (ii) and (iii) *inclined*. A horizontal (resp. inclined) edge $`\{k,l\}`$ with $`k<l`$ is directed from $`k`$ to $`l`$ if and only if $`\epsilon (i_k)=+1`$ (resp. $`\epsilon (i_k)=1`$). We shall write $`(kl)\mathrm{\Sigma }(𝐢)`$ if $`kl`$ is a directed edge of $`\mathrm{\Sigma }(𝐢)`$.
We now associate to each $`n[1,m]`$ a transvection $`\tau _n=\tau _{n,𝐢}:^m^m`$ defined as follows: if $`\tau _n(\xi _1,\mathrm{},\xi _m)=(\xi _1^{},\mathrm{},\xi _m^{})`$ then $`\xi _k^{}=\xi _k`$ for $`kn`$, and
(2.2)
$$\xi _n^{}=\xi _n\underset{(kn)\mathrm{\Sigma }(𝐢)}{}C_{kn}\xi _k+\underset{(nl)\mathrm{\Sigma }(𝐢)}{}C_{ln}\xi _l$$
(note that if $`G`$ is simply-laced then all the coefficients $`C_{kn}`$ and $`C_{ln}`$ in (2.2)are equal to $`1`$, so (2.2) becomes formula (2.4) in ). We call an index $`n[1,m]`$ *$`𝐢`$-bounded* if $`n^{}>0`$. Let $`\mathrm{\Gamma }_𝐢`$ denote the group of linear transformations of $`^m`$ generated by the transvections $`\tau _n`$ for all $`𝐢`$-bounded indices $`n[1,m]`$. Let $`\mathrm{\Gamma }_𝐢(𝔽_2)`$ denote the group of linear transformations of the $`𝔽_2`$-vector space $`𝔽_2^m`$ obtained from $`\mathrm{\Gamma }_𝐢`$ by reduction modulo $`2`$ (recall that $`𝔽_2`$ is the $`2`$-element field).
We are finally ready to formulate our main result.
###### Theorem 2.2.
For every reduced word $`𝐢R(u,v)`$, the connected components of $`L^{u,v}()`$ are in a natural bijection with the $`\mathrm{\Gamma }_𝐢(𝔽_2)`$-orbits in $`𝔽_2^m`$.
Note that in Theorem 2.2 we only need the modulo $`2`$ reductions of transvections $`\tau _n`$, so the formula (2.2) could be simplified as follows:
$$\xi _n^{}=\xi _n+\underset{(k,n)\mathrm{\Sigma }(𝐢)}{}C_{kn}\xi _k.$$
We prefer the form (2.2) because it is suggested by the construction of toric charts in $`L^{u,v}`$ which is our main ingredient in proving Theorem 2.2.
## 3. Main lemma
As before, let $`G`$ be a simply connected connected complex semisimple group with the Dynkin graph $`\mathrm{\Pi }`$. We fix a pair $`(u,v)W\times W`$, let $`m=\mathrm{}(u)+\mathrm{}(v)`$, and fix a reduced word $`𝐢=(i_1,\mathrm{},i_m)R(u,v)`$.
###### Lemma 3.1.
There exist regular functions $`M_1,\mathrm{},M_m`$ on $`L^{u,v}`$ with the following properties:
(1) If $`k[1,m]`$ is not $`𝐢`$-bounded then $`M_k`$ vanishes nowhere on $`L^{u,v}`$.
(2) The map $`(M_1,\mathrm{},M_m):L^{u,v}^m`$ restricts to a biregular isomorphism $`U_𝐢_0^m`$, where $`U_𝐢`$ is the locus of all $`xL^{u,v}`$ such that $`M_k(x)0`$ for all $`k[1,m]`$.
(3) For every $`𝐢`$-bounded $`n[1,m]`$, the rational function $`M_k^{}`$ defined by
(3.1)
$$M_n^{}M_n=\underset{(kn)\mathrm{\Sigma }_𝐢}{}M_k^{C_{kn}}+\underset{(nl)\mathrm{\Sigma }_𝐢}{}M_l^{C_{ln}}$$
is regular on $`L^{u,v}`$.
(4) For every $`𝐢`$-bounded $`n[1,m]`$, the map $`(M_1,\mathrm{},M_{n1},M_n^{},M_{n+1},\mathrm{},M_m):L^{u,v}^m`$ restricts to a biregular isomorphism $`U_{n,𝐢}_0^m`$, where $`U_{n,𝐢}`$ is the locus of all $`xL^{u,v}`$ such that $`M_n^{}(x)0`$ and $`M_k(x)0`$ for all $`k[1,m]\{n\}`$.
(5) The functions $`M_k`$ and $`M_n^{}`$ take real values on $`L^{u,v}()`$, and the biregular isomorphisms in (2) and (4) restrict to biregular isomorphisms $`U_𝐢()_0^m`$ and $`U_{n,𝐢}()_0^m`$.
The functions $`M_k=M_{k,𝐢}`$ in Lemma 3.1 were introduced in \[4, (4.13)\]. We recall the definition and prove Lemma 3.1 in the next section; in the rest of this section we show that it implies Theorem 2.2. To be more precise, we shall prove that the bijection in Theorem 2.2 can be defined as follows. For every $`\xi =(\xi _1,\mathrm{},\xi _m)𝔽_2^m`$, let $`U_𝐢(\xi )`$ denote the set of all $`xU_𝐢()`$ such that $`(1)^{\xi _k}M_k(x)>0`$ for all $`k`$. For every $`YL^{u,v}()`$, let $`\overline{Y}`$ denote the closure of $`Y`$ in $`L^{u,v}()`$ in the real topology.
###### Theorem 3.2.
The correspondence $`\mathrm{\Omega }_{\xi \mathrm{\Omega }}\overline{U_𝐢(\xi )}`$ is a bijection between $`\mathrm{\Gamma }_𝐢(𝔽_2)`$-orbits in $`𝔽_2^m`$ and connected components of $`L^{u,v}()`$.
We split the proof of Theorem 3.2 into several lemmas. Let us abbreviate $`X=L^{u,v}`$, and let $`[X]`$ be the ring of regular functions on $`X`$. Since $`X`$ is isomorphic to a Zariski open subset of $`^m`$, the ring $`[X]`$ is a unique factorization domain. By property (1) in Lemma 3.1, if $`k`$ is not $`𝐢`$-bounded then $`M_k`$ is an invertible element of $`[X]`$.
###### Lemma 3.3.
A Laurent monomial $`P=M_1^{d_1}\mathrm{}M_m^{d_m}`$ is a regular function on $`X`$ if and only if $`d_n0`$ for any $`𝐢`$-bounded $`n`$.
###### Proof.
The “if” part is trivial. To prove the “only if” part, fix an $`𝐢`$-bounded index $`n`$, and consider the restriction of $`P`$ to the Zariski open subset $`U_{n,𝐢}X`$. By property (4) in Lemma 3.1, if $`P[X]`$ then $`M_n^{d_n}`$ is a regular function on $`U_{n,𝐢}`$ and so it must be a Laurent polynomial in $`M_1,\mathrm{},M_{n1},M_n^{},M_{n+1},\mathrm{},M_m`$. In view of (3.1), this implies that $`d_n0`$, as desired. ∎
###### Lemma 3.4.
For every $`𝐢`$-bounded $`n`$, the function $`M_n`$ is an irreducible element of $`[X]`$.
###### Proof.
Notice that every $`P[X]`$ restricts to a regular function on the Zariski open subset $`U_𝐢X`$. By property (2) in Lemma 3.1, $`P`$ is a Laurent polynomial in $`M_1,\mathrm{},M_m`$. It follows that if $`M_n`$ is the product of two regular functions $`P`$ and $`Q`$ then both $`P`$ and $`Q`$ must be Laurent monomials in $`M_1,\mathrm{},M_m`$. By Lemma 3.3, one of the factors $`P`$ and $`Q`$ must be a Laurent monomial in the variables $`M_k`$ for $`k`$ not $`𝐢`$-bounded, hence is an invertible element of $`[X]`$. Therefore, $`M_n`$ is irreducible. ∎
###### Lemma 3.5.
For every $`𝐢`$-bounded $`n`$, the function $`M_n^{}`$ is equal to some irreducible element $`M_n^{\prime \prime }[X]`$ times a Laurent monomial in $`M_1,\mathrm{},M_{n1},M_{n+1},\mathrm{},M_m`$.
###### Proof.
Let $`P[X]`$ be an irreducible factor of $`M_n^{}`$. Restricting $`P`$ to $`U_{n,𝐢}`$, we conclude that $`P`$ is a Laurent monomial in $`M_1,\mathrm{},M_{n1},M_n^{},M_{n+1},\mathrm{},M_m`$. Restricting $`P`$ to $`U_𝐢`$ and using property (2) in Lemma 3.1, we see that $`P`$ must be also a Laurent polynomial in $`M_1,\mathrm{},M_m`$. By (3.1), this implies that the exponent of $`M_n^{}`$ in $`P`$ written as a Laurent monomial in $`M_1,\mathrm{},M_{n1},M_n^{},M_{n+1},\mathrm{},M_m`$ must be nonnegative. It follows that there is an irreducible factor $`M_n^{\prime \prime }`$ of $`M_n^{}`$ which is equal to $`M_n^{}`$ times a Laurent monomial in $`M_1,\mathrm{},M_{n1},M_{n+1},\mathrm{},M_m`$, while the rest of the factors are just Laurent monomials in $`M_1,\mathrm{},M_{n1},M_{n+1},\mathrm{},M_m`$. ∎
We set $`U=U_𝐢_nU_{n,𝐢}`$.
###### Lemma 3.6.
The complement $`XU`$ is the locus of all $`xX`$ such that $`M_n(x)=M_k(x)=0`$ for two distinct $`𝐢`$-bounded indices $`n`$ and $`k`$, or $`M_n(x)=M_n^{\prime \prime }(x)=0`$ for some $`𝐢`$-bounded $`n`$. The variety $`XU`$ has (complex) codimension $`2`$ in $`X`$.
###### Proof.
Suppose $`xXU`$. Since $`xU_𝐢`$, property (1) in Lemma 3.1 implies that $`M_n(x)=0`$ for some $`𝐢`$-bounded $`n`$. Since $`xU_{n,𝐢}`$, it follows that either $`M_n^{\prime \prime }(x)=0`$, or $`M_k(x)=0`$ for some $`𝐢`$-bounded $`kn`$.
The converse inclusion is obvious. Finally, the statement that $`XU`$ has codimension $`2`$ in $`X`$ is clear since $`XU`$ is the union of finitely many subvarieties, each given by two (distinct) irreducible equations. ∎
Now consider the real part $`U()=U_𝐢()_nU_{n,𝐢}()`$. By Lemma 3.6 and property (5) in Lemma 3.1, the complement $`X()U()`$ has real codimension $`2`$ in $`X()`$. Therefore, the connected components of $`X()`$ (in the real topology) are closures of the connected components of $`U()`$. It remains to show that Theorem 3.2 holds with $`X()`$ replaced by $`U()`$. For a subset $`YU()`$ we now denote by $`\overline{Y}`$ the closure of $`Y`$ in $`U()`$. The role of transvections $`\tau _n`$ is explained by the following lemma.
###### Lemma 3.7.
Let $`\xi ^{(1)}`$ and $`\xi ^{(2)}`$ be two distinct vectors in $`𝔽_2^m`$. Then $`\overline{U(\xi ^{(1)})}\overline{U(\xi ^{(2)})}\mathrm{}`$ if and only if $`\xi ^{(2)}=\tau _n(\xi ^{(1)})`$ for some $`𝐢`$-bounded index $`n`$.
###### Proof.
Suppose $`xU()`$ belongs to the intersection $`\overline{U(\xi ^{(1)})}\overline{U(\xi ^{(2)})}`$. Then $`M_k(x)=0`$ whenever $`\xi _k^{(1)}\xi _k^{(2)}`$. Using Lemma 3.6, we see that there is a unique $`n`$ such that $`\xi _n^{(1)}\xi _n^{(2)}`$; furthermore, this index $`n`$ is $`𝐢`$-bounded, and $`M_n^{}(x)0`$. Since any neighborhood of $`x`$ intersects both $`U(\xi ^{(1)})`$ and $`U(\xi ^{(2)})`$, it follows that the two monomials on the right hand side of (3.1) must have opposite signs at $`x`$. Let us write $`\xi _k=\xi _k^{(1)}=\xi _k^{(2)}`$ for $`kn`$. Then we have
$$\xi _n^{(2)}\xi _n^{(1)}=1=\underset{(nl)\mathrm{\Sigma }(𝐢)}{}C_{ln}\xi _l\underset{(kn)\mathrm{\Sigma }(𝐢)}{}C_{kn}\xi _k.$$
Comparing this with (2.2), we conclude that $`\xi ^{(2)}=\tau _n(\xi ^{(1)})`$, as claimed.
Conversely, suppose $`\xi ^{(2)}=\tau _n(\xi ^{(1)})\xi ^{(1)}`$, and let $`\xi _k=\xi _k^{(1)}=\xi _k^{(2)}`$ for $`kn`$. Then
$$\underset{(nl)\mathrm{\Sigma }(𝐢)}{}C_{ln}\xi _l\underset{(kn)\mathrm{\Sigma }(𝐢)}{}C_{kn}\xi _k.$$
This implies that there exists a point $`xU_{n,𝐢}()`$ such that $`(1)^{\xi _k}M_k(x)>0`$ for all $`kn`$, and the right hand side of (3.1) vanishes at $`x`$. Then any neighborhood of $`x`$ contains points with the signs of all $`M_k`$ for $`kn`$ unchanged and with the right hand side of (3.1 positive (as well as negative). Thus, $`x\overline{U(\xi ^{(1)})}\overline{U(\xi ^{(2)})}`$, and we are done. ∎
Now we are ready to complete the proof of Theorem 3.2. Let $`\mathrm{\Omega }`$ be a $`\mathrm{\Gamma }_𝐢(𝔽_2)`$-orbit in $`𝔽_2^m`$, and consider the corresponding closed subset $`Y_\mathrm{\Omega }=_{\xi \mathrm{\Omega }}\overline{U_𝐢(\xi )}`$ of $`U()`$. Each $`U_𝐢(\xi )`$ is a copy of $`_{>0}^m`$ and so is connected. Using the “if” part of Lemma 3.7, we conclude that $`Y_\mathrm{\Omega }`$ is connected (since the closure of a connected set and the union of two non-disjoint connected sets are connected as well). On the other hand, by the “only if” part of the same lemma, all the sets $`Y_\mathrm{\Omega }`$ are pairwise disjoint. Thus, they are the connected components of $`U()`$, and we are done.
## 4. Proof of Lemma 3.1
### 4.1. The functions $`M_k`$
We start by recalling the definition of the functions $`M_k=M_{k,𝐢}`$ given in \[4, (4.13)\]. First of all, recall that the *weight lattice* $`P`$ of $`G`$ can be thought of as the group of rational multiplicative characters of $`H`$ written in the exponential notation: a weight $`\gamma P`$ acts by $`aa^\gamma `$. The lattice $`P`$ is also identified with the additive group of all $`\gamma 𝔥^{}`$ such that $`\gamma (\alpha _i^{})`$ for all $`i\mathrm{\Pi }`$. Thus, $`P`$ has a $``$-basis $`\{\omega _i:i\mathrm{\Pi }\}`$ of *fundamental weights* given by $`\omega _j(\alpha _i^{})=\delta _{i,j}`$.
We now recall from the definition of *generalized minors*. Denote by $`G_0=N_{}HN`$ the open subset of elements $`xG`$ that have Gaussian decomposition; this (unique) decomposition will be written as $`x=[x]_{}[x]_0[x]_+`$. For $`u,vW`$ and $`i\mathrm{\Pi }`$, the *(generalized) minor* $`\mathrm{\Delta }_{u\omega _i,v\omega _i}`$ is the regular function on $`G`$ whose restriction to the open set $`\overline{u}G_0\overline{v}^1`$ is given by
(4.1)
$$\mathrm{\Delta }_{u\omega _i,v\omega _i}(x)=(\left[\overline{u}^1x\overline{v}\right]_0)^{\omega _i}.$$
As shown in , $`\mathrm{\Delta }_{u\omega _i,v\omega _i}`$ depends on the weights $`u\omega _i`$ and $`v\omega _i`$ alone, not on the particular choice of $`u`$ and $`v`$. It is easy to see that the generalized minors are distinct irreducible elements of the ring $`[G]`$ of regular functions on $`G`$. In the special case $`G=SL_n`$, the generalized minors are nothing but the ordinary minors of a matrix.
According to \[4, Proposition 4.3\], an element $`xG^{u,v}`$ belongs to $`L^{u,v}`$ if and only if $`[\overline{u}^1x]_0=1`$, or equivalently if $`\mathrm{\Delta }_{u\omega _i,\omega _i}(x)=1`$ for any $`i[1,r]`$.
We fix a pair $`(u,v)W\times W`$ and a double reduced word $`𝐢=(i_1,\mathrm{},i_m)R(u,v)`$. Recall from \[4, Definition 4.6 and Theorem 4.7\] that there is a biregular isomorphism $`\psi ^{u,v}`$ between $`L^{u,v}`$ and $`L^{v,u}`$ given by
(4.2)
$$\psi ^{u,v}(x)=[(\overline{v}x^\iota )^1]_+\overline{v}([\overline{u}^1x]_+)^\iota ;$$
here $`xx^\iota `$ is the involutive antiautomorphism of $`G`$ given by
(4.3)
$$\phi _i(\begin{array}{cc}a& b\\ c& d\end{array})^\iota =\phi _i(\begin{array}{cc}d& b\\ c& a\end{array}).$$
Recall that the length $`m`$ of $`𝐢`$ is equal to $`\mathrm{}(u)+\mathrm{}(v)`$. For $`k[1,m]`$, denote
(4.4)
$$u_k=\underset{l=m,\mathrm{},k}{}_{i_l\mathrm{\Pi }}s_{|i_l|},v_{<k}=\underset{l=1,\mathrm{},k1}{}_{i_l\mathrm{\Pi }}s_{i_l}.$$
This notation means that in the first (resp. second) product in (4.4), the index $`l`$ is decreasing (resp. increasing); for example, if $`\mathrm{\Pi }=\{1,2,3\}`$ and $`𝐢=(2,1,3,3,2,1,2,1,1)`$, then, say, $`u_7=s_1s_2`$ and $`v_{<7}=s_1s_3s_2`$.
Following \[4, (4.13)\], we define a regular function $`M_k=M_{k,𝐢}`$ on $`L^{u,v}`$ by
(4.5)
$$M_k(x)=M_{k,𝐢}(x)=\mathrm{\Delta }_{v_{<k}\omega _{|i_k|},u_k\omega _{|i_k|}}(\psi ^{u,v}(x)).$$
### 4.2. Properties (1), (2) and (5)
To prove property (1) in Lemma 3.1, notice that if $`k`$ is not $`𝐢`$-bounded and $`|i_k|=i`$ then (4.5) turns into $`M_{k,𝐢}(x)=\mathrm{\Delta }_{\omega _i,u^1\omega _i}(\psi ^{u,v}(x))`$. It remains to show that $`\mathrm{\Delta }_{\omega _i,u^1\omega _i}`$ vanishes nowhere on $`L^{v,u}`$. In fact, a stronger statement holds: $`\mathrm{\Delta }_{\omega _i,u^1\omega _i}`$ vanishes nowhere on $`B_{}uB_{}`$. This follows from the definition (4.1) and the well-known inclusion $`B_{}uB_{}u^1G_0`$ (cf. \[6, Proposition 2.10\]).
As for property (2) in Lemma 3.1, it follows from the solution to the so-called factorization problem given in \[6, Theorem 1.9\] (or rather from its modification in \[4, Theorem 4.8\]). To formulate it, we need some notation.
For every $`i\mathrm{\Pi }`$ and $`t_0`$, we denote
$$x_i(t)=\phi _i(\begin{array}{cc}1& t\\ 0& 1\end{array}),y_i(t)=\phi _i(\begin{array}{cc}1& 0\\ t& 1\end{array}),t^{\alpha _i^{}}=\phi _i(\begin{array}{cc}t& 0\\ 0& t^1\end{array});$$
following , we also denote
(4.6)
$$x_i(t)=y_i(t)t^{\alpha _i^{}}=\phi _i(\begin{array}{cc}t^1& 0\\ 1& t\end{array}).$$
For any word $`𝐢=(i_1,\mathrm{},i_m)`$ in the alphabet $`\stackrel{~}{\mathrm{\Pi }}`$, let us define the *product map* $`x_𝐢:_0^mG`$ by
(4.7)
$$x_𝐢(t_1,\mathrm{},t_m)=x_{i_1}(t_1)\mathrm{}x_{i_m}(t_m).$$
For $`k[1,m]`$, we denote $`k^+=\mathrm{min}\{l:l>k,|i_l|=|i_k|\}`$, so that $`k^+`$ is the next occurrence of an index $`\pm i_k`$ in $`𝐢`$; if $`k`$ is the last occurrence of $`\pm i_k`$ in $`𝐢`$ then we set $`k^+=m+1`$. We also adopt the convention that $`M_{m+1}(x)=1`$.
The following reformulation of Theorem 4.8 in provides a refinement of property (2) in Lemma 3.1.
###### Theorem 4.1.
Let $`𝐢=(i_1,\mathrm{},i_m)`$ be a double reduced word for $`(u,v)`$, and let $`(M_1,\mathrm{},M_m)_0^m`$. Then there is a unique $`xL^{u,v}`$ such that $`M_{k,𝐢}(x)=M_k`$ for $`k[1,m]`$. This element $`x`$ has the form $`x=x_𝐢(t_1,\mathrm{},t_m)`$, with the factorization parameters $`t_k`$ given by: if $`i_k\mathrm{\Pi }`$ then
(4.8)
$$t_k=M_k/M_{k^+};$$
if $`i_k\mathrm{\Pi }`$ then
(4.9)
$$t_k=\frac{1}{M_kM_{k^+}}\underset{l:l^{}<k<l}{}M_l^{a_{|i_l|,i_k}}.$$
###### Remark 4.2.
We see that the parameters $`t_1,\mathrm{},t_m`$ in the factorization $`x=x_𝐢(t_1,\mathrm{},t_m)`$ are related to $`M_1,\mathrm{},M_m`$ by an invertible monomial transformation. The inverse of this monomial transformation can be computed explicitly: a direct calculation shows that
(4.10)
$$M_k=M_{k,𝐢}(x)=\underset{lk}{}t_l^{\epsilon (i_l)v_{<l}^1v_{<k}\omega _{|i_k|}(\alpha _{|i_l|}^{})}.$$
Finally, property (5) in Lemma 3.1 is clear since each $`M_k`$ is just a Laurent monomial in the factorization parameters $`t_1,\mathrm{},t_m`$, while each $`M_n^{}`$ is the sum of two Laurent monomials; therefore they take real values when all $`t_k`$ are real.
### 4.3. Property (3)
To prove property (3) in Lemma 3.1, we shall construct a new family of regular functions on the whole group $`G`$. Let $`𝐢=(i_1,\mathrm{},i_m)`$ be a reduced word for $`(u,v)W\times W`$ such that $`|i_1|=|i_m|=i`$ for some $`i\mathrm{\Pi }`$, and $`|i_k|i`$ for $`1<k<m`$. Let $`E_\pm =\{k[2,m1]:\epsilon (i_k)=\pm 1\}`$, and let $`J_\pm =\{i\}\{|i_k|:kE_\pm ,k^+=m+1\}\mathrm{\Pi }`$. Let $`\mathrm{\Delta }^{}=\mathrm{\Delta }_𝐢^{}`$ be the rational function on $`G`$ defined by one of the following four equations.
Case 1. If $`i_1=i_m=i`$ then
$`\mathrm{\Delta }^{}\mathrm{\Delta }_{s_i\omega _i,\omega _i}=\mathrm{\Delta }_{\omega _i,\omega _i}{\displaystyle \underset{kE_+}{}_{k^+E_{}\{m+1\}}}\mathrm{\Delta }_{v_k\omega _{i_k},u_{>k}\omega _{i_k}}^{a_{i_k,i}}`$
$`+\mathrm{\Delta }_{v\omega _i,\omega _i}{\displaystyle \underset{kE_+}{}_{k^{}E_{}\{0\}}}\mathrm{\Delta }_{v_{<k}\omega _{i_k},u_{>k}\omega _{i_k}}^{a_{i_k,i}}.`$
Case 2. If $`i_1=i`$ and $`i_m=i`$ then
$`\mathrm{\Delta }^{}\mathrm{\Delta }_{s_i\omega _i,s_i\omega _i}=\mathrm{\Delta }_{\omega _i,s_i\omega _i}\mathrm{\Delta }_{s_i\omega _i,\omega _i}{\displaystyle \underset{kE_+}{}_{k^+E_{}}}\mathrm{\Delta }_{v_k\omega _{i_k},u_{>k}\omega _{i_k}}^{a_{i_k,i}}`$
$`+{\displaystyle \underset{kE_+}{}_{k^{}E_{}\{0\}}}\mathrm{\Delta }_{v_{<k}\omega _{i_k},u_{>k}\omega _{i_k}}^{a_{i_k,i}}{\displaystyle \underset{j\mathrm{\Pi }J_+}{}}\mathrm{\Delta }_{v\omega _j,\omega _j}^{a_{ji}}.`$
Case 3. If $`i_1=i`$ and $`i_m=i`$ then
$`\mathrm{\Delta }^{}\mathrm{\Delta }_{\omega _i,\omega _i}=\mathrm{\Delta }_{\omega _i,u^1\omega _i}\mathrm{\Delta }_{v\omega _i,\omega _i}{\displaystyle \underset{kE_+}{}_{k^{}E_{}}}\mathrm{\Delta }_{v_{<k}\omega _{i_k},u_{>k}\omega _{i_k}}^{a_{i_k,i}}`$
$`+{\displaystyle \underset{kE_{}}{}_{k^{}E_+\{0\}}}\mathrm{\Delta }_{v_{<k}\omega _{|i_k|},u_k\omega _{|i_k|}}^{a_{|i_k|,i}}{\displaystyle \underset{j\mathrm{\Pi }J_{}}{}}\mathrm{\Delta }_{v\omega _j,\omega _j}^{a_{ji}}.`$
Case 4. If $`i_1=i_m=i`$ then
$`\mathrm{\Delta }^{}\mathrm{\Delta }_{\omega _i,s_i\omega _i}=\mathrm{\Delta }_{\omega _i,u^1\omega _i}{\displaystyle \underset{kE_{}}{}_{k^+E_+\{m+1\}}}\mathrm{\Delta }_{v_{<k}\omega _{|i_k|},u_{>k}\omega _{|i_k|}}^{a_{|i_k|,i}}`$
$`+\mathrm{\Delta }_{\omega _i,\omega _i}{\displaystyle \underset{kE_{}}{}_{k^{}E_+\{0\}}}\mathrm{\Delta }_{v_{<k}\omega _{|i_k|},u_k\omega _{|i_k|}}^{a_{|i_k|,i}}.`$
###### Theorem 4.3.
In each of the above four cases, $`\mathrm{\Delta }^{}=\mathrm{\Delta }_𝐢^{}`$ is a regular function on $`G`$.
Before proving Theorem 4.3, we show that it implies property (3) in Lemma 3.1. Let $`(u,v)`$ be an arbitrary pair of elements of $`W`$, and fix a reduced word $`𝐢=(i_1,\mathrm{},i_m)R(u,v)`$. Let $`n`$ be an $`𝐢`$-bounded index in $`[1,m]`$, and let $`𝐢^{}`$ denote the subword $`(i_n^{},\mathrm{},i_n)`$ of $`𝐢`$. We claim that the rational function $`M_n^{}`$ on $`L^{u,v}`$ defined by (3.1) is given by
(4.11)
$$M_n^{}(x)=M_{n,𝐢}^{}(x)=\mathrm{\Delta }_𝐢^{}^{}(\overline{u_{>n}}^1\psi ^{u,v}(x)\overline{v_{<n^{}}}),$$
and so is regular. To see this, let us evaluate the defining equation for $`\mathrm{\Delta }_𝐢^{}^{}`$ at the point $`\overline{u_{>n}}^1\psi ^{u,v}(x)\overline{v_{<n^{}}}`$. Remembering the definition (4.1) of generalized minors, and the definition (4.5) of the functions $`M_k`$, a direct check shows that, in each of the above four cases, the corresponding equality turns into the equation (3.1) with $`M_n^{}`$ given by (4.11).
It remains to prove Theorem 4.3. Our main tool will be the following identity established in \[6, Theorem 1.17\]:
(4.12)
$$\mathrm{\Delta }_{v^{}\omega _i,u^{}\omega _i}\mathrm{\Delta }_{v^{}s_i\omega _i,u^{}s_i\omega _i}\mathrm{\Delta }_{v^{}s_i\omega _i,u^{}\omega _i}\mathrm{\Delta }_{v^{}\omega _i,u^{}s_i\omega _i}=\underset{j\mathrm{\Pi }\{i\}}{}\mathrm{\Delta }_{u^{}\omega _j,v^{}\omega _j}^{a_{ji}}$$
for any $`u^{},v^{}W`$ and $`i\mathrm{\Pi }`$ such that $`\mathrm{}(u^{}s_i)=\mathrm{}(u^{})+1`$ and $`\mathrm{}(v^{}s_i)=\mathrm{}(v^{})+1`$.
To prove that $`\mathrm{\Delta }_𝐢^{}`$ is regular on $`G`$, we first consider the case when $`𝐢`$ is “non-mixed,” i.e., $`k<l`$ for each $`kE_{}`$ and $`lE_+`$. Then the defining equation for $`\mathrm{\Delta }^{}=\mathrm{\Delta }_𝐢^{}`$ simplifies as follows. Denote $`S_\pm =\{|i_k|:kE_\pm \}\mathrm{\Pi }`$.
Case 1 (non-mixed). If $`i_1=i_m=i`$ then
$`\mathrm{\Delta }^{}\mathrm{\Delta }_{s_i\omega _i,\omega _i}=\mathrm{\Delta }_{\omega _i,\omega _i}{\displaystyle \underset{jS_+}{}}\mathrm{\Delta }_{v\omega _j,\omega _j}^{a_{ji}}+\mathrm{\Delta }_{v\omega _i,\omega _i}{\displaystyle \underset{jS_+}{}}\mathrm{\Delta }_{\omega _j,\omega _j}^{a_{ji}}.`$
Case 2 (non-mixed). If $`i_1=i`$ and $`i_m=i`$ then
$`\mathrm{\Delta }^{}\mathrm{\Delta }_{s_i\omega _i,s_i\omega _i}=\mathrm{\Delta }_{\omega _i,s_i\omega _i}\mathrm{\Delta }_{s_i\omega _i,\omega _i}+{\displaystyle \underset{j\mathrm{\Pi }\{i\}}{}}\mathrm{\Delta }_{\omega _j,\omega _j}^{a_{ji}}.`$
Case 3 (non-mixed). If $`i_1=i`$ and $`i_m=i`$ then
$`\mathrm{\Delta }^{}\mathrm{\Delta }_{\omega _i,\omega _i}=\mathrm{\Delta }_{\omega _i,u^1\omega _i}\mathrm{\Delta }_{v\omega _i,\omega _i}{\displaystyle \underset{jS_+S_{}}{}}\mathrm{\Delta }_{\omega _j,\omega _j}^{a_{ji}}`$
$`+{\displaystyle \underset{jS_{}}{}}\mathrm{\Delta }_{\omega _j,u^1\omega _j}^{a_{ji}}{\displaystyle \underset{j\mathrm{\Pi }(S_{}S_+)}{}}\mathrm{\Delta }_{v\omega _j,\omega _j}^{a_{ji}}.`$
Case 4 (non-mixed). If $`i_1=i_m=i`$ then
$`\mathrm{\Delta }^{}\mathrm{\Delta }_{\omega _i,s_i\omega _i}=\mathrm{\Delta }_{\omega _i,u^1\omega _i}{\displaystyle \underset{jS_{}}{}}\mathrm{\Delta }_{\omega _j,u^1\omega _j}^{a_{ji}}+\mathrm{\Delta }_{\omega _i,\omega _i}{\displaystyle \underset{jS_{}}{}}\mathrm{\Delta }_{\omega _j,u^1\omega _j}^{a_{ji}}.`$
By (4.12), in Case 2 we have $`\mathrm{\Delta }^{}=\mathrm{\Delta }_{\omega _i,\omega _i}`$. Cases 1 and 4 are equivalent to each other in view of the identity $`\mathrm{\Delta }_{\gamma ,\delta }(x^T)=\mathrm{\Delta }_{\delta ,\gamma }(x)`$, where $`xx^T`$ is the involutive antiautomorphism of $`G`$ given by
$$\phi _i(\begin{array}{cc}a& b\\ c& d\end{array})^T=\phi _i(\begin{array}{cc}a& c\\ b& d\end{array})$$
(see \[6, Proposition 2.7\]). It remains to show that $`\mathrm{\Delta }^{}`$ is regular in each of the cases 1 and 3.
Let us start with Case 3. Multiplying both monomials on the right hand side of the corresponding equation with the monomial
$$\underset{j\mathrm{\Pi }\{i\}(S_+S_{})}{}\mathrm{\Delta }_{\omega _j,\omega _j}^{a_{ji}}=\underset{j\mathrm{\Pi }\{i\}S_{}}{}\mathrm{\Delta }_{\omega _j,u^1\omega _j}^{a_{ji}}\underset{jS_{}S_+}{}\mathrm{\Delta }_{v\omega _j,\omega _j}^{a_{ji}}$$
and using (4.12), we obtain
$$\mathrm{\Delta }_{\omega _i,u^1\omega _i}\mathrm{\Delta }_{v\omega _i,\omega _i}\underset{j\mathrm{\Pi }\{i\}}{}\mathrm{\Delta }_{\omega _j,\omega _j}^{a_{ji}}+\underset{j\mathrm{\Pi }\{i\}}{}\mathrm{\Delta }_{\omega _j,u^1\omega _j}^{a_{ji}}\underset{j\mathrm{\Pi }\{i\}}{}\mathrm{\Delta }_{v\omega _j,\omega _j}^{a_{ji}}$$
$$=\mathrm{\Delta }_{\omega _i,u^1\omega _i}\mathrm{\Delta }_{v\omega _i,\omega _i}(\mathrm{\Delta }_{\omega _i,\omega _i}\mathrm{\Delta }_{s_i\omega _i,s_i\omega _i}\mathrm{\Delta }_{s_i\omega _i,\omega _i}\mathrm{\Delta }_{\omega _i,s_i\omega _i})$$
$$+(\mathrm{\Delta }_{\omega _i,\omega _i}\mathrm{\Delta }_{s_i\omega _i,u^1\omega _i}\mathrm{\Delta }_{s_i\omega _i,\omega _i}\mathrm{\Delta }_{\omega _i,u^1\omega _i})(\mathrm{\Delta }_{\omega _i,\omega _i}\mathrm{\Delta }_{v\omega _i,s_i\omega _i}\mathrm{\Delta }_{v\omega _i,\omega _i}\mathrm{\Delta }_{\omega _i,s_i\omega _i})$$
$$=\mathrm{\Delta }_{\omega _i,\omega _i}det\left(\begin{array}{ccc}\mathrm{\Delta }_{\omega _i,u^1\omega _i}& \mathrm{\Delta }_{\omega _i,s_i\omega _i}& \mathrm{\Delta }_{\omega _i,\omega _i}\\ \mathrm{\Delta }_{s_i\omega _i,u^1\omega _i}& \mathrm{\Delta }_{s_i\omega _i,s_i\omega _i}& \mathrm{\Delta }_{s_i\omega _i,\omega _i}\\ 0& \mathrm{\Delta }_{v\omega _i,s_i\omega _i}& \mathrm{\Delta }_{v\omega _i,\omega _i}\end{array}\right).$$
Since all “principal minors” $`\mathrm{\Delta }_{\omega _j,\omega _j}`$ are distinct irreducible elements of $`[G]`$, it follows that $`\mathrm{\Delta }_{\omega _i,\omega _i}`$ is relatively prime with $`_{j\mathrm{\Pi }\{i\}(S_+S_{})}\mathrm{\Delta }_{\omega _j,\omega _j}^{a_{ji}}`$. Therefore,
$$\mathrm{\Delta }^{}=det\left(\begin{array}{ccc}\mathrm{\Delta }_{\omega _i,u^1\omega _i}& \mathrm{\Delta }_{\omega _i,s_i\omega _i}& \mathrm{\Delta }_{\omega _i,\omega _i}\\ \mathrm{\Delta }_{s_i\omega _i,u^1\omega _i}& \mathrm{\Delta }_{s_i\omega _i,s_i\omega _i}& \mathrm{\Delta }_{s_i\omega _i,\omega _i}\\ 0& \mathrm{\Delta }_{v\omega _i,s_i\omega _i}& \mathrm{\Delta }_{v\omega _i,\omega _i}\end{array}\right)/\underset{j\mathrm{\Pi }\{i\}(S_+S_{})}{}\mathrm{\Delta }_{\omega _j,\omega _j}^{a_{ji}}$$
is a regular function on $`G`$, as required.
The argument in Case 1 is similar (and simpler). Let us only give the final answer: the function $`\mathrm{\Delta }^{}`$ is now given by
$$\mathrm{\Delta }^{}=(\mathrm{\Delta }_{\omega _i,\omega _i}\mathrm{\Delta }_{v\omega _i,s_i\omega _i}\mathrm{\Delta }_{v\omega _i,\omega _i}\mathrm{\Delta }_{\omega _i,s_i\omega _i})/\underset{j\mathrm{\Pi }\{i\}S_+}{}\mathrm{\Delta }_{\omega _j,\omega _j}^{a_{ji}},$$
and it is again a regular function on $`G`$, as required.
We shall deduce the general case in Theorem 4.3 from the non-mixed case just considered. Note that every reduced word $`𝐢`$ in each of the cases 1 – 4 is obtained from the corresponding non-mixed word by a sequence of $`2`$-*moves* each of which interchanges a pair of consecutive indices $`i_k`$ and $`i_{k+1}`$ with $`kE_{}`$ and $`k+1E_+`$. It suffices to show that if $`𝐢^{}`$ is obtained from $`𝐢`$ by such a move then the regularity of $`\mathrm{\Delta }_𝐢^{}`$ implies that of $`\mathrm{\Delta }_𝐢^{}^{}`$. We shall only treat Case 1; the argument in the other three cases is the same.
Let $`P_1`$ and $`P_2`$ (resp. $`P_1^{}`$ and $`P_2^{}`$) be two monomials in the right hand side of the defining equation for $`\mathrm{\Delta }_𝐢^{}`$ (resp. for $`\mathrm{\Delta }_𝐢^{}^{}`$). Thus, we assume that $`P_1+P_2`$ is divisible by $`\mathrm{\Delta }_{s_i\omega _i,\omega _i}`$ in $`[G]`$, and need to show that the same is true for $`P_1^{}+P_2^{}`$. Let us abbreviate $`v^{}=v_{<k}`$ and $`u^{}=u_{>k+1}`$, where $`k`$ and $`k+1`$ are two positions involved in the $`2`$-move that turns $`𝐢`$ into $`𝐢^{}`$. It is clear from the definitions that $`P_1^{}=P_1`$ and $`P_2^{}=P_2`$ unless $`i_k=i_{k+1}=j`$ for some $`j\mathrm{\Pi }`$ such that $`a_{ji}<0`$. In the latter case, we have
$$P_1^{}=P_1\left(\frac{\mathrm{\Delta }_{v^{}s_j\omega _j,u^{}s_j\omega _j}}{\mathrm{\Delta }_{v^{}\omega _j,u^{}s_j\omega _j}^{\delta _1}\mathrm{\Delta }_{v^{}s_j\omega _j,u^{}\omega _j}^{\delta _2}}\right)^{a_{ji}},P_2^{}=P_2\left(\frac{\mathrm{\Delta }_{v^{}\omega _j,u^{}s_j\omega _j}^{1\delta _1}\mathrm{\Delta }_{v^{}s_j\omega _j,u^{}\omega _j}^{1\delta _2}}{\mathrm{\Delta }_{v^{}\omega _j,u^{}\omega _j}}\right)^{a_{ji}},$$
where $`\delta _1=1`$ (resp. $`\delta _2=0`$) if $`k^{}>1`$ and $`i_k^{}=j`$ (resp. $`(k+1)^+<m`$ and $`i_{(k+1)^+}=j`$), otherwise $`\delta _1=0`$ (resp. $`\delta _2=1`$). Since the common denominator $`(\mathrm{\Delta }_{v^{}\omega _j,u^{}s_j\omega _j}^{\delta _1}\mathrm{\Delta }_{v^{}s_j\omega _j,u^{}\omega _j}^{\delta _2}\mathrm{\Delta }_{v^{}\omega _j,u^{}\omega _j})^{a_{ji}}`$ of $`P_1^{}`$ and $`P_2^{}`$ is relatively prime with $`\mathrm{\Delta }_{s_i\omega _i,\omega _i}`$, it remains to show that
$$P_1(\mathrm{\Delta }_{v^{}\omega _j,u^{}\omega _j}\mathrm{\Delta }_{v^{}s_j\omega _j,u^{}s_j\omega _j})^{a_{ji}}+P_2(\mathrm{\Delta }_{v^{}\omega _j,u^{}s_j\omega _j}\mathrm{\Delta }_{v^{}s_j\omega _j,u^{}\omega _j})^{a_{ji}}$$
is divisible by $`\mathrm{\Delta }_{s_i\omega _i,\omega _i}`$. Since $`P_1+P_2`$ is divisible by $`\mathrm{\Delta }_{s_i\omega _i,\omega _i}`$, it suffices to show that
$$(\mathrm{\Delta }_{v^{}\omega _j,u^{}\omega _j}\mathrm{\Delta }_{v^{}s_j\omega _j,u^{}s_j\omega _j})^{a_{ji}}(\mathrm{\Delta }_{v^{}\omega _j,u^{}s_j\omega _j}\mathrm{\Delta }_{v^{}s_j\omega _j,u^{}\omega _j})^{a_{ji}}$$
is divisible by $`\mathrm{\Delta }_{s_i\omega _i,\omega _i}`$. This in turn follows from the fact that
$$\mathrm{\Delta }_{v^{}\omega _j,u^{}\omega _j}\mathrm{\Delta }_{v^{}s_j\omega _j,u^{}s_j\omega _j}\mathrm{\Delta }_{v^{}\omega _j,u^{}s_j\omega _j}\mathrm{\Delta }_{v^{}s_j\omega _j,u^{}\omega _j}$$
is divisible by $`\mathrm{\Delta }_{s_i\omega _i,\omega _i}`$. But the last expression can be factored according to (4.12), and one of the factors is $`\mathrm{\Delta }_{s_i\omega _i,\omega _i}^{a_{ij}}`$. This completes the proofs of Theorem 4.3 and property (3) in Lemma 3.1.
### 4.4. Property (4)
Let us fix a reduced word $`𝐢=(i_1,\mathrm{},i_m)R(u,v)`$, and an $`𝐢`$-bounded index $`n[2,m]`$. Let $`|i_n|=i\mathrm{\Pi }`$. Let $`^m`$ denote the $`m`$-dimensional vector space with coordinates $`M_1,\mathrm{},M_{n1},M_n^{},M_{n+1},\mathrm{},M_m`$. Let $`t_1,\mathrm{},t_m`$ be rational functions on $`^m`$ given by (4.8) and (4.9), where $`M_n`$ is determined from (3.1). By Theorem 4.1, the map
$$\pi :(M_1,\mathrm{},M_{n1},M_n^{},M_{n+1},\mathrm{},M_m)x_𝐢(t_1,\mathrm{},t_m)$$
is a birational isomorphism $`^mL^{u,v}`$ inverse to the map
$$x(M_1(x),\mathrm{},M_{n1}(x),M_n^{}(x),M_{n+1}(x),\mathrm{},M_m(x)).$$
To prove property (4) in Lemma 3.1, it suffices to show that $`\pi `$ restricts to a regular map $`_0^mL^{u,v}`$.
Let us first show that $`\pi `$ restricts to a regular map $`_0^mG`$. In view of (4.8) and (4.9), if $`k<n^{}`$ or $`k>n`$ then $`t_k`$ is a Laurent monomial in the variables $`M_l`$ with $`ln`$. Thus we only need to show that the product $`x_{i_n^{}}(t_n^{})\mathrm{}x_{i_n}(t_n)`$ is a regular function on $`_0^m`$. Without loss of generality, we can assume that $`n^{}=1`$. For each $`k=2,\mathrm{},n`$, we define $`p_k`$ and a rational map $`y_k:^mG`$ as follows:
$$p_k=t_1\underset{1<l<k}{}_{\epsilon (i_l)=1}t_l^{a_{|i_l|,i}},$$
while $`y_k=x_{i_1}(p_k)x_{i_k}(t_k)x_{i_1}(p_{k+1})^1`$ for $`k<n`$, and $`y_n=x_{i_1}(p_n)x_{i_n}(t_n)`$. Then we have $`x_{i_1}(t_1)\mathrm{}x_{i_n}(t_n)=y_2\mathrm{}y_n`$. Thus it suffices to show that each $`y_k`$ is a regular function on $`_0^m`$. As in section 4.3, we denote $`E_\pm =\{l[2,n1]:\epsilon (i_l)=\pm 1\}`$. Let us first prove that $`y_n`$ is a regular function on $`_0^m`$. We have four cases to consider.
Case 1: $`i_1=i_n=i`$. Then
$$y_n=x_i(p_n)x_i(t_n)=\phi _i(\begin{array}{cc}1& p_n\\ 0& 1\end{array})\phi _i(\begin{array}{cc}1& t_n\\ 0& 1\end{array})=\phi _i(\begin{array}{cc}1& p_n+t_n\\ 0& 1\end{array}).$$
It remains to show that $`p_n+t_n`$ is a regular function on $`_0^m`$. This follows by a direct calculation using (4.8), (4.9) and (3.1): we obtain that
$$p_n+t_n=\frac{M_n^{}}{M_1M_{n^+}}\underset{l>n}{}_{l^{}E_{}\{0\}}M_l^{C_{ln}}\underset{lE_{}}{}_{l^{}E_+}M_l^{C_{ln}}$$
is a Laurent monomial in $`M_1,\mathrm{},M_{n1},M_n^{},M_{n+1},\mathrm{},M_m`$.
Case 2: $`i_1=i,i_n=i`$. Then
$$y_n=x_i(p_n)x_i(t_n)=\phi _i(\begin{array}{cc}1& p_n\\ 0& 1\end{array})\phi _i(\begin{array}{cc}t_n^1& 0\\ 1& t_n\end{array})=\phi _i(\begin{array}{cc}p_n+t_n^1& p_nt_n\\ 1& t_n\end{array}).$$
It remains to show that each of $`t_n`$, $`p_nt_n`$, and $`p_n+t_n^1`$ is a regular function on $`_0^m`$. By a direct calculation, $`p_nt_n`$ and $`p_n+t_n^1`$ are Laurent monomials in $`M_1,\mathrm{},M_{n1},M_n^{},M_{n+1},\mathrm{},M_m`$, while $`t_n=M_n/M_{n^+}`$ is the sum of two such Laurent monomials; in fact, we have
$$p_n+t_n^1=\frac{M_n^{}}{M_1}\underset{lE_{}}{}_{l^{}E_+}M_l^{C_{ln}}.$$
Case 3: $`i_1=i,i_n=i`$. Then
$$y_n=x_i(p_n)x_i(t_n)=\phi _i(\begin{array}{cc}p_n^1& 0\\ 1& p_n\end{array})\phi _i(\begin{array}{cc}1& t_n\\ 0& 1\end{array})=\phi _i(\begin{array}{cc}p_n^1& p_n^1t_n\\ 1& p_n+t_n\end{array}).$$
It remains to show that each of $`p_n^1`$, $`p_n^1t_n`$, and $`p_n+t_n`$ is a regular function on $`_0^m`$. By a direct calculation, $`p_n^1t_n`$ and $`p_n+t_n`$ are Laurent monomials in $`M_1,\mathrm{},M_{n1},M_n^{},M_{n+1},\mathrm{},M_m`$, while $`p_n^1`$ is the sum of two such Laurent monomials; in fact, we have
$$p_n+t_n=\frac{M_n^{}}{M_{n^+}}\underset{l>n}{}_{l^{}E_{}}M_l^{C_{ln}}.$$
Case 4: $`i_1=i_n=i`$. Then
$$y_n=x_i(p_n)x_i(t_n)=\phi _i(\begin{array}{cc}p_n^1& 0\\ 1& p_n\end{array})\phi _i(\begin{array}{cc}t_n^1& 0\\ 1& t_n\end{array})=\phi _i(\begin{array}{cc}p_n^1t_n^1& 0\\ p_n+t_n^1& p_nt_n\end{array}).$$
It remains to show that each of $`(p_nt_n)^{\pm 1}`$ and $`p_n+t_n^1`$ is a regular function on $`_0^m`$. By a direct calculation, both $`p_nt_n`$ and $`p_n+t_n^1`$ are Laurent monomials in $`M_1,\mathrm{},M_{n1},M_n^{},M_{n+1},\mathrm{},M_m`$; in fact, we have
$$p_n+t_n^1=M_n^{}\underset{lE_{}}{}_{l^{}E_+\{0\}}M_l^{C_{ln}}.$$
Now let $`1<k<n`$, and suppose $`|i_k|=j\mathrm{\Pi }`$. To show that $`y_k`$ is a regular function on $`_0^m`$, we have to consider another four cases.
Case 1: $`i_1=i,i_k=j`$. Then $`y_k=x_i(p_k)x_j(t_k)x_i(p_k)`$. Let $`p=p_k^1,q=p_k^{a_{ij}}t_k`$. Clearly, both $`p`$ and $`q`$ are regular functions on $`_0^m`$. The desired regularity of $`y_k`$ becomes a consequence of the following lemma.
###### Lemma 4.4.
For any two distinct $`i,j\mathrm{\Pi }`$, the map $`_0\times N`$ given by $`(p,q)x_i(p^1)x_j(p^{a_{ij}}q)x_i(p^1)`$ extends to a regular map $`^2N`$.
In order not to interrupt the exposition, we will prove this lemma in the end of this section.
Case 2: $`i_1=i,i_k=j`$. Then
$$y_k=x_i(p_k)x_j(t_k)x_i(p_kt_k^{a_{ji}})^1=x_j(t_k)=x_j(M_k/M_{k^+})$$
(see \[4, Proposition 7.2\]), which is a regular function on $`_0^m`$.
Case 3: $`i_1=i,i_k=j`$. Then
$$y_k=x_i(p_k)x_j(t_k)x_i(p_k)^1=x_j(p_k^{a_{ij}}t_k)$$
(see \[4, Proposition 7.2\]), which is a regular function on $`_0^m`$ since $`p_k^{a_{ij}}t_k`$ is a Laurent monomial in $`M_1,\mathrm{},M_{n1},M_{n+1},\mathrm{},M_m`$.
Case 4: $`i_1=i,i_k=j`$. Then
$$y_k=x_i(p_k)x_j(t_k)x_i(p_kt_k^{a_{ji}})^1.$$
Using (4.6) and the commutation relation \[6, (2.5)\], we can rewrite $`y_k`$ as follows:
$$y_k=y_i(p_k)y_j(p_k^{a_{ij}}t_k)y_i(p_k)t_k^{s_i\alpha _j^{}}.$$
The “Cartan factor” $`t_k^{s_i\alpha _j^{}}`$ is clearly a regular function on $`_0^m`$. As for the first factor $`y_i(p_k)y_j(p_k^{a_{ij}}t_k)y_i(p_k)`$, after applying the automorphism $`xx^\iota T`$ of $`G`$, it becomes $`x_i(p_k)x_j(p_k^{a_{ij}}t_k)x_i(p_k)`$, and its regularity follows from Lemma 4.4 with $`p=p_k^1`$ and $`q=t_k`$.
We have proved (modulo Lemma 4.4) that the map
$$\pi :(M_1,\mathrm{},M_{n1},M_n^{},M_{n+1},\mathrm{},M_m)x_𝐢(t_1,\mathrm{},t_m)$$
is a regular map $`_0^mG`$. To complete the proof of property (4), it remains to show that the image of $`\pi `$ is contained in $`L^{u,v}`$. By Theorem 4.1, this image is contained in the closure of $`L^{u,v}`$. Recall from that $`L^{u,v}`$ is determined inside its closure by the conditions $`\mathrm{\Delta }_{\omega _j,v^1\omega _j}(x)0`$ for all $`j`$. The results in also imply that, for any $`xL^{u,v}`$, we have $`\mathrm{\Delta }_{\omega _j,v^1\omega _j}(x)=(\mathrm{\Delta }_{\omega _j,u^1\omega _j}(\psi ^{u,v}(x)))^1=M_{k(j)}(x)^1`$, where $`k(j)`$ is the first occurrence of the index $`\pm j`$ in $`𝐢`$. It follows that
$$\mathrm{\Delta }_{\omega _j,v^1\omega _j}(\pi (M_1,\mathrm{},M_{n1},M_n^{},M_{n+1},\mathrm{},M_m))=M_{k(j)}^10$$
on $`_0^m`$, and we are done.
Proof of Lemma 4.4. The proof below was suggested by N. Reshetikhin; it is simpler than the author’s original proof. Let $`\{e_i:i\mathrm{\Pi }\}`$ be the standard generators of the Lie algebra $`𝔫=\mathrm{Lie}(N)`$; thus, we have $`x_i(t)=\mathrm{exp}(te_i)`$ for every $`i\mathrm{\Pi }`$ and $`t`$. Let $`\mathrm{Ad}:G\mathrm{Aut}(𝔤)`$ be the adjoint representation of $`G`$, and $`\mathrm{ad}:𝔤\mathrm{End}(𝔤)`$ be the differential of $`\mathrm{Ad}`$; recall that these representations satisfy $`\mathrm{exp}(\mathrm{Ad}(x)e)=x\mathrm{exp}(e)x^1`$ and $`\mathrm{exp}(\mathrm{ad}(e))=\mathrm{Ad}(\mathrm{exp}(e))`$ for $`xG`$ and $`e𝔤`$. It follows that
$$x_i(p^1)x_j(p^{a_{ij}}q)x_i(p^1)=\mathrm{exp}(\mathrm{Ad}(x_i(p^1))p^{a_{ij}}qe_j)$$
$$=\mathrm{exp}(\mathrm{exp}(\mathrm{ad}(p^1e_i))p^{a_{ij}}qe_j)=\mathrm{exp}\left(q\underset{n0}{}\frac{p^{a_{ij}n}}{n!}\mathrm{ad}(e_i)^n(e_j)\right)$$
$$=\mathrm{exp}\left(q\underset{n=0}{\overset{a_{ij}}{}}\frac{p^{a_{ij}n}}{n!}\mathrm{ad}(e_i)^n(e_j)\right),$$
which is obviously regular in $`p`$ and $`q`$ (the last equality follows from Serre’s relation $`\mathrm{ad}(e_i)^{1a_{ij}}(e_j)=0`$).
Property (4) in Lemma 3.1, and Theorems 3.2 and 2.2 are finally proved.
## 5. Some examples and applications
### 5.1. Cones of regular monomials
Let us again fix a pair $`(u,v)W\times W`$, and a reduced word $`𝐢=(i_1,\mathrm{},i_m)R(u,v)`$. In view of Theorem 4.1, a generic element $`xL^{u,v}`$ has the form $`x=x_𝐢(t_1,\mathrm{},t_m)`$, so the factorization parameters $`t_k`$ are well-defined rational functions on $`L^{u,v}`$ given by (4.8) and (4.9). Combining these formulas with Lemma 3.3 yields the following corollary.
###### Proposition 5.1.
A Laurent monomial $`t_1^{a_1}\mathrm{}t_m^{a_m}`$ is a regular function on $`L^{u,v}`$ if and only if
(5.1)
$$\epsilon (i_n)a_na_n^{}+\underset{n^{}<k<n}{}_{\epsilon (i_k)=1}C_{nk}a_k0$$
for any $`𝐢`$-bounded $`n[1,m]`$.
Two special cases are worth mentioning. If $`v=e`$ then $`\epsilon (i_k)=1`$ for all $`k`$, and the inequalities (5.1) take the form $`a_na_n^{}`$. If $`u=e`$ then $`\epsilon (i_k)=1`$ for all $`k`$, and the inequalities (5.1) take the form
$$a_na_n^{}+\underset{n^{}<k<n}{}C_{nk}a_k0;$$
the cone defined by these inequalities appeared in a different context in , and also in .
### 5.2. Intersections of two open opposite Schubert cells
Let us illustrate Theorem 2.2 by the case when $`u=e`$ and $`w=w_0`$, the longest element in $`W`$. In this case, $`L^{u,v}`$ is biregularly isomorphic to the intersection of two open opposite Schubert cells $`C_{w_0}w_0C_{w_0}`$, where $`C_{w_0}=(Bw_0B)/B`$ is the open Schubert cell in the flag variety $`G/B`$. These opposite cells appeared in the literature in various contexts, and were studied (in various degrees of generality) in . Let $`C`$ denote the number of connected components of $`L^{e,w_0}()`$; to emphasize the dependency on $`G`$, we shall write $`C=C(X_r)`$, where $`X_r=A_r,B_r,\mathrm{},G_2`$ is the type of $`G`$ in the Cartan-Killing classification.
The numbers $`C(A_r)`$ were determined in : it turns out that $`C(A_1)=2,C(A_2)=6,C(A_3)=20,C(A_4)=52`$, and $`C(A_r)=32^r`$ for $`r5`$. Theorem 2.2 allows us to extend this result to all other simply-laced types.
###### Proposition 5.2.
If $`X_r`$ is one of the types $`A_r(r5),D_r(r4),E_6,E_7`$, or $`E_8`$ then $`C(X_r)=32^r`$.
###### Proof.
Following \[15, Definition 3.10\], we say that a graph is $`E_6`$-compatible if it is connected, and it contains an induced subgraph with $`6`$ vertices isomorphic to the Dynkin graph $`E_6`$ (see Fig. 1).
Combining Theorem 2.2 with \[15, Corollary 3.12\], we obtain the following sufficient condition for the equality $`C(X_r)=32^r`$: it holds provided $`G`$ is simply-laced, and there exists $`𝐢R(w_0)`$ such that the induced subgraph of $`\mathrm{\Sigma }(𝐢)`$ (see Definition 2.1) on the set of all $`𝐢`$-bounded vertices is $`E_6`$-compatible. In this condition was checked for the type $`A_5`$. Therefore, it also holds for any simply-laced Dynkin graph that contains an induced subgraph of type $`A_5`$, that is, for $`A_r(r5),D_r(r6),E_6,E_7`$, and $`E_8`$. It remains to check this condition for the type $`D_4`$ (the statement for $`D_5`$ then follows). Let $`\mathrm{\Pi }=\{1,2,3,4\}`$ with the branching vertex $`3`$. Take the reduced word $`𝐢=(1,2,3,1,2,3,4,3,1,2,3,4)R(w_0)`$. By inspection, the induced subgraph of $`\mathrm{\Sigma }(𝐢)`$ with $`𝐢`$-bounded vertices $`4,5,9,10,11`$, and $`12`$ is isomorphic to the Dynkin graph $`E_6`$, and we are done. ∎
The numbers $`C(B_2)`$ and $`C(G_2)`$ were determined in : it turns out that $`C(B_2)=8`$ and $`C(G_2)=11`$. Theorem 2.2 gives a simpler way to prove these answers. In the case of $`B_2`$, take $`𝐢=(j,i,j,i)`$ with $`a_{ij}=2`$ and $`a_{ji}=1`$. Then $`\mathrm{\Gamma }_𝐢(𝔽_2)`$ is the group of transformations of $`𝔽_2^4`$ generated by $`\tau _3:\xi _3\xi _3+\xi _1`$ and $`\tau _4:\xi _4\xi _4+\xi _3+\xi _2`$. It is easy to see that the action of $`\mathrm{\Gamma }_𝐢(𝔽_2)`$ in $`𝔽_2^4`$ has $`8`$ orbits: four fixed points $`0000,0001,0110`$, and $`0111`$, two $`2`$-element orbits $`0010\stackrel{\tau _4}{}0011`$ and $`0100\stackrel{\tau _4}{}0101`$, and two $`4`$-element orbits $`1000\stackrel{\tau _3}{}1010\stackrel{\tau _4}{}1011\stackrel{\tau _3}{}1001`$ and $`1110\stackrel{\tau _3}{}1100\stackrel{\tau _4}{}1101\stackrel{\tau _3}{}1111`$.
The case of $`G_2`$ is treated in a similar fashion. Take $`𝐢=(j,i,j,i,j,i)`$ with $`a_{ij}=3`$ and $`a_{ji}=1`$. Then $`\mathrm{\Gamma }_𝐢(𝔽_2)`$ is the group of transformations of $`𝔽_2^6`$ generated by the transvections $`\tau _n(3n6)`$ acting by $`\tau _n:\xi _n_{|kn|2}\xi _k`$. It is easy to see that the action of $`\mathrm{\Gamma }_𝐢(𝔽_2)`$ in $`𝔽_2^6`$ has $`11`$ orbits: four fixed points $`000000,001001,001110`$, and $`000111`$; six $`8`$-element orbits, and one $`12`$-element orbit. The $`8`$-element orbits are depicted in Fig. 2 (one has to take the first depicted orbit together with its translates by the $`3`$ non-zero fixed vectors; and the second depicted orbit together with its translate by the vector $`001110`$); the $`12`$-element orbit is depicted in Fig. 3.
###### Remark 5.3.
Computing the numbers $`C(B_r)`$ and $`C(C_r)`$ for $`r3`$ seems to be a challenging problem. Since the transvections $`\tau _n`$ are no longer symplectic in this case, one cannot use \[15, Corollary 3.12\] (at least, not in a straightforward way).
### 5.3. Dual canonical basis for the type $`B_2`$
In conclusion, we briefly discuss a potential application of the above results. Let $`G/N`$ be the *base affine space* for $`G`$. It is well-known that the ring of regular functions $`[G/N]`$ (that is, regular functions on $`G`$ invariant under right translations by elements of $`N`$) is the multiplicity-free sum of all irreducible finite-dimensional representations of $`G`$. Let $`B`$ denote the *dual canonical basis* in $`[G/N]`$ (more precisely, $`B`$ is the “classical limit” of the dual canonical basis in the $`q`$-deformed ring $`_q[G/N]`$). Despite much progress in studying properties of the canonical bases, an explicit construction of $`B`$ still remains to be found. It is known that $`B`$ contains all “Plücker coordinates” $`P_\gamma =\mathrm{\Delta }_{\gamma ,\omega _i}`$ for $`i\mathrm{\Pi }`$ and $`\gamma W\omega _i`$. We suspect that $`B`$ also contains all functions $`\mathrm{\Delta }_𝐢^{}`$ in Theorem 4.3 corresponding to reduced words $`𝐢`$ consisting of elements of $`\mathrm{\Pi }`$. Thus, these functions together with the Plücker coordinates $`P_\gamma `$ are among the building blocks for $`B`$.
As an illustration, consider the case when $`G`$ is of type $`B_2`$, i.e., $`\mathrm{\Pi }=\{i,j\}`$ with $`a_{ij}=2`$ and $`a_{ji}=1`$. The basis $`B`$ in this case was found in (even before the “official” discovery of canonical bases). Translating the results in into our present notation, we obtain the following.
There are $`8`$ Plücker coordinates: $`P_{\omega _i},P_{\omega _j},P_{s_i\omega _i},P_{s_j\omega _j},P_{s_js_i\omega _i},P_{s_is_j\omega _j},P_{w_0\omega _i}`$, and $`P_{w_0\omega _j}`$. Let us also denote $`Q_{\omega _j}=\mathrm{\Delta }_{(i,j,i)}^{}`$ and $`Q_{2\omega _i}=\mathrm{\Delta }_{(j,i,j)}^{}`$; thus, these functions are defined from the equations
(5.2)
$$Q_{\omega _j}P_{s_i\omega _i}=P_{s_is_j\omega _j}P_{\omega _i}+P_{\omega _j}P_{w_0\omega _i}$$
and
(5.3)
$$Q_{2\omega _i}P_{s_j\omega _j}=P_{s_js_i\omega _i}^2P_{\omega _j}+P_{\omega _i}^2P_{w_0\omega _j}.$$
The main result of can be now summarized as follows.
###### Proposition 5.4.
The dual canonical basis $`B`$ of $`[G/N]`$ consists of all monomials in $`10`$ variables $`P_{\omega _i},\mathrm{},P_{w_0\omega _j},Q_{\omega _j},Q_{2\omega _i}`$ with the following property: if this monomial contains variables in two vertices of the “magical hexagon” in Fig. 4 then these two vertices are adjacent.
We see that $`B`$ is the union (not disjoint!) of six families of elements corresponding to the edges of the hexagon in Fig. 4: each family consists of all monomials in six variables $`P_{\omega _i},P_{\omega _j},P_{w_0\omega _i},P_{w_0\omega _j},P,Q`$, where $`P`$ and $`Q`$ lie in two adjacent vertices of the hexagon.
Note that the equations (5.2) and (5.3) can be now interpreted as expansions in the basis $`B`$ of two “forbidden” monomials corresponding to diagonals of the hexagon. There are 7 more such identities corresponding to the remaining 7 diagonals:
$$Q_{\omega _j}P_{s_js_i\omega _i}=P_{w_0\omega _j}P_{\omega _i}+P_{s_j\omega _j}P_{w_0\omega _i};$$
$$Q_{2\omega _i}P_{s_is_j\omega _j}=P_{w_0\omega _i}^2P_{\omega _j}+P_{s_i\omega _i}^2P_{w_0\omega _j};$$
$$P_{s_i\omega _i}P_{s_js_i\omega _i}=P_{\omega _i}P_{w_0\omega _i}+Q_{2\omega _i};$$
$$P_{s_j\omega _j}P_{s_is_j\omega _j}=P_{\omega _j}P_{w_0\omega _j}+Q_{\omega _j}^2;$$
$$P_{s_i\omega _i}P_{s_j\omega _j}=P_{s_js_i\omega _i}P_{\omega _j}+P_{\omega _i}Q_{\omega _j};$$
$$P_{s_js_i\omega _i}P_{s_is_j\omega _j}=P_{s_i\omega _i}P_{w_0\omega _j}+P_{w_0\omega _i}Q_{\omega _j};$$
$$Q_{2\omega _i}Q_{\omega _j}=P_{s_js_i\omega _i}P_{w_0\omega _i}P_{\omega _j}+P_{s_i\omega _i}P_{\omega _i}P_{w_0\omega _j}.$$
## Acknowledgements
I thank Arkady Berenstein, Nicolai Reshetikhin, Boris and Misha Shapiro, and Alek Vainshtein for helpful discussions.
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# Beyond storage capacity in a single model neuron: Continuous replica symmetry breaking
## I Introduction
In his seminal paper Hopfield reformulated Little’s model of auto-associative memory in terms of an energy function. By this act, the field of the statistical mechanics of neural networks was plowed and sown in, and proved itself since then remarkably fertile. The network is an interconnected set of McCulloch-Pitts neurons , the latter being perhaps the biologically least realistic model of a nerve cell. In its simplest version the model neuron can be in one of only two states, “firing” or “quiescent”, it is nonlinear, and admits a large number of connections from other units. Despite the oversimplification in its node cells, the network can exhibit complex behavior and functions as a reasonable model of content addressable memory. From the practical viewpoint, however, artificial neural networks are generically not superior to other methods in the task of storage and associative retrieval .
An important ingredient of this type of models of neural networks is what is called in statistical mechanics quenched disorder. For example, in the Little-Hopfield model the synaptic coupling strengths, made up of random numbers by the Hebb rule , can be held fixed for the duration of the neural dynamical process. This is the basis for the analogy between spin glass models, the archetypical example of a many-body problem with quenched disorder, and neural networks . Methods borrowed from the theory of spin glasses, in particular from the infinite range interaction Sherrington-Kirkpatrick model, yielded a harvest of results on a variety of neural model systems, as well as on other problems whose motivation came from outside of physics but could be formulated as disordered statistical mechanical systems .
It is safe to say that the technique inherited from spin glass theory and the most widely used in the statistical mechanical approach to neural networks in equilibrium is the replica method . It was first applied thoroughly for infinite range interaction spin glasses, thus it is especially suited for networks where neurons have a large number of connections. In its simplest version, with so called replica symmetry, it can be straightforwardly adopted to many a neural problem .
Another family of artificial neural networks consists of layered feed-forward networks , which accept a number of inputs and for a given set of couplings produce the output as a function of the inputs. Such networks were introduced for the purpose of generalization, i. e., rule extraction, from examples of input-output pairs. Then the a priori unknown target rule is to be reconstructed by adaptive changes in the couplings, called training. Since the introduction of the backpropagation algorithm for that purpose, feed-forward networks found wide usage .
The statistical physics of neural modeling gained an impetus of lasting effect from the work by Gardner and Derrida on the storage problem of a single neuron, called also simple perceptron in that context. This is the simplest version of feedforward networks. While in the Little-Hopfield statistical mechanical system the quenched variables are the couplings and a microstate is a configuration of neural states, the roles in feed-forward networks are reversed. Now the space of synaptic couplings is considered as the configuration space within Boltzmannian thermodynamics and the examples appear as quenched parameters. The error for one example measures the difference between the actual output of the network and the required output. The errors on all training examples add up to form the Hamiltonian, i. e., the cost function, of the statistical mechanical model. Within the canonical statistical mechanical approach the temperature plays its usual roles. On the one hand, it is the Lagrange multiplier associated with a preset value of the error, on the other hand, it is the amplitude of noise if a gradient-descent-like dynamics of the couplings is used so as to reach optimal configuration. The thermodynamical limit is achieved by admitting a large number of adjustable couplings, but for that it is not necessary to have many neurons. The approach of Gardner and Derrida was successful in what is called equilibrium learning, when a Hamiltonian can be associated with the problem. However, statistical physical methods are proving themselves useful also in studying off-equilibrium learning algorithms .
A central quantity of a feed-forward network is its storage capacity, i. e., how many random input-output examples the network can reproduce without error. In terms of the statistical mechanical approach this is in its original formulation a zero temperature problem. The subspace of couplings that reproduce a given set of patterns is called version space, its volume, related to the ground state ($`T=0`$) entropy, vanishes beyond capacity. Since the statistical mechanical solution of the region below capacity of a single neuron by Gardner and Derrida a number of results have been obtained about storage properties of feed-forward networks, see for example Refs. . Nevertheless, if the task is to minimize the number of incorrectly stored examples, beyond capacity the problem has not been solved. Technically this is because below capacity for completely random examples replica symmetry holds, while beyond it no finite replica symmetry breaking scheme yields thermodynamically stable solution .
In the present paper we reconsider the problem of storage of random patterns, technically generalize Parisi’s solution of the Sherrington-Kirkpatrick model, and obtain beyond capacity a phase reminiscent to the frustrated ground state of the Sherrington-Kirkpatrick model. That phase continues for $`T>0`$ into the analog of the low temperature, or Parisi, phase of the spin glass.
Our work is motivated not only by the problem of storing random patterns. Generalization has also been successfully analyzed by statistical mechanical methods, see Ref. for reviews. The storage problem below capacity is analogous to equilibrium learning of a learnable task, where the network is compatible with all possible examples, there is no frustration in either systems. For instance, equilibrium generalization properties of the perceptron when the examples are generated by another one, can be understood for arbitrary number of examples within replica symmetry . However, a given feed-forward network may be unable to reproduce a complicated target function on all possible examples . If the target is unlearnable then the network is presumed to get into a frustrated phase if a sufficiently large number of examples are used. Thus the properties of a network beyond capacity are of foremost interest from the viewpoint of rule extraction as well.
The single neuron may be considered as the hydrogen atom of neural problems and studied for its own interest. It is the unit of the Little-Hopfield network, where the symmetry in the couplings has been given up and the couplings of different neurons are considered as independent. As to a feed-forward network, even if the whole network operates without error, its units may still be strained beyond their individual capacities . Thus the description of a single neuron beyond its storage capacity is of importance also from the viewpoint of networked neurons. Furthermore, a close analogy exists between the behavior of model neurons beyond capacity and the glassy, frustrated, phase of disordered spin systems . Therefore, the understanding of the way a single neuron works may have ramifications beyond the field of artificial neural networks.
We firstly review in Sec. II the statistical mechanics of storage, recall the basic thermodynamical quantities and the formula for the replica free energy of a single neuron. In Sec. III we summarize the main ingredients of Parisi’s approach and obtain the free energy functional that we propose describes the equilibrium problem. This way some background is given to our previous communication , wherein we identified a spin glass phase of Parisi type in the high temperature limit. The recipe for the calculation of expectation values by means of Green functions is explained in Sec. IV, producing among other the formulas for the local stability distribution and the stationarity conditions. Sec. V is devoted to the scaling for low temperatures, enabling us to put the extremization of the free energy functional on a computer, and in the end a simulation is discussed.
## II The storage problem and its replica free energy
The model neuron, or perceptron, under consideration is
$`\xi `$ $`=`$ $`\text{sign}(h),`$ (2)
$`h`$ $`=`$ $`N^{1/2}{\displaystyle _{k=1}^N}J_kS_k,`$ (3)
where $`𝐉`$ is the synaptic coupling vector, $`𝐒`$ the input and $`\xi `$ the output. Patterns are given as input-output data,
$$\{𝐒^\mu ,\xi ^\mu \}_{\mu =1}^M.$$
(4)
In the simplest setup the $`S_k^\mu `$-s are independently drawn from any distribution with unit variance and zero average, and $`\xi ^\mu =\pm 1`$ with equal probability. We introduce the local stability parameter as
$`\mathrm{\Delta }^\mu =h^\mu \xi ^\mu ,`$ (5)
where $`h^\mu `$ is given by (3) with $`S_k^\mu `$. If the neuron generates $`\xi ^\mu `$ in response to $`𝐒^\mu `$, i. e.$`\mathrm{\Delta }^\mu >0`$, then we say that the $`\mu `$-th pattern is stored. If $`\mathrm{\Delta }^\mu `$ is a large positive number then high stability of storage against changes in either the couplings or the inputs can be assumed. Large stability is associated with large basin of attraction in memory networks . Given the ensemble of patterns, the local stability parameter obeys some distribution $`\rho (\mathrm{\Delta })`$ . If the number of patterns $`M`$ is of order $`N`$ then it is useful to introduce the relative number of examples
$$\alpha =M/N,$$
(6)
called also load parameter. Since an overall positive factor of $`𝐉`$ does not change the output, we set the norm of $`𝐉`$ to $`\sqrt{N}`$, expressed by the prior distribution
$$w(𝐉)=C_N\delta \left(N\left|𝐉\right|^2\right).$$
(7)
This is called spherical constraint. The factor $`C_N`$ normalizes $`w(𝐉)`$ to unity, it has no thermodynamical significance besides setting the zero point of the entropy scale. Given the distribution of patterns and the length of $`𝐉`$ it can be easily seen that the normalization in (3) results in $`h`$ values of typically $`O(1)`$.
Storage with minimal error can be formulated as an optimization task by our introducing an error measure. If we treat all patterns in the same way we obtain what is called the equilibrium problem. The associated Hamiltonian, or cost function, is
$$=\underset{\mu =1}{\overset{M}{}}V(\mathrm{\Delta }^\mu ),$$
(8)
where the potential $`V(\mathrm{\Delta }^\mu )`$ gives the error on a single pattern $`𝐒^\mu ,\xi ^\mu `$. Obviously, $`V(y)=0`$ for $`y>0`$ in the original storage problem. One can also impose a bound $`\kappa `$ for the local stability, i. e.$`V(y)`$ is set to be zero for $`y>\kappa `$, if $`\kappa >0`$ this is obviously stricter than the original storage criterion. Generically, $`V(y)`$ should be monotonically decreasing for $`y<\kappa `$. In this paper specific results will be presented on the error counting, or Gardner-Derrida , potential
$$V(y)=\theta (\kappa y),$$
(9)
where $`\theta (y)`$ is the Heaviside function. Given the potential (9) the aim is to minimize the number of patterns whose stability is below the bound $`\kappa `$. In the theoretical framework we shall keep the general form $`V(y)`$.
The partition function of the optimization task is
$$Z=d^NJw(𝐉)\mathrm{exp}\left(\beta \underset{\mu =1}{\overset{M}{}}V(\mathrm{\Delta }^\mu )\right),$$
(10)
with $`\beta =1/T`$. Quenched average, $`\mathrm{}_{\text{qu}}`$, is defined as the mean over the patterns. We deal with the idealized equilibrium of the system, when for large $`N`$ the free energy, the energy, and the entropy are assumed to approach their quenched average. This property of self-averaging was proved rigorously only in special cases, see for example , but it is the widely used basis in studies of the equilibrium thermodynamics of disordered systems .
The replica method consists in writing the mean free energy per coupling as
$$f=\underset{N\mathrm{}}{lim}\frac{\mathrm{ln}Z_{\text{qu}}}{N\beta }=\underset{N\mathrm{}}{lim}\underset{n0}{lim}\frac{1Z^n_{\text{qu}}}{nN\beta }.$$
(11)
Denoting the thermal average with the Boltzmann weight in (10) by $`\mathrm{}_{\text{th}}`$ the mean error per pattern can be written as
$$\epsilon =V(\mathrm{\Delta })_{\text{th}}_{\text{qu}}=\frac{1}{\alpha }\frac{\beta f}{\beta }.$$
(12)
The entropy is
$$s=\beta (\alpha \epsilon f).$$
(13)
For $`T=0`$ and $`\epsilon =0`$ the volume of version space, i. e., the space of couplings that perfectly reproduce the examples is obtained as $`\mathrm{\Omega }=\mathrm{exp}(Ns)`$. In general, $`Ns`$ has the usual meaning of the logarithm of the volume with given error $`\epsilon `$.
Introducing the overlap matrix $`𝖰`$ of synaptic vectors
$$\left[𝖰\right]_{ab}q_{ab}=\frac{1}{N}\underset{k=1}{\overset{N}{}}J_{ak}J_{bk},$$
(14)
where the replica indices $`a,b`$ go from $`1`$ to $`n`$, we can express the free energy as the result of the minimization of the replica free energy
$`f`$ $`=`$ $`\underset{n0}{lim}{\displaystyle \frac{1}{n}}\underset{𝖰}{\mathrm{min}}f(𝖰)`$ (16)
$`f(𝖰)`$ $`=`$ $`f_s(𝖰)+\alpha f_e(𝖰),`$ (17)
$`f_s(𝖰)`$ $`=`$ $`(2\beta )^1\mathrm{ln}\text{det}𝖰,`$ (18)
$`f_e(𝖰)`$ $`=`$ $`{\displaystyle \frac{1}{\beta }}\mathrm{ln}{\displaystyle \frac{d^nxd^ny}{(2\pi )^n}\mathrm{exp}\left(\beta _{a=1}^nV(y_a)+i\mathrm{𝐱𝐲}\frac{1}{2}𝐱𝖰𝐱\right)}.`$ (19)
The subscripts $`s`$ and $`e`$ stand for entropic and energy-like terms. The entropic contribution (18) arises because of the spherical constraint and the definition (14), it is indeed independent of the potential, while (19) depends on it.
A central role is played by the probability density for the local stabilities
$$\rho (\mathrm{\Delta })=\delta \left(\mathrm{\Delta }h^1\xi ^1\right)_{\text{th}}_{\text{qu}},$$
(20)
where (3) is understood. Due to the symmetry of the Hamiltonian with respect to the permutation of patterns we could choose $`\mu =1`$ for convenience. It will turn out to be useful to interpret the integrand in (19) as an effective Boltzmann weight and denote the average over this measure as
$$\mathrm{},$$
(21)
where the $`n0`$ limit is implied. A straightforward replica calculation shows (see Ref. for a pedagogic presentation) that the local stability distribution can be rewritten as
$$\rho (\mathrm{\Delta })=\delta \left(\mathrm{\Delta }y_1\right).$$
(22)
Here the subscript could be any replica index, for convenience we chose $`1`$. Comparison with (20) allows an intuitive interpretation for the replica average $`\mathrm{}`$, namely, this corresponds to the combined average over thermal and quenched fluctuations. From (20,22) it follows that the combined thermal and quenched average in (12) boils down to
$$\epsilon =V\left(y_1\right)=𝑑y\rho (y)V(y).$$
(23)
The free energy (14) was calculated within the replica symmetric ansatz for the error counting potential (9) and the capacity, i. e., the maximal $`\alpha `$ with $`\epsilon =0`$ at $`T=0`$ was determined in . It has been shown that beyond capacity the replica symmetric solution is thermodynamically unstable . One and two step replica symmetry breaking solutions were presented, while Ref. proved that no finite step symmetry breaking ansatz can possibly be thermodynamically stable. We presented in a variational free energy functional without derivation that incorporated continuous replica symmetry breaking, but gave concrete results only in the high temperature, large $`\alpha `$ limit. In what follows we provide some background to the general theory and in the end properties of the ground state ($`T=0`$) beyond capacity will also be described.
## III The Parisi scheme
In this section we show how to evaluate the replica free energy $`f_e(𝖰)`$ of (19) within Parisi’s ansatz. The $`R`$ step replica symmetry breaking form is
$$𝖰=\underset{r=0}{\overset{R+1}{}}\left(q_rq_{r1}\right)𝖴_{m_r}𝖨_{n/m_r},$$
(24)
where $`k`$ subscripts $`k`$-dimensional matrices, $`𝖨_k`$ is the identity operator, all elements of $`𝖴_k`$ equal $`1`$, $``$ marks the direct product, and
$`q_1`$ $`=`$ $`0q_0q_1\mathrm{}q_Rq_{R+1}=1,`$ (26)
$`m_{R+1}`$ $`=`$ $`1m_Rm_{R1}\mathrm{}m_1m_0=n,`$ (27)
where the integer $`m_r`$ is a divisor of $`m_{r1}`$. The $`n0`$ limit can be performed smoothly if instead of $`m_r`$ we use $`x_r=(nm_r)/(n1)`$ for parametrization. Thus for arbitrary $`n>0`$ we have the ordering
$$x_{R+1}=1x_Rx_{R1}\mathrm{}x_1x_0=0.$$
(28)
We consider the $`x_r`$-s fixed along the $`n0`$ limiting process, whence follows the formal $`n`$-dependence of the $`m_r(n)`$-s, and for $`n=0`$ we get $`x_r=m_r(0)`$. The inspection of the first few $`R=0,1,2`$ cases allows, in the spirit of Parisi’s , the generalization of the energy term (19) in the replica free energy to arbitrary $`R`$ as
$`f_e(𝐪,𝐱)`$ $`=`$ $`\underset{n0}{lim}{\displaystyle \frac{1}{n}}f_e(𝖰)`$ (29)
$`=`$ $`{\displaystyle \frac{1}{\beta x_1}}{\displaystyle Dz_0\mathrm{ln}Dz_1}`$ (31)
$`\times \left[{\displaystyle Dz_2\mathrm{}\left[Dz_{R+1}\mathrm{exp}\left\{\beta V\left(\underset{r=0}{\overset{R+1}{}}z_r\sqrt{q_rq_{r1}}\right)\right\}\right]^{\frac{x_R}{x_{R+1}}}\mathrm{}}\right]^{\frac{x_1}{x_2}},`$
where
$$Dz=\frac{dze^{\frac{1}{2}z^2}}{\sqrt{2\pi }}.$$
(32)
This is the analog of Parisi’s formula for the Sherrington-Kirkpatrick model, Eq. (11) in ; a comprehensive derivation will be presented elsewhere . The energy term (19) has become a function of the parameters in (26,28). The evaluation of (29) can be done by iteration,
$`\psi _{r1}(y)`$ $`=`$ $`{\displaystyle Dz\psi _r\left(y+z\sqrt{q_rq_{r1}}\right)^{\frac{x_r}{x_{r+1}}}},`$ (34)
$`\psi _{R+1}(y)`$ $`=`$ $`e^{\beta V(y)},`$ (35)
where $`x_{R+2}=1`$ is understood. Then the sought free energy term is obtained as
$$f_e(𝐪,𝐱)=\frac{1}{\beta x_1}Dz\mathrm{ln}\psi _0\left(z\sqrt{q_0}\right).$$
(36)
where $`𝐪=(q_0,\mathrm{},q_R)`$ and $`𝐱=(x_1,\mathrm{},x_R)`$.
The above iteration can be redressed as a partial differential equation. Parisi’s order parameter function $`x(q)`$ is a concatenation of the $`𝐪`$ and $`𝐱`$ as
$$x(q)=\underset{r=0}{\overset{R}{}}(x_{r+1}x_r)\theta (qq_r),$$
(37)
where $`x_1=0`$. Next we introduce $`\psi (q,y)`$ such that at the discontinuities
$`\psi (q_r^{+0},y)`$ $`=`$ $`\psi _r(y),`$ (39)
$`\psi (q_r^0,y)`$ $`=`$ $`\psi (q_r^{+0},y)^{\frac{x(q_r^0)}{x(q_r^{+0})}},`$ (40)
that is, $`\psi (q,y)^{1/x(q)}`$ is continuous in $`q`$. Furthermore, along the plateau in the open interval $`(q_{r1},q_r)`$
$$\psi (q,y)=Dz\psi (q_r^0,y+z\sqrt{q_rq}).$$
(41)
It is easy to show that the field so defined satisfies the partial differential equation (PDE)
$`_q\psi (q,y)`$ $`=`$ $`{\displaystyle \frac{1}{2}}_y^2\psi (q,y)+{\displaystyle \frac{\dot{x}(q)}{x(q)}}\psi (q,y)\mathrm{ln}\psi (q,y),`$ (43)
$`\psi (1,y)`$ $`=`$ $`e^{\beta V(y)},`$ (44)
which evolves from $`q=1`$ to $`q=0`$. Indeed, along the plateaus $`\dot{x}(q)=0`$ when only the first term on the r. h. s. of (43) remains, thus producing (41). Near jumps of $`x(q)`$ the second term dominates, and at a fixed $`y`$ the resulting ordinary differential equation in the variable $`q`$ is separable. Hence it follows that $`\psi (q,y)^{1/x(q)}`$ is continuous in $`q`$, thus (40) is recovered. An equivalent field can be defined by
$$f(q,y)=\frac{\mathrm{ln}\psi (q,y)}{\beta x(q)},$$
(45)
satisfying for a continuous potential $`V(y)`$ the PDE
$`_qf(q,y)`$ $`=`$ $`{\displaystyle \frac{1}{2}}_y^2f(q,y)+{\displaystyle \frac{1}{2}}\beta x(q)\left(_yf(q,y)\right)^2,`$ (47)
$`f(1,y)`$ $`=`$ $`V(y).`$ (48)
The fact that $`f`$ denotes the free energy (16) as well as the field $`f(q,y)`$ should not cause misunderstandings. Equation (47) with initial condition $`\mathrm{ln}2\mathrm{cosh}\beta y`$ has been discovered by Parisi while studying the Sherrington-Kirkpatrick model.
If the potential $`V(y)`$ is not continuous, the PDE (45) holds only from any $`q^{}<1`$ onward where $`\mathrm{ln}\psi (q^{},y)`$ is continuous in $`y`$. In the generic case such is $`q_R`$, so the evolution along the first plateau from $`1`$ to $`q_R`$, where $`x(q)1`$, is to be done explicitly as
$$f(q_R,y)=\frac{\mathrm{ln}\psi (q_R,y)}{\beta x_{R+1}}=\beta ^1\mathrm{ln}Dze^{\beta V\left(y+z\sqrt{1q_R}\right)},$$
(49)
and for $`q<q_R`$ the PDE (47) with the initial condition (49) can be used. Alternatively, one can introduce the field $`m(q,y)`$ as
$$\psi (q,y)m(q,y)=\frac{_y\psi (q,y)}{\beta x(q)},$$
(50)
which equals $`_yf(q,y)`$ when continuity in $`y`$ holds. Then the PDE (47) should be replaced by
$$_qf(q,y)=\frac{1}{2}_ym(q,y)+\frac{1}{2}\beta xm^2(q,y),$$
(51)
while the initial condition (48) can be kept.
Finally we obtain the sought free energy term (19) as a functional of the order parameter function
$$f_e[x(q)]=\underset{n0}{lim}\frac{1}{n}f_e(𝖰)=f(0,0).$$
(52)
It should be emphasized that the above PDE-s do not require infinite refining of the partition by $`q_r`$-s of the interval $`(0,1)`$. They are valid for discrete as well as continuous replica symmetry breaking schemes, i. e., they admit $`x(q)`$ with steps and plateaus, as well as strictly monotonically increasing continuous segments.
The entropic term (18) can also be cast in the form of the energy term (19) with the substitution
$$e^{\beta V(y)}=\sqrt{2\pi }\delta (y).$$
(53)
If we conceive the Dirac delta as a Gaussian with small variance, the initial condition (48) becomes a quadratic function and the PDE (47) can be solved analytically. The analogue of (52) for the entropic term, after going with the variance to zero, can be cast into
$$f_s[x(q)]=\underset{n0}{lim}\frac{1}{n}f_s(𝖰)=\frac{1}{2\beta }_0^1𝑑q\left[\frac{1}{D(q)}\frac{1}{1q}\right],$$
(54)
where
$$D(q)=_q^1𝑑\overline{q}x(\overline{q}).$$
(55)
The free energy of the neuron is then obtained as
$$f=\begin{array}{c}\text{ }\\ \text{max}\\ \text{x(q)}\end{array}f\left[x(q)\right],$$
(56)
where the free energy functional is
$$f[x(q)]=f_s[x(q)]+\alpha f_e[x(q)],$$
(57)
with $`f_s`$ and $`f_e`$ defined in Eqs. (54,52). It is due to the $`n0`$ limit that the maximization in (56) replaces the minimization by the matrix elements of $`𝖰`$ in (16), see, e. g., Ref. .
## IV Linear response theory, stationarity conditions, and expectation values
The least obvious part of the extremization condition (56) is the variation of $`f(0,0)`$ by $`x(q)`$. This can be calculated from linear response theory for the PDE (47). Moreover, linear response theory yields a technique to calculate replica averages as introduced in (21), essential for the evaluation of physical quantities.
The Green function for the PDE (47) can be introduced formally as
$$𝒢(q_1,y_1;q_2,y_2)=\frac{\delta f(q_1,y_1)}{\delta f(q_2,y_2)},$$
(58)
whence $`𝒢(q_1,y_1;q_2,y_2)=0`$ for $`q_1>q_2`$. The Green function for the Parisi solution of the Sherrington-Kirkpatrick model has been studied in Refs. . In the fore and hind variable pairs the Green function $`𝒢(q_1,y_1;q_2,y_2)`$ satisfies the respective PDE-s
$`_{q_1}𝒢={\displaystyle \frac{1}{2}}_{y_1}^2𝒢+\beta x(q_1)m(q_1,y_1)_{y_1}𝒢\delta (q_1q_2)\delta (y_1y_2),`$ (60)
$`_{q_2}𝒢={\displaystyle \frac{1}{2}}_{y_2}^2𝒢+\beta x(q_2)_{y_2}\left[m(q_2,y_2)𝒢\right]+\delta (q_1q_2)\delta (y_1y_2),`$ (61)
where $`m`$ is given in (50). The first equation without the Dirac delta excitation is the linearization of the PDE (47). The minus sign of the Dirac deltas follows from the fact that (60) evolves towards decreasing “time” $`q`$. The homogeneous part of (61) is obtained from the requirement that
$$𝒢(q_1,y_1;q_3,y_3)=𝑑y_2𝒢(q_1,y_1;q_2,y_2)𝒢(q_2,y_2;q_3,y_3)$$
(62)
does not depend on $`q_2`$, and the plus sign of the inhomogeneous term is due to the fact that evolution goes towards increasing $`q`$. The homogeneous parts of two PDE-s (61) and (61) are called adjoint to each other.
For the sake of simplicity we gave formula (58) for the case of continuous potential $`V(y)`$. If the potential is discontinuous then the definition (58) should and can be appropriately modified when a $`q`$-argument is near $`1`$, but the PDE-s (58) for the Green functions hold as they are.
The significance of the Green function is in that it helps to solve the linear PDE with the source term $`h(q,y)`$
$$_q\vartheta (q,y)=\frac{1}{2}_y^2\vartheta (q,y)+\beta x(q)m(q,y)_y\vartheta (q,y)+h(q,y),$$
(63)
as
$$\vartheta (q,y)=𝑑y_1𝒢_\phi (q,y;1,y_1)\vartheta (1,y_1)_q^1𝑑q_1𝑑y_1𝒢_\phi (q,y;q_1,y_1)h(q_1,y_1).$$
(64)
A prominent role will be played by
$$P(q,y)=𝒢(0,0;q,y),$$
(65)
which solves the PDE (61) with $`q_1=y_1=0`$, i. e., with initial condition
$$P(0,y)=\delta (y).$$
(66)
This function first appeared in the context of the Sherrington-Kirkpatrick model in Ref. . Note that the PDE for $`P(q,y)`$ is in fact a Fokker-Planck equation, producing a nonnegative solution and conserving the norm $`𝑑yP(q,y)1`$ for all $`q`$-s. This suggests the intuitive interpretation of $`P(q,y)`$ as probability density of $`y`$.
Now we are in the position to calculate the variation of the free energy functional (57). As to the energy term (52), by varying the functions $`f(q,y)`$ and $`x(q)`$ in the PDE (47) one obtains (63) with $`\vartheta =\delta f`$, $`\vartheta (1,y)=0`$, and $`h=\frac{1}{2}\beta (_yf)^2\delta x`$. Hence (64) gives at $`q=0,y=0`$ the sought $`\delta f(0,0)/\delta x`$. The variation of $`f_s[x(q)]`$ can be calculated straightforwardly, and, with the notation (65), we arrive at
$$F(q,[x(q)])=\frac{2}{\beta }\frac{\delta f[x(q)]}{\delta x(q)}=\left(_0^q\frac{d\overline{q}}{\beta ^2D(\overline{q})^2}\alpha 𝑑yP(q,y)m(q,y)^2\right).$$
(67)
The stationarity condition in case $`x(q)`$ can be freely varied is thus
$$F(q,[x(q)])=0.$$
(68)
If in an interval $`I`$ the $`x(q)`$ is supposed to have a plateau, we differentiate by the plateau value of $`x(q)`$ to get
$$_I𝑑xF(q,[x(q)])=0.$$
(69)
In isolated points $`q_r`$ where plateaus meet (68) should hold pointwise. This summarizes the stationarity conditions for an arbitrary order parameter function, i. e., arbitrary replica symmetry broken scheme of Parisi type.
We have seen in Sec. II instances when the double, thermal and quenched, average could be replaced by the replica average (21). Replica averages can be calculated by the Green function technique as described below. The procedure can be viewed as the generalization of the groundbreaking results from Refs. , where the local magnetization and some of its moments in the Parisi phase of the Sherrington-Kirkpatrick model were evaluated.
A simple case is when $`A(y_a)`$ is to be calculated for an arbitrary function $`A(y)`$. Because of the symmetry with respect to the permutation of single replica indices we have $`A(y_a)=n^1_{a=1}^nA(y_a)`$. This quantity can be easily evaluated if one replaces in (19) $`V(y)`$ by $`V(y)+\lambda A(y)`$, thus obtains $`f_e(𝖰;\lambda )`$, and calculates its initial slope at $`\lambda =0`$. Reversing the limits $`n0`$ and $`\lambda 0`$ then using the first equality of (29) and Eq. (52) we get
$`A(y_a)`$ $`=`$ $`\underset{n0}{lim}{\displaystyle \frac{1}{n}}{\displaystyle \frac{f_e(𝖰;\lambda )}{\lambda }}|_{\lambda =0}={\displaystyle \frac{f(0,0;\lambda )}{\lambda }}|_{\lambda =0}`$ (70)
$`=`$ $`{\displaystyle 𝑑y\frac{\delta f(0,0)}{\delta f(1,y)}A(y)}={\displaystyle 𝑑yP(1,y)A(y)}.`$ (71)
The third equality comes from the fact that $`\lambda `$ is in the initial condition for $`f(q,y)`$ at $`q=1`$, and the last one comes from the definitions for the Green function (58) and for $`P(q,y)`$ (65). Again, for the sake of brevity we gave the derivation for continuous potential $`V(y)`$, however, the result holds also for discontinuous ones. Immediately follows from (22) the formula for the probability density of the local stabilities
$$\rho (y)=P(1,y),$$
(72)
and thus from (23)
$$\epsilon =𝑑yP(1,y)V(y).$$
(73)
From the practical viewpoint interesting are effective averages of products like
$$A_1(y_{a_1})A_2(y_{a_2})\mathrm{}A_k(y_{a_k}).$$
(74)
One can show by using elementary properties of the Fourier transformation that the average of a product of $`x_a`$-s over the effective Boltzmann weight on the r. h. s. of (19) can be expressed as averages over functions of variables $`y_a`$-s. Thus the knowledge how to evaluate (74) also resolves the problem of averages of polynomials in $`x_a`$-s. The latter quantities are of importance because they appear when the replica free energy is differentiated in terms of $`q_{ab}`$-s.
Here we shall only describe the recipe for calculating (74), details can be found in Ref. . If $`k=2`$ then the average depends on $`q=q_{a_1a_2}`$ and is given by the formula
$`C_{12}(q)`$ $`=`$ $`A_1(y_{a_1})A_2(y_{a_2})`$ (75)
$`=`$ $`{\displaystyle 𝑑y𝑑y_2𝑑y_3P(q,y)𝒢(q,y;1,y_1)𝒢(q,y;1,y_2)A_1(y_1)A_2(y_2)}.`$ (76)
Of such type is $`f_e(𝖰)/q_{a_1a_2}`$ where $`A_1(y)=A_2(y)=im(1,y)`$, see (50) for definition, whence we obtain the second term in $`F(q,[x(q)])`$ given in Eq. (67). This is related to the fact that the stationarity condition can also be obtained by first differentiating the replica free energy (17) by $`q_{ab}`$ and then equating the result to zero. For $`k=3`$ suppose without restricting generality that
$$q=q_{a_1a_3}=q_{a_2a_3}<\overline{q}=q_{a_1a_2}.$$
(77)
Then we have
$`C_{123}(q,\overline{q})`$ $`=`$ $`A_1(y_{a_1})A_2(y_{a_2})A_3(y_{a_3})`$ (78)
$`=`$ $`{\displaystyle 𝑑y𝑑\overline{y}𝑑y_1𝑑y_2𝑑y_3P(q,y)𝒢(q,y;1,y_3)𝒢(q,y;\overline{q},\overline{y})𝒢(\overline{q},\overline{y};1,y_1)𝒢(\overline{q},\overline{y};1,y_2)}`$ (80)
$`\times A_1(y_1)A_2(y_2)A_3(y_3).`$
Special versions of the above formulas, for the case of the second and third moments of the magnetization in the Sherrington-Kirkpatrick model, were worked out in Refs. . The integrals in (78) admit a simple graphic representation as shown on Fig. 1.
The cases considered above are the effective average for $`k=1`$, given by Eq. (70), represented by one line, and for $`k=2`$, as calculated in (76), represented by a fork with one handle and two branches. Averages of more than three functions can be analogously constructed, for a given $`k`$ a graph has $`k+1`$ “legs”. Obviously there are two topologically possible graphs for $`k=4`$, depending on the overlaps $`q_{a_ia_j}`$, and more for larger $`k`$-s.
The ability to calculate $`k=4`$ effective averages allows us to study linear stability of the replica free energy (17) at the stationary $`𝖰`$. Using the results of Ref. on ultrametric matrices we expressed the so called replicon eigenvalues in terms of Green functions. While a general proof of the fact that there are no negative eigenvalues in the case of continuous replica symmetry breaking, i. e., when $`x(q)`$ has a continuously increasing segment, is not available, in the high temperature limit we confirmed the absence of linear instability against replicons whenever we encountered such a stationary state. For any temperatures we recovered analytically the zero eigenvalues, corresponding to Goldstone modes, as well as the lowest order Ward-Takahashi identities predicted by algebra .
The generalization of (74) to non-factorizable functions is straightforward. For example, such functions would simply replace the products of $`A_k`$-s in (76) and (78).
## V Low temperature results
In Ref. we have shown that in the high temperature limit, i. e., for $`\alpha ,\beta \mathrm{}`$ with $`\gamma =\alpha \beta ^2`$ finite, the problem simplifies to the extent that if $`x(q)`$ has a continuously increasing segment, this can be given in a closed analytic form. In that limit the problem becomes equivalent to the spherical, multi-$`p`$-spin interaction spin glass , where a similar observation has been made. Four different phases have been found for the error counting potential (9): for small $`\gamma `$ replica symmetry holds, and for $`|\kappa |<2`$ and large $`\gamma `$ there is a Parisi phase with a single continuously increasing segment of $`x(q)`$ between the trivial plateaus $`x0`$ and $`x1`$. When $`|\kappa |>2`$ there is also a narrow one-step replica symmetry broken regime, and for large enough $`\gamma `$-s equilibrium is characterized by an $`x(q)`$ that is a concatenation of a nontrivial plateau and a continuously increasing segment. While for $`T>0`$ there cannot be error free storage, it is plausible to conceive the replica symmetric regime as the continuation of the $`T=0`$ phase of perfect storage and the symmetry broken phases as the analog of the regime at $`T=0`$ beyond capacity.
In the case of a $`V(y)`$ potential that vanishes for $`y>\kappa `$ the limit of capacity is given for $`\kappa 0`$ at $`T=0`$ by
$$\alpha _c(\kappa )=\left(_{\mathrm{}}^\kappa Dt(\kappa t)^2\right)^1,$$
(81)
as it follows from the replica symmetric solution when $`q1`$. This formula also gives the limit of the de Almeida-Thouless (AT) local stability in the case of the potential (9). For $`\kappa =0`$ one has $`\alpha _c=2`$, and, for increasing $`\kappa `$, $`\alpha _c(\kappa )`$ understandably decreases. As it has been already mentioned, beyond capacity none of the finite-step replica symmetry breaking schemes gives a locally stable equilibrium state . Thus in this regime the order parameter function is no longer of the step-like form of (37), rather it has a continuously increasing part. This makes it necessary to numerically solve the extremization problem (56).
The ground state (T=0) has its special scaling properties. The PDE (47) stays meaningful if for $`q<1`$ the function $`\beta x(q)`$ does not diverge, implying that $`x(q)`$ goes to zero. Given the meaning of $`x(q)`$ as the probability of a $`q`$ being in the interval $`(0,q)`$ , we can say that at $`T=0`$ the overlap $`q`$ is $`1`$ with probability 1 for all $`\alpha >\alpha _c`$. Nevertheless, a physically meaningful order parameter is obtained after scaling by $`\beta `$. Firstly we introduce for $`T>0`$ the parameters of the classic Parisi shape as
$$\begin{array}{cc}x(q)0\hfill & \text{ if }0q<q_{(0)}\hfill \\ \dot{x}(q)>0\text{ and }x_{(0)}<x(q)<x_{(1)}\hfill & \text{ if }q_{(0)}<q<q_{(1)}\hfill \\ x(q)1\hfill & \text{ if }q_{(1)}<q1.\hfill \end{array}$$
(82)
The scaled quantities
$`q(t)`$ $`=`$ $`q_{(1)}\left(q_{(1)}q_{(0)}\right)\left(1(1+q_{(1)})t+q_{(1)}t^2\right),0t1,`$ (84)
$`\xi (t)`$ $`=`$ $`\beta x(q(t))\dot{q}(t),`$ (85)
$`\eta `$ $`=`$ $`\beta \left(1q_{(1)}\right),`$ (86)
$`\mathrm{\Delta }(t)`$ $`=`$ $`\beta D(q(t))={\displaystyle _t^1}\xi (\overline{t})𝑑\overline{t}+\eta ,`$ (87)
are expectedly regular even in the $`T0`$ limit, when $`q_{(1)}1`$. Note that there is an arbitrariness in the parametrization by $`t`$, the main features being that $`q(0)=q_{(0)}`$, $`q(0)=q_{(0)}`$, $`q(1)=q_{(1)}`$, and $`\dot{q}(1)=0`$. With this parametrization the PDE (45) becomes
$`_tf(t,y)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\dot{q}(t)_y^2f(t,y)+{\displaystyle \frac{1}{2}}\xi (t)\left(_yf(t,y)\right)^2,`$ (89)
$`f(1,y)`$ $`=`$ $`{\displaystyle \frac{1}{\beta }}\mathrm{ln}{\displaystyle Dze^{\beta V\left(y+z\sqrt{1q_{(1)}}\right)}}.`$ (90)
For $`T=0`$ the initial condition becomes
$$f(t=1,y)|_{T=0}=\underset{\overline{y}}{\mathrm{min}}\left(V(\overline{y})+\frac{(y\overline{y})^2}{2\eta }\right)=\{\begin{array}{cc}1\hfill & \text{if }y\kappa \sqrt{2\eta }\hfill \\ \frac{(\kappa y)^2}{2\eta }\hfill & \text{if }\kappa \sqrt{2\eta }y\kappa \hfill \\ 0\hfill & \text{if }y\kappa ,\hfill \end{array}$$
(91)
where (9) was substituted to get the second equality. By Gaussian integration hence the replica symmetric solution can be obtained . Note that $`f(t=1,y)`$ is a continuous function even though $`V(y)`$ is the step function. Although we used the same symbols for functions of $`q`$ and $`t`$, misunderstanding are avoided by our marking which argument we mean.
The PDE for $`P(q,y)`$ follows from the definition (65) and from the evolution equation of the Green function (61). The latter can be properly rescaled according to (82), yielding in principle the solution $`P(t,y)`$. The probability density of local stabilities is obtained by evolving $`P(t=1,y)=P(q=q_{(1)},y)`$ to $`P(q=1,y)`$. At $`T=0`$ with (9) we have
$$\rho (\mathrm{\Delta })=\{\begin{array}{cc}P(1,\mathrm{\Delta })\hfill & \text{if }\mathrm{\Delta }\kappa \sqrt{2\eta }\hfill \\ 0\hfill & \text{if }\kappa \sqrt{2\eta }<\mathrm{\Delta }<\kappa \hfill \\ P(1,\mathrm{\Delta })+\delta (\mathrm{\Delta }\kappa )_{\kappa \sqrt{2\eta }}^\kappa 𝑑\overline{y}P(1,\overline{y})\hfill & \text{if }\mathrm{\Delta }\kappa ,\hfill \end{array}$$
(92)
where $`P(1,\mathrm{\Delta })=P(t=1,\mathrm{\Delta })`$ is understood. For arbitrarily small $`T>0`$ the gap in the support of $`\rho (\mathrm{\Delta })`$ immediately vanishes. For details we again refer to .
The numerical extremization was done with complementing the free energy functional (57) by constraints. The PDE (V) was added giving rise to a Lagrange multiplier field. This field can be shown to be just $`P(t,y)`$, see Refs. . Further technical requirements are $`x(q_{(0)})0`$, $`x(q_{(1)})1`$, and $`\dot{x}(q)0`$ for $`q_{(0)}<q<q_{(1)}`$, which were taken into account by soft constraints. A few results for the Gardner-Derrida potential (9), with $`\kappa =0`$, $`\alpha =3`$, are displayed on Figs. 2 and 3 for various low temperatures. Note that the point $`(\kappa =0,\alpha =3)`$ lies beyond capacity, while beyond about $`\beta ^1=T=0.2`$ the RS solution satisfies the AT stability condition .
On Fig. 2 the scaled order parameter function $`\beta x(q)`$ is shown. For small temperatures the replica symmetric solution is AT unstable, and we indeed obtain the Parisi form (82) for $`x(q)`$. At $`q_{(0)}`$ near $`0.75`$ the functions $`x(q)`$ jump to zero and remains there as $`q`$ further decreases. The upper plateau with $`x(q)1`$ starts at $`q_{(1)}`$. The consistency of the scaling (82) is confirmed by our finding at $`T=0`$ finite values for $`\beta x(q)`$ if $`q<1`$ and for $`\eta 0.26`$. Interestingly, the curved segment of $`\beta x(q)`$ does not change much with increasing temperature, the main effect being the decrease of $`q_{(1)}`$.
FIG. 2. Scaled order parameter function $`x(q)`$ for $`\kappa =0`$, $`\alpha =3`$ at $`T=0`$ (solid), $`T=0.01`$ (dashed), and $`T=0.1`$ (dotted).
The local stability distribution $`\rho (\mathrm{\Delta })`$ of (22) is displayed on Fig. 3. It was numerically obtained for $`T=0`$ from (92) and for $`T>0`$ from the original formula (72), with the same parameter values as in Fig. 3. For the sake of better visibility of the other details, the very high peaks near $`\mathrm{\Delta }=\kappa =0`$ for $`T=0.1`$ and $`T=0.01`$ as well as the corresponding $`\delta `$-peak for $`T=0`$ have been omitted from this plot. For $`|\mathrm{\Delta }|>3`$ the curves approach zero very quickly. A true gap with $`\rho (\mathrm{\Delta })=0`$ to the left of $`\mathrm{\Delta }=0`$ develops only if $`T=0`$, while for the positive $`T`$-values $`\rho (\mathrm{\Delta })`$ is positive albeit small there. While the gap at $`T=0`$ is present in $`R=0,1,2`$ step replica symmetry breaking schemes, see Refs. respectively, in all these cases a jump appears near the lower edge. This can be associated with the thermodynamic instability of those saddle points . Our present solution gives linearly vanishing $`\rho (\mathrm{\Delta })`$ at the lower edge, signaling the absence of replicon instability .
FIG. 3. Density of local stabilities $`\rho (\mathrm{\Delta })`$ from theory for $`\kappa =0`$, $`\alpha =3`$ at $`T=0`$ (solid), $`T=0.01`$ (dashed), and $`T=0.1`$ (dotted).
In order to compare theory with practice we performed a medium scale simulation. The standard Hebbian algorithm was modified by Wendemuth to provide convergence for negative stabilities. Since we chose $`\kappa =1`$, the final steps during stabilization of a pattern went on with $`\mathrm{\Delta }>0`$, so in our case the modification was not essential. The algorithm goes as follows. Firstly random patterns (4) are generated uniformly from an interval centered about zero and normalized as $`_{k=1}^N(S_k^\mu )^2=N`$, and all outputs $`\xi ^\mu `$ are taken uniformly $`1`$. This does not restrict generality since $`S_k^\mu `$ have random signs. The initial coupling vector is set to be proportional to
$$J_k(0)\underset{\mu =1}{\overset{M}{}}S_k^\mu ,$$
(93)
such that its Eucledian norm is $`N`$. In the $`t`$-th step of the algorithm one calculates the local stabilities
$$\mathrm{\Delta }^\mu (t)=\frac{𝐉(t)𝐒^\mu }{|𝐉(t)|}$$
(94)
and selects the least unstable pattern, i. e., the one with the largest $`\mathrm{\Delta }^\mu (t)<\kappa `$. Let us denote its index by $`\mu _0`$, whose argument $`t`$ we omit. Next one augments the couplings as follows. If $`\mathrm{\Delta }^{\mu _{0(t)}}(t)>0`$ then
$$𝐉(t+1)=𝐉(t)+\lambda 𝐒^{\mu _0},$$
(95)
and if $`\mathrm{\Delta }^{\mu _0(t)}(t)<0`$ we have following Wendemuth
$$𝐉(t+1)=𝐉(t)+\lambda \left(𝐒^{\mu _0}+𝐉(t)\frac{N/|𝐉(t)|\mathrm{\Delta }^{\mu _0}(t)}{|𝐉(t)|\mathrm{\Delta }^{\mu _0}(t)}\right).$$
(96)
Here $`\lambda `$ is the gain parameter, the overall scale of increments of $`𝐉`$. In Ref. the gain parameter was $`\lambda =N^{3/2}`$, after experimentation we chose $`\lambda =N^1`$. Such an increase in the gain parameter did not endanger, rather sped up convergence. At time $`t+1`$ we again look for the least unstable pattern, and so on. The algorithm goes on until it gets stuck with one pattern that we are not able to stabilize in a reasonable time. The intuitive idea behind the algorithm is that since an unstable pattern counts as error irrespective of the distance of the stability parameter $`\mathrm{\Delta }^\mu `$ from $`\kappa `$, one assumes that it is the easiest to stabilize the pattern with $`\mathrm{\Delta }^\mu `$ closest to $`\kappa `$. So one may hope that thus the largest possible number of patterns can be stabilized.
The program ran on $`28`$ PC-s in parallel, each having an AMD K6 processor of $`333`$ MHz, during about one day. Fig. 4 shows both the theoretical curve and the results of the simulation for $`\alpha =\kappa =1`$, a point known to fall beyond the capacity curve (81). The full line is the result of numerical extremization of the free energy functional in the way Fig. 3 was obtained. The Dirac delta peak of the theoretical probability density at $`\kappa `$ is not illustrated. The discontinuous lines represent the histograms for the local stabilities from simulation for two sizes, $`M=N=500`$ and $`1000`$, after normalization.
FIG. 4. Density of local stabilities $`\rho (\mathrm{\Delta })`$ at $`\alpha =\kappa =1`$, axes as in Fig. 3. The theoretical prediction is given by the full line. The two empirical densities are normalized histograms, taken with $`M=N=500`$ and $`1000`$.
The closeness of the two histograms demonstrates that size effects were probably not the cause for the systematic difference between theory and the numerical experiment. A possible ground for the discrepancy is that the algorithm may have been halted prematurely. However, the time necessary for the stabilization of patterns was allowed to grow for each subsequent pattern, and the algorithm was ended only when stabilization did not occur even within the multiple of such an extrapolated time. Another possible reason for the deviation may be that the algorithm got stuck in a “local optimum” without being able to globally maximize the number of stable patterns. In this regard several modified initial conditions were tested but the number of stabilized patterns did not grow in the end. A source of concern can be that the built in random number generator of the C compiler was used; we did not test other routines for this purpose. As to the algorithm, despite its intuitive appeal, there is no proof that it would be able to globally minimize the Hamiltonian (8) with error measure (9). Furthermore, it is likely that with the present parameters the learning task is an NP-complete problem , thus explaining imperfect convergence.
We emphasize that the results represent a significant improvement with respect to the earlier simulation in Ref. . The error per example $`\epsilon `$ found in is about $`0.21`$, while the present data correspond to $`0.15`$ and theory predicts $`0.1358`$. Thus the deviation between simulation and theory has been decreased by 80%. That means that we stabilized more patterns than , although, given the difference between the theoretical and simulation results, we still could not find the global optimum. Furthermore, an important feature of the density $`\rho (\mathrm{\Delta })`$ is that it should continuously vanish, with a nonzero slope, at the lower edge of the gap. This property is reproduced by the simulation data, in a sharper fashion with the larger $`M=N=1000`$ size, while the value of the edge remains slightly overestimated.
###### Acknowledgements.
The authors acknowledge support by OTKA grant No. T017272 (G. Gy.) and from a special grant for young scientists from the University of Augsburg and by the State of Bavaria within the postgraduate scheme Graduiertenkolleg GRK283 “Nonlinear Problems in Analysis, Geometry, and Physics” (P. R.). It is a pleasure to thank F. Csikor and Z. Fodor for offering us for the simulation their PC farm, supported by grants Nos. OTKA-T22929 and FKFP-0128/1997. Thanks are due to F. Pázmándi for his pointing out and discussing with us the problem of a discontinuity in the potential.
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# DIAGONALIZATION OF A BOSONIC QUADRATIC FORM USING CCM: APPLICATION ON A SYSTEM WITH TWO INTERPENETRATING SQUARE LATTICE ANTIFERROMAGNETS
## 1 Introduction – The Model
It is always possible to diagonalize quadratic bosonic forms (which appear frequently in physics) using a Bogoliubov transformation , but it is a tedious analytical task to find one for a complicated form with many different magnons. Therefore we want to show here how the coupled cluster method (CCM), one of the most powerful and universal techniques in quantum many-body theory (s. and references therein), can be used in a straightforward scheme to find the exact ground state of such a form.
To be concrete we consider the 2D spin $`1/2`$ Heisenberg model
$$H=J_{AA}\underset{iA_1,jA_2}{}𝐒_i𝐒_j+J_{BB}\underset{iB_1,jB_2}{}𝐒_i𝐒_j+J_{AB}\underset{iA,jB}{}𝐒_i𝐒_j,$$
(1)
which is related to the situation in Ba<sub>2</sub>Cu<sub>3</sub>O<sub>4</sub>Cl<sub>2</sub> , a layered quantum antiferromagnet showing significant differences to its parent cuprats like La<sub>2</sub>CuO<sub>4</sub> (see e.g. for recent experiments). In contrast to La<sub>2</sub>CuO<sub>4</sub> we have two different types of Cu-sites in the Cu-O-planes, namely there are additional Cu(B) atoms located in the centre of every second Cu(A)-O<sub>2</sub> square. Within the Cu(A) subsystem we have a strong 180 Cu-O-Cu superexchange yielding to strong antiferromagnetic couplings ($`J_{AA}`$) between Cu(A) atoms, whereas the couplings within the Cu(B) subsystem ($`J_{BB}`$) and between the subsystems ($`J_{AB}`$) are weaker. A recent calculation of $`J_{AA}`$, $`J_{BB}`$ , finding $`J_{AA}10J_{BB}`$ (both antiferromagnetic) agrees with the experimental values . There are also some arguments for a ferromagnetic $`|J_{AB}|J_{BB}`$.
In the classical ground state (1) shows for $`|J_{AB}|2\sqrt{J_{AA}J_{BB}}`$ a Néel like order for the two subsystems A and B, where the energy is degenerated with respect to the angle $`\phi `$ between the spins of these two subsystems.
## 2 The Method
In this paper we study the ground state properties of (1), using a four-magnon linear spin wave approximation around the classical ground state, i.e. for each of the four sublattices $`A_1,A_2,B_1,B_2`$ of the two coupled bipartite antiferromagnetic square lattices we introduce different bosonic operators. Thus we get for (1)
$$H=\frac{2N}{3}s^2\left(2J_{AA}+J_{BB}\right)+\underset{𝐤}{}H_𝐤,\text{with}$$
(2)
$$\begin{array}{c}H_𝐤=4J_{AA}s\left(a_{1𝐤}^+a_{1𝐤}+a_{2𝐤}^+a_{2𝐤}\gamma _{𝐤AA}\left[a_{1𝐤}^+a_{2𝐤}^++a_{1𝐤}a_{2𝐤}\right]\right)\hfill \\ +2J_{BB}s\left(b_{1𝐤}^+b_{1𝐤}+b_{2𝐤}^+b_{2𝐤}\gamma _{𝐤BB}\left[b_{1𝐤}^+b_{2𝐤}^++b_{1𝐤}b_{2𝐤}\right]\right)\hfill \\ +J_{AB}s(1+\mathrm{cos}\phi )/2\left(b_{1𝐤}^+a_{1𝐤}+b_{1𝐤}a_{1𝐤}^+b_{2𝐤}^+a_{1𝐤}^+b_{2𝐤}a_{1𝐤}\right)\gamma _{𝐤AB}^1\hfill \\ +J_{AB}s\left(1\mathrm{cos}\phi \right)/2\left(b_{2𝐤}^+a_{2𝐤}+b_{2𝐤}a_{2𝐤}^+b_{1𝐤}^+a_{2𝐤}^+b_{1𝐤}a_{2𝐤}\right)\gamma _{𝐤AB}^2\hfill \\ +J_{AB}s\left(1\mathrm{cos}\phi \right)/2\left(b_{2𝐤}^+a_{1𝐤}+b_{2𝐤}a_{1𝐤}^+b_{1𝐤}^+a_{1𝐤}^+b_{1𝐤}a_{1𝐤}\right)\gamma _{𝐤AB}^1\hfill \\ +J_{AB}s\left(1+\mathrm{cos}\phi \right)/2\left(b_{1𝐤}^+a_{2𝐤}+b_{1𝐤}a_{2𝐤}^+b_{2𝐤}^+a_{2𝐤}^+b_{2𝐤}a_{2𝐤}\right)\gamma _{𝐤AB}^2,\hfill \end{array}$$
(3)
using the lattice structure factors $`\gamma _{𝐤AA}=\mathrm{cos}(k_x/2)\mathrm{cos}(k_y/2)`$, $`\gamma _{𝐤BB}=(\mathrm{cos}k_x+\mathrm{cos}k_y)/2`$ and $`\gamma _{𝐤AB}^{1(2)}=\mathrm{cos}(k_{x(y)}/2)`$.
As stated, we use the coupled cluster method (CCM) to find the exact ground state of (2). To do this we notice the following property of $`H`$
$$\underset{𝐤}{}H_𝐤=\underset{𝐤}{}(H_𝐤+H_𝐤)/2\underset{𝐤}{}H_𝐤^{};[H_𝐤^{},H_𝐤^{}^{}]_{}=0𝐤,𝐤^{}.$$
(4)
Hence it is possible to treat each $`H_𝐤^{}`$ seperately within the CCM, since they all commute with each other. So we have to deal with a bosonic system with eight different bosonic operators $`a_{1\pm 𝐤}`$, $`a_{2\pm 𝐤}`$, $`b_{1\pm 𝐤}`$, $`b_{2\pm 𝐤}`$ denoted with $`a_1,\mathrm{},a_8`$.
The ket and bra ground state of such a system (i.e a many-mode bosonic field theory with bosonic operators $`a_i`$, $`a_i^+`$ in the Hamiltonian) in CCM-SUB$`l`$ approximation is given by
$$\begin{array}{cc}|\mathrm{\Psi }=e^S|0,\hfill & S=_{i_1,i_2,\mathrm{},i_l}A_{i_1,i_2,\mathrm{},i_l}a_{i_1}^+a_{i_2}^+\mathrm{}a_{i_l}^+,\hfill \\ \stackrel{~}{\mathrm{\Psi }}|=0|\stackrel{~}{S}e^S,\hfill & \stackrel{~}{S}=1+_{i_1,i_2,\mathrm{},i_l}\stackrel{~}{A}_{i_1,i_2,\mathrm{},i_l}a_{i_1}a_{i_2}\mathrm{}a_{i_l}\hfill \end{array},$$
(5)
where $`|0`$ is the bosonic vacuum state (i.e. $`a_i|0=0`$), and $`A_{i_1\mathrm{}}`$ and $`\stackrel{~}{A}_{i_1\mathrm{}}`$ are the CCM correlation coefficients. These coefficients are calculated by two systems of equations (one of them is a system of nonlinear equations).
$$\frac{\overline{H}}{\stackrel{~}{A}_{i_1\mathrm{}i_l}}=0,\frac{\overline{H}}{A_{i_1\mathrm{}i_l}}=0,\overline{H}=\stackrel{~}{\mathrm{\Psi }}|H|\mathrm{\Psi },$$
(6)
using the expectation value ($`\overline{H}`$) of the Hamiltonian, i.e. the ground state energy.
Note, that the CCM-SUB2 approximation (i.e. having only quadratic terms of bosonic operators in $`S`$ and $`\stackrel{~}{S}`$ (5)) gives the exact ground state of a quadratic bosonic Hamiltonian, since the ground state wave function of such a Hamiltonian has the form $`|\mathrm{\Psi }=\mathrm{exp}[_{ij}f_{ij}a_i^+a_j^+]|0`$, which can easily be shown using a Bogoliubov transformation (see appendix). Therefore the CCM correlation operator $`S`$ (and $`\stackrel{~}{S}`$ respectively) (5) consist of products of two bosonic creation operators only, all other coefficients $`A_{i_1,\mathrm{},i_l}`$ are zero; so we just have to use SUB2.
To calculate the CCM equations (6) easily using computer algebra, we make use of the Bargmann representation
$$a^+z,a\frac{d}{dz},|01,0|f(a,a^+)|0f(\frac{d}{dz},z)|_{z=0},$$
(7)
which maps the original many-mode bosonic field theory into the corresponding (classical) field theory of complex functions in a particular normed space. So instead of bosonic operators we just have to handle with (complex) numbers and differential operators, which is much easier. Once the (partial nonlinear) equations are obtained they can be solved numerically.
## 3 Results and conclusions
We apply the CCM-scheme described above to calculate the exact ground state of (3) and by doing this getting a spin wave approximate ground state of the model (1). We discuss the energy as a function of the angle between spins of the two subsystems $`A`$ and $`B`$ and the correlation between spins of different subsystems as a function of $`J_{AB}`$ (Fig.1). We find as a typical order from disorder effect, that the degeneracy of the ground state with respect to the angle $`\phi `$ is lifted by quantum fluctuations and a collinear ordering ($`\phi =0,\pi `$) is stabilized. This can clearly be seen by the energy vs. $`\phi `$ picture in Fig.1 and by the correlation $`\text{S}_i\text{S}_j_{A,B}`$ vs. $`J_{AB}`$, which is zero in the classical case, independent of the value of $`J_{AB}`$ (for $`|J_{AB}|2\sqrt{J_{AA}J_{BB}}`$). In the quantum case however that correlation does depend on $`J_{AB}`$, showing again an order effect induced by quantum fluctuations.
In addition we find a lowering of the magnetic order within the subsystems $`A`$ and particular $`B`$ by frustrating $`J_{AB}`$ in the quantum case.
## Acknowledments
This work has been supported by the DFG (Project Nr. Ri 615/7-1).
## Appendix A Proof that CCM-SUB2 gives exact ground state
Using the fact, that a Bogoliubov transformation $`\beta _\nu =_\mu (u_{\mu \nu }^{}a_\mu v_{\mu \nu }^{}a_\mu ^+)`$ exactly diagonalize a quadratic bosonic Hamiltonian with the bosonic operators $`a_i`$, $`a_i^+`$, one can easily show that its ground state must have the form $`|\mathrm{\Psi }=\mathrm{exp}[_{ij}f_{ij}a_i^+a_j^+]|0`$, by showing that $`\beta _\nu |\mathrm{\Psi }=0`$ $`\nu `$. We use the Bargmann representation (7) and get
$$\beta _\nu |\mathrm{\Psi }\stackrel{!}{=}0\underset{\mu }{}\left(u_{\mu \nu }^{}\frac{d}{dz_\mu }v_{\mu \nu }^{}z_\mu \right)\mathrm{exp}\left[\underset{ij}{}f_{ij}z_iz_j\right]\stackrel{!}{=}0z_i$$
$$\underset{\mu }{}(u_{\mu \nu }^{}2\underset{i}{}f_{i\mu }z_iv_{\mu \nu }^{}z_\mu )\stackrel{!}{=}0z_i,2\underset{\mu }{}f_{i\mu }u_{\mu \nu }^{}\stackrel{!}{=}v_{i\nu }^{}$$
and this last matrix equation is allways fulfilled for some $`f_{i\mu }`$.
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# 1 Introduction
## 1 Introduction
Horizontal air showers have been studied for many years for several different reasons . Much of the relevance of horizontal air showers induced by cosmic rays is in the understanding of the background against which high energy neutrino showers could be detected. Although no neutrino events can be expected in the Haverah Park data set, even for the most optimistic neutrino flux predictions, it has recently been shown that relevant bounds could be obtained with the Auger Observatories .
Horizontal cosmic ray showers are of great interest in their own right for two principle reasons. Firstly the acceptance of an air-shower array could be doubled if events above $`60^{}`$ can be adequately analysed. Secondly because, uniquely, they deal with surviving particles that are created very close to the shower core, they complement the information obtained in the study of near vertical cosmic ray showers. Moreover understanding the azimuthal asymmetries at large zenith angles \[6,7,8,9,10 (hereafter Paper I)\] can lead to significant improvement of ultra high energy shower analysis at moderate zenith angles ($`30^{}60^{}`$) .
The Haverah Park array, being made of 1.2 m deep water-Čerenkov tanks , is quite possibly the array detector so far constructed which is best suited on geometrical considerations for the analysis of very large inclined showers. Moreover it can be considered as an early prototype of the Auger Observatories which will employ water-Čerenkov tanks of identical depth. The quantitative aspects of our results are very specific to the water-Čerenkov technique.
During the 14 years of operation of the Haverah Park array which are considered here nearly 10000 air showers were recorded for $`\theta >60^{}`$. The analysis of these data was difficult because of the complex geomagnetic field effects which distort the circular symmetry of air showers, the relatively small size of the array (12 km<sup>2</sup>) and the limited computing power then available, both for data analysis and simulation. A single remarkable event, perhaps of 10<sup>20</sup> eV, at $`85^{}`$ from the zenith which triggered 20 of the array detectors was extensively discussed , but no systematic study was made with those data for which the initial analysis gave a zenith angle $`\theta >60^{}`$.
Here we present the results of the first systematic analysis of horizontal showers from primaries of energy greater than $`10^{17}`$ eV at Haverah Park with zenith angles exceeding $`60^{}`$, comparing them to those expected from the known cosmic ray spectrum. Preliminary results of this work were reported in . For the calculation of the expected rates we make use of a parameterization for the muon number densities at ground level described elsewhere (Paper I). We take into account two possible compositions (proton and iron nuclei) and use alternative models for the high energy interactions namely QGSJET and SIBYLL 1.6 .
The article is organized as follows: In section 2 we discuss the relative signals and fluctuations expected from the electromagnetic (electrons and photons) and muon components in horizontal air showers (HAS) induced by hadrons and in section 3 we address the muon component of HAS. In section 4 we discuss the detector simulation paying particular attention to the origin of the different parts of the signal from the individual tanks, stressing the differences between vertical and horizontal particle signals. In section 5 we describe in detail the procedure implemented to obtain a prediction for the horizontal shower rate and in section 6 we give the results obtained in the different cases. Our conclusions are presented in section 7.
## 2 Muon and Electromagnetic Components in Horizontal Air Showers
An air shower induced by a proton or a nucleus can be qualitatively understood as a pion shower that is continuously feeding an electromagnetic component (photons, electrons, and positrons) through $`\pi ^0`$ decay, and both a muonic and a neutrino component through charged pion decay. At ground level the typical number density ratio of photons to electrons in a vertical shower is stable. For a low energy threshold of 90 keV this ratio has a value typically between $`20:1`$ and $`30:1`$ for proton or iron as the distance to shower axis $`r`$ changes from 100 m to 1 km. This value is characteristic of electromagnetic showers. The ratio of electrons to muons on the other hand depends strongly on $`r`$. For example for a $`10^{19}`$eV proton shower it drops from $`100:1`$ to $`3:1`$ as $`r`$ rises in the same interval.
In spite of outnumbering the muon component, the average electron and photon energies are typically in the few MeV range, compared to GeV and above for muons. As a result the relative contributions to the signal in a water-Čerenkov tank are not so different because the water-Čerenkov technique is on average much more efficient at converting single muons into signals than electrons or photons.
Direct measurements with Fly’s Eye have shown that above $`10^{17}`$ eV the shower maximum is usually between 600 g cm<sup>-2</sup> and sea level ($`1000`$ g cm$`{}_{}{}^{2})`$. As the zenith angle varies from the vertical, $`\theta =0^{}`$, to the horizontal, $`\theta =90^{}`$, direction, the slant matter depth rises to $`36000`$ g cm<sup>-2</sup> and showers are observed well past maximum. The behavior of the electromagnetic and muon components beyond shower maximum is shown in Fig. 1A. While the electromagnetic component of an air shower becomes exponentially attenuated with depth, the muons which do not decay propagate practically unattenuated to the ground, except for energy loss and geomagnetic deviations. As a result the ratio of the electromagnetic to the muon component of an air shower drops as the zenith angle increases up to $`60^{}`$. Above this angle the ratio levels out because the muons themselves produce electromagnetic particles. The remaining electromagnetic component is mainly due to muon decay and to a smaller extent to hadronic interactions, pair production, and bremsstrahlung.
We have studied the effects of muon decay through Monte Carlo simulation using AIRES . Indeed even at $`60^{}`$ the electromagnetic component of a $`10^{19}`$ eV proton shower which can be directly associated to $`\pi ^0`$ decay is already low and confined within a relatively small region of about 200 m around shower axis. At higher zenith angles this component can be neglected and the muon decay contribution becomes stable, roughly at a proportion of 0.8 electrons for each muon, more than three times below the number obtained for vertical showers 1 km away from the shower axis. Unlike the electromagnetic component from pion decay in vertical showers, the lateral distribution follows that of the muons rather closely. This is not difficult to understand as the decay muons give rise to relatively small electromagnetic sub-showers that preserve the muon spatial distribution. The energy distribution of the electromagnetic component is essentially the same as that of an electromagnetic air shower (i.e. for vertical showers) for the same reason.
The contribution of muon interactions to the electromagnetic component can be simply estimated by considering the muon energy spectrum of a single shower and folding it analytically with the bremsstrahlung cross section and the Greisen parameterization, see . For an $`80^{}`$ zenith and $`10^{19}`$eV proton shower the total number of electrons and positrons ($`N_e`$) obtained is about $`2.510^4`$. These are mostly due to the muons in the energy range between 30 GeV and 500 GeV. As this component arises also from electromagnetic sub-showers its energy distribution should also reflect that of electromagnetic cascades. If we multiply, conservatively, the total number of electrons by a factor three to account for the two other muon interactions, it is still a factor of $``$50 below the number of muons in the shower (see figure 1A).
The contribution to the signal from the electromagnetic component at high zenith angles is in the end small and dominated by muon decay. It has been simulated using WTANK using the muon and electromagnetic signals from different zenith angle showers generated with AIRES as an input. In Fig. 1B we show the simulated ratio of electromagnetic to muon signal in a Čerenkov tank of 1.2 m depth (as used in Haverah Park and planned for the Auger Observatory) as a function of distance to the shower axis for a vertical shower compared to two inclined showers. The shower particles have been fed through the tank simulation as if they were coming from the vertical direction to eliminate geometric tank effects. The results illustrate the genuine decrease in the electromagnetic to muon signal ratio due only to the shower composition varying with zenith angle. Even at distances from the shower axis of order $`1.5`$ km the ratio is about a factor 3 smaller than for vertical showers.
As will be shown in section 4, in the Haverah Park detectors the relative electromagnetic signal drops even further as the zenith angle increases because of geometric effects. At large zenith angles the muon track length is enhanced compared with that in the vertical direction and the output signal increases accordingly. In the analysis that follows we will consider a track dependent, but otherwise constant correction due to the electromagnetic component from muon decay, in agreement with , and neglect that from all muon interactions in the atmosphere.
## 3 The muon component of horizontal showers
The behavior of the muonic component at high zenith angles becomes extremely complex because of magnetic field effects. It is described in detail in Paper I where an accurate model is presented to account for the average muon number densities at ground level. The inputs for this model are the lateral distribution function (LDF), the average muon energy as a function of radius ($`(r)`$) and the mean distance to the muon production point, all evaluated in the absence of magnetic effects. The model is based on an anticorrelation between the average muon energy and distance of the muon from the core. We have generated these inputs, using AIRES, for two primaries proton and iron, and two hadronic models Quark Gluon String Model (QGSJET) and SIBYLL and for primary energies in the range $`10^{16}`$eV to $`10^{20}`$ eV. The muon energy distribution at a fixed distance to the shower axis is assumed to have a log-normal distribution of width 0.4 (Paper I).
Both the $`(r)`$ and the shape of the LDF are essentially invariant over this large energy interval and only mildly dependent on zenith angle as shown in Fig. 2. As a result the dependence on shower energy of the muon number density distributions with the magnetic field can be parameterized with an absolute normalization to a high level of precision. The dependence of the normalization factor with energy can be obtained by monitoring the total muon number in the showers. The results are plotted in Fig. 3 for four different zenith angles. The normalization or energy scaling can be taken into account by a relation of the following form:
$$N_\mu =N_0E^\beta $$
(1)
where $`N_0`$ and $`\beta `$ are constants for a given model and mass composition as shown in Table 1.
Fluctuations can enhance trigger rates for air showers produced by lower energy primaries because of the steep cosmic ray spectrum . The fluctuations to high number of particles allow the more numerous low energy showers to trigger. We have studied muon number fluctuations at ground level and how they depend on shower development (mean muon production height) and average muon energy. We have found that the mean energy correlates strongly with production height but that most of the number density fluctuations can be accounted for by fluctuations in the total muon number. This is mainly due to fluctuations in the neutral to charged pion ratios in the first interactions. According to the results shown in Table (1) of Paper I it is sufficient to implement a $`20\%`$ RMS fluctuation in the average total muon number as obtained with Eq. 1.
Although the actual distributions obtained with simulation do show a long tail to low muon numbers (up to a factor of 4 reduction), such showers are not expected to be relevant for triggering the array precisely because they have so few muons. As a result it is sufficient to assume a gaussian distribution.
## 4 Detector simulation
### 4.1 The Haverah Park Array
The trigger rate of an air-shower array at large zenith angles is extremely sensitive to the geometry of the array. Factors such as the shape and relative height of detectors become very important for such showers. Figure 4A shows the layout of the Haverah Park array. The relative heights and orientations of the four A-site detectors, the triggering detectors, are shown in figure 4B. A gradient across the array is apparent and this has a significant effect on the observed azimuthal distribution (see section 6). Figure 4C shows the positions of individual tanks within a detector hut. The signals from 15 of the 16 tanks, each of area 2.25 m<sup>2</sup>, were summed to provide the signal used in the trigger. One tank in each hut was used to provide a low gain signal. See for a more detailed description of the array.
Water-Čerenkov densities are expressed in terms of the mean signal from a vertical muon (1 vertical equivalent muon or VEM). It has been shown that this signal was equivalent to approximately 14 photoelectrons (pe) .
The formation of a trigger was conditional on: 1. A density of $`>`$0.3 VEM m<sup>-2</sup> in the central detector (A1) 2. At least 2 of the 3 remaining A-site detectors recording a signal of $`>`$0.3 VEM m<sup>-2</sup>.
The rates of the triggering detectors were monitored daily. Over the life of the experiment, after correction for atmospheric pressure effects, the rates of the detectors were stable to better than 5%.
### 4.2 Čerenkov signals of vertical particles
The calculation of the water-Čerenkov signal from horizontal showers is complex. It is informative to consider first the simpler case of vertical showers. A GEANT simulation of the propagation of vertical electrons, gammas, and muons through Haverah Park tanks has been performed (WTANK ). The wavelength dependences of the physical properties of the tank have been taken from but normalised to the following peak values which are the best estimates for the Haverah Park tanks:
* reflectivity of the walls - 83%
* absorption length of the water - 15 m
* photomultiplier (PMT) quantum efficiency - 22%
The Thorn/EMI, 9618 PMT wavelength acceptance function is taken from the manufacturers specifications.
The results of this simulation are summarized in figure 5. It can be seen that the values used enable us to reproduce the measured signal from an average vertical muon (assuming a mean muon energy of 1 GeV ).
The signal from a vertical muon is composed primarily of the Čerenkov light emitted from the muon track. However there is a significant contribution from Čerenkov light emitted by $`\delta `$-ray electrons (2 pe of the 14 pe total for an average vertical muon).
By contrast the mean energy of electrons and gamma-rays in vertical showers is below 10 MeV. Convolving the energy spectrum of electrons and gammas with the response given by figure 5, we find mean signals of 2.6 and 0.9 photoelectrons for electrons and gamma-rays respectively.
### 4.3 Signals produced by horizontal showers
Three factors complicate the picture for horizontal showers. Firstly for horizontal showers it is possible for Čerenkov photons to fall directly onto the PMT without reflection from the tank walls (we refer to such photons as “direct light”). Secondly the azimuthal asymmetry of the Haverah Park tanks (and detector huts) becomes increasingly important at large zenith angles. Thirdly the mean muon energy grows with zenith angle. For 87 showers the mean muon energy is 200 GeV. For higher energy muons the probability of interaction in the tank is much greater. The production of secondary electrons via pair-production, bremsstrahlung, nuclear interactions (collectively referred to as PBN interactions), and electron knock-on ($`\delta `$-rays) is therefore enhanced. For example the correction due to $`\delta `$-ray production increases from 2 pe at typical vertical muon energies of 1 GeV to around 3 pe for $`>10`$ GeV.
Signal enhancements due to direct light and muon interactions appear in only a fraction of events and hence do not greatly affect the peak of the photoelectron distribution. They do however add long tails to this distribution which are of great importance in calculation of the array trigger rate. Figure 6 shows the different contributions to the water-Čerenkov signal of a horizontal muon as a function of energy. A single Haverah Park tank is considered with muons propagating parallel to the short (1.24 m) side of the tank.
For showers of $`>60^{}`$ the electromagnetic component is almost entirely a product of muon decay as explained in section 2. Convolving the number and energy spectrum of electrons and gammas with the tank response shown in figure 5 we find that the electromagnetic component contributes 2.8 pe for each muon corresponding to 20$`\pm `$2 % of the vertical equivalent muon signal.
## 5 Implementation
We have simulated the rate of the Haverah Park array as a function of the zenith angle following the steps described below. First of all we simulated 100, 10<sup>19</sup> eV, showers for zenith angles from $`60^{}`$ to $`88^{}`$ in steps of $`2^{}`$, without magnetic field, for each of the four combinations of primary mass and interaction model as described above. Use is made of the ideas developed in Paper I.
For each zenith angle the following procedure is applied:
1. The inputs of the analytical model (LDF, $`(r)`$ and $`<d>`$) are parameterized.
2. The analytical model is applied to these parameterizations to generate muon density maps of the showers in the transverse plane taking into account magnetic field effects.
3. The density maps for different primary energies are obtained with the energy scaling relationship given by Eq. 1 and the parameters shown in Table 1.
4. Different azimuthal angles are considered by simple rotation algorithms in the approximation explained in section 3 of Paper I.
For each of these zenith angles we have calculated the triggering probability of the Haverah Park array in 40 energy bins between $`10^{16}`$ eV and $`10^{20}`$eV and 18 azimuth bins as follows:
1. The shower is directed on to the array N times with random core locations on the ground plane up to X km from the centre of the array. N and X are varied according to the primary energy and zenith we are dealing with. Typical values are $`N=10000`$ and $`X=8`$ km.
2. Each time a shower is directed at the array, the total muon number ($`N_\mu `$) is fluctuated with a gaussian distribution of spread 0.2$`N_\mu `$ to take into account shower fluctuations.
3. The density in the ground plane at the location of each of the trigger detectors is read from the muon density maps with an appropriate projection, taking into account the corrections due to the different heights of the detectors and the effects of the magnetic field as explained in section 7 of Paper I.
4. The corresponding signal in each of the trigger detectors is generated (see next subsection).
5. The trigger condition of the Haverah Park array is tested.
6. If an event is deemed to trigger the array then the primary energy and the combination of detectors contributing to the trigger are recorded.
The triggering probability derived in this way has been convolved with a primary energy spectrum , derived from Akeno and Haverah Park data, to obtain the event rate at each zenith and azimuth angle. The final step before comparison with data is to convolve the obtained zenith angle distribution with the appropriate measurement errors. The arrival directions of all the data that had zenith angle greater than $`60^{}`$ have been reanalysed assuming a plane front for the muons and using all available timing information. Previously only times from the four central triggering detectors were used to compute the arrival directions . Here up to 16 times were used in the analysis with the additional data being from the detectors $``$2 km from A1 in the ground plane. The uncertainly in arrival direction for each event has accordingly been significantly reduced. Fig. 7 shows the average uncertainty in zenith angle of the reconstructed events as a function of zenith angle. The uncertainty in azimuth angle is approximately constant and is $`1^{}`$. The events considered here have not been analysed for core position. Accordingly some of the very large zenith angle events may have fallen a long way outside the boundary of the array and consequently have a larger uncertainty in zenith angle than suggested by fig. 7.
We have rejected about $`3\%`$ of the events having an error exceeding $`4.5^{}`$ in the reconstructed zenith angle $`\theta `$. The RMS error in zenith angle shown in fig. 7 is used to *smear* the calculated distribution.
### 5.1 Implementation of the signal in detector
The signal is generated as follows:
1. The projected area of the detector in the shower plane is calculated.
2. Given the local muon density and the projected area, we sample the number of incident muons from a Poisson distribution.
3. The track length of each muon through the detector is sampled from a distribution obtained analytically from the detector geometry (see figure 4C).<sup>1</sup><sup>1</sup>1This distribution accounts for the possibility that a single muon may traverse several tanks.
4. The contribution of indirect Čerenkov light from the incident muons and from $`\delta `$-ray electrons is calculated from the sampled track lengths (12 pe for each 1.2 m of track, with an additional 3 pe/1.2 m to account for the signal from $`\delta `$-rays as described in the previous section).
5. The signal from direct light on the PMTs is related to the detector geometry in a more complex way and is implemented using WTANK to simulate the passage of muons through the whole detector for a range of zenith and azimuth angles.
6. The probability of PBN interactions in the detector is calculated from the track length and the muon energy (obtained from the $`(r)`$ with a zenith angle dependent correction for the magnetic field).
7. If an interaction is deemed to occur then the energy of the resulting em-cascade is sampled from the appropriate differential cross-section. The signal produced by this cascade is calculated using an analytical approach.
8. The electromagnetic component of the shower is approximated by the addition of 2.8 pe for each muon as discussed in the previous section.
The contributions to the signal depend in different ways on the zenith and azimuth. The mean signal for those contributions proportional to the track are constant with azimuth and zenith, because of the compensating effect of the projected area of the detector. In fact the mean signal is simply proportional to the volume of the detector. For larger zenith angles fewer muons produce the same signal so Poisson fluctuations become more important. The mean signal produced by direct light increases with zenith and has a complex dependence on azimuth. The mean signal from the electromagnetic component is proportional to the projected area, and hence decreases with zenith and has a simpler azimuthal dependence.
We note that the use of distributions (rather than mean values) for different contributions to the signal is essential for accurate calculation of the event rate.
## 6 Results and comparison with data
### 6.1 Energy distribution of showers that trigger the array
In Fig. 8 we show the energy distribution of showers which trigger the array as calculated with the QGSJET and SIBYLL models and for proton and iron primaries at two zenith angles. The mass dependence seen in Fig. 8 is a consequence of the greater number of muons produced by heavier primaries. The choice of interaction model has a some what smaller effect on the energy response.
Fig. 9 shows the mean and width of the energy distribution as a function of zenith angle (for QGSJET with proton primaries). The median energy changes from $`6\times 10^{17}`$ eV at 70 to $`5\times 10^{18}`$ eV at 84.
### 6.2 Azimuthal distribution
One of the most stringent tests for the simulation of horizontal air showers is the prediction of the azimuthal distributions. Minor corrections in the altitude of the detectors, or the orientation of the tanks can lead to large differences in the rates for inclined showers. The total event rate is dominated by events triggering only 3 of the 4 A-site detectors. The azimuthal distribution is best understood by separating these 3-fold events into 3 sets, each set requiring that a particular detector *does not* trigger. The recorded azimuthal distributions for each of these three sets (A2 out, A3 out and A4 out) are shown in Fig. 10. The calculated distributions (using QGSJET with primary protons) are shown normalised to the data for comparison of the shape. The agreement is reasonable with the exception of the set “A2 out”. There are many possible reasons for this disagreement. In particular we will in future consider reflection and absorption in the ground around the detectors, but we are encouraged by the overall result.
### 6.3 Zenith angle distribution
In Fig. 11 we show the contributions to the total rate (from each of the processes considered) as a function of zenith angle. The QGSJET model and proton primaries are assumed. Similar plots for other hadronic models and primaries differ mainly in their normalization. It should be noted that inclusion of only the muon track contribution would lead to underestimating the total rate by a factor $`3`$ ($`4`$) at 65 (85). The effects of delta rays, direct light, electromagnetic corrections, and pair production and bremsstrahlung are all very significant in the calculation of the overall rate. The effect of fluctuations in the total muon number in a shower represents at most a $`20\%`$ rate enhancement because of the large width of the energy response at a given angle.
The left hand plot of Fig. 12 shows the number of events detected at Haverah Park as a function of the zenith angle compared to the expected number of events for two different models and two mass compositions. The error bars on the data points are statistical. The right hand plot illustrates the effect of the uncertainty in the primary energy spectrum on our result. The maximum and minimum predictions from the spectrum described in (using one sigma errors) are shown together with the result obtained using the primary spectrum parameterized in (all for proton primaries using the QGSJET model).
We stress that no normalization was done with the simulated events. The shape of the simulated distribution is in good agreement with data between $`60^{}`$ and $`70^{}`$ while at higher zenith angles there is a mild disagreement, at most a $``$30% effect. Within the uncertainties resulting from the zenith angle measurement, the primary mass, the interaction models and the primary energy spectrum, the simulated overall rate of showers above $`60^{}`$ is completely consistent with the data.
For most assumptions of primary spectrum and interaction model an intermediate primary mass in the decade around $`10^{18}`$ eV seems to be required to fit the data. Further work is in progress to select higher energy events by considering detectors (B – H of Fig. 4A) that are not in the trigger. For example the event discussed in and believed to be of energy $`10^{20}`$eV struck 20 of the detectors.
## 7 Conclusion
We have presented the first comparison of $`>10^{17}`$ eV air-shower data above $`60^{}`$ with simulation. The agreement is reasonable and allows us to state that the detection of HAS induced by hadronic primaries in an array of water-Čerenkov tanks can be explained in terms of the muonic component. The effects of all the considered contributions to the tank signals have been shown to be extremely important in understanding the total rate.
Our final results differ from the preliminary ones reported earlier in several respects. We have taken into consideration the electromagnetic component from muon decay separately. Previously this contribution was taken together with the correction due to $`\delta `$rays. We have further improved the implementation of the PBN corrections taking into account their increase with the muon energy. We have also simulated the ground number density profiles in a completely different fashion making use of the model described in Paper I.
The general shape agreement between data and simulation both for the zenith angle and azimuthal angle differential rates, together with agreement in the absolute normalization makes us confident that the geometrical considerations have been taken into account correctly, at least at the 20% level. This, together with the general agreement with our preliminary rate simulation, which used Monte Carlo shower simulation rather than the analytical approach described in Paper I, provides strong supporting evidence for this approach as an important tool for HAS studies.
The normalization of the shower rate is sensitive to the model used for high energy hadronic interactions. For proton showers the absolute normalization of the rate differs by up to a factor of two for the two models considered. The absolute rate has also been shown to be very sensitive to the primary composition. There is about a factor 2.2 (3) enhancement, reasonably independent of zenith angle, for the differential rate normalization if the primary particles are assumed to be iron nuclei in the QGSJET (SIBYLL) model. It is interesting to note that current air-shower models are not able to reproduce the observations of high energy particles close to the shower axis above 10<sup>16</sup> eV . This may be a possible explanation of the $``$30% disagreement above 80 degrees (where the surviving muons all originate close to the shower core).
These effects illustrate that HAS measurements may provide another tool with which to study both primary composition and hadronic interactions in ultra high energy cosmic rays. Future detectors such as the Auger Observatories should benefit significantly from the study of cosmic ray showers at large zenith angles.
Acknowledgements: We thank Xavier Bertou for helping us with the angular reanalysis of the data described in and Gonzalo Parente for suggestions after reading the manuscript. This work was partly supported by a joint grant from the British Council and the Spanish Ministry of Education (HB1997-0175), by Xunta de Galicia (XUGA-20604A98), by CICYT (AEN99-0589-C02-02) and by PPARC(GR/L40892).
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# Topological defects: fossils of an anisotropic era?
## I Introduction
It is well know that the ‘Hot Big Bang’ model , despite its numerous successes, is plagued by a number of ‘initial conditions’ problems, of which the horizon, flatness and unwanted relic ones are the best known. The standard way to solve them is to invoke an epoch of cosmological inflation , a relatively brief period of exponential (or quasi-exponential) cosmological expansion. The way inflation solves these problems is, loosely speaking, by erasing all traces of earlier epochs and re-setting the universe to a rather simple state. Indeed, inflation is so efficient in this task that a number of people have wondered if one can ever hope to probe the physics of a pre-inflationary epoch.
There are, however, a small number of possible pre-inflationary relics. For example, the recent work by Turok and collaborators shows that curvature can, in some sense, survive inflation. Another class of inflationary survivors are topological defects , formed at phase transitions either before or during inflation ). It is known (see, e.g. ) that one needs about 20 e-foldings of inflation<sup>*</sup><sup>*</sup>*The exact number is of course model-dependent. to solve the monopole problem. One can reverse the argument and say that monopoles can survive about 20 e-foldings of inflation. The inflationary epoch itself will obviously push the monopoles outside the horizon, but the subsequent evolution of the universe tends to make them come back inside, so if the inflationary epoch is not too long they can still have important cosmological consequences.
Cosmic strings are even more successful, being able to survive about 50 e-foldings. The reason for this difference is that their non-trivial dynamics makes them come back inside the horizon faster than one might naively have expected. The above two numbers are typical, but there are specific models where defects can survive even longer. One example is that of open inflation scenarios . In this case the universe undergoes two different inflationary epochs—roughly speaking, a period of ‘old inflation’ followed by one of ‘new inflation’. As pointed out by Vilenkin , one can expect that defects will form between the two inflationary epochs. In this case, a collaboration including the present authors has recently shown that not only will cosmic strings survive the entire second inflationary epoch, regardless of how long it lastsNote that in these models the duration of the second inflationary epoch is fixed by the present value of the density of the universe ., but they will in fact be back inside the horizon by the time of equal matter and radiation densities. In such models, monopoles can survive up to about 30 e-foldings.
Now, given that defects seem to be so successful surviving inflation, and that one expects them to be frozen out while they are outside the horizon, one can think of a further interesting possibility. For the best-studied case of cosmic strings, it is well known that the scaling properties of the network depend on the background cosmology . Moreover, in some cases (typically when their evolution is friction-dominated) they can retain a ‘memory’ of the initial conditions, or the general properties of the cosmology in which they find themselves at early times, for quite a large number of orders of magnitude in time . It is therefore conceivable that if such an imprint of an early cosmological epoch is retained by a defect network which manages to survive inflation, we might still be able to observe it today.
We believe that this is a general feature of defect models, and a number of non-trivial pieces of information about the very early universe can probably be preserved in this way. In the present paper we will restrict ourselves to a simple example. We will discuss the possibility of a domain wall network retaining information about an early anisotropic phase of the universe. There are very strong constraints on the mass of domain walls formed after inflation, due to the fact that their density decays more slowly than the radiation and matter densities. However, these can be evaded by walls forming before or during inflation. In a subsequent paper, we will discuss the more interesting, but also more complicated, case of cosmic strings.
The plan of the paper is as follows. In section II we briefly describe our background (Bianchi I) cosmology and the basic evolutionary properties of the domain walls. In particular, we focus on the approach to isotropy during inflation, which is discussed through both analytic arguments and numerical simulations. We emphasise that these simulations do not include the defects. However, they serve an important purpose, as they are used in the subsequent discussion to show that the timescale needed for isotropization is compatible with the ‘survival’ on anisotropic defect networks.
We provide a description of our numerical simulations of domain wall evolution in section III. These are analogous to those of Press, Ryden & Spergel , and the interested reader is referred to this paper for a more detailed discussion of some relevant numerical issues. Here defect networks are evolved in an isotropic, matter-dominated (ie, post-inflationary universe), and their main purpose is to show that isotropic and anisotropic networks will evolve in different ways, so two such networks can in principle be observationally distinguished as they re-enter the horizon. Our main results are presented and discussed in section IV, and finally we present our conclusions and discuss future work in section V.
Throughout this paper we will use fundamental units in which $`c=1`$.
## II Evolution equations for domain walls
We consider the evolution of a network of domain walls in a k=0 anisotropic universe of Bianchi type I with line element :
$$ds^2=dt^2X^2(t)dx^2Y^2(t)dy^2Z^2(t)dz^2$$
(1)
where $`X(t)`$,$`Y(t)`$ and $`Z(t)`$ are the cosmological expansion factors in the $`x`$, $`y`$ and $`z`$ directions respectively, and $`t`$ is the physical time. The dynamics of a scalar filed $`\varphi `$ is determined by the Lagrangian density,
$$=\frac{1}{4\pi }\left(\frac{1}{2}\varphi _{,\alpha }\varphi ^{,\alpha }+V(\varphi )\right),$$
(2)
where we will take $`V(\varphi )`$ to be the generic $`\varphi ^4`$ potential with two degenerate minima given by
$$V(\varphi )=V_0\left(\frac{\varphi ^2}{\varphi _0^2}1\right)^2.$$
(3)
This obviously admits domain wall solutions . By varying the action
$$S=𝑑td^3x\sqrt{g},$$
(4)
with respect to $`\varphi `$ we obtain the field equation of motion:
$$\frac{^2\varphi }{t^2}+\theta \frac{\varphi }{t}_{}^{}{}_{}{}^{2}\varphi =\frac{V}{\varphi }.$$
(5)
where
$$_{}^{}{}_{}{}^{2}=\frac{1}{X^2}\frac{^2}{x^2}+\frac{1}{Y^2}\frac{^2}{y^2}+\frac{1}{Z^2}\frac{^2}{z^2},$$
(6)
with $`\theta (t)=\dot{W}/W`$ and $`W(t)=XYZ`$. The dynamics of the universe is described by the Einstein field equations. Here we shall seek perfect fluid solutions. The time component of the Einstein equation then becomes
$$\dot{\theta }+A^2+B^2+C^2=\frac{1}{2}k(\rho +3p),$$
(7)
while the spatial components give
$$\dot{A}+\theta A=\dot{B}+\theta B=\dot{C}+\theta C=\frac{1}{2}k(\rho p),$$
(8)
with $`A=\dot{X}/X`$, $`B=\dot{Y}/Y`$ and $`C=\dot{Z}/Z`$, $`\theta =A+B+C`$ and $`k=8\pi G/c^2`$ and $`i=1,2,3`$. It is straightforward to combine equations (7,8) to obtain:
$$AB+BC+CA=k\rho .$$
(9)
In the following discussion we will make the simplification that $`X(t)=Z(t)`$ (and therefore $`A=C`$) and consider the dynamics of the universe during an inflationary phase with $`\rho =p=const`$. In this case $`H^2k\rho /3=const`$ and the Einstein field equations (7,8,9) imply:
$$\dot{A}+\frac{3}{2}(A^2H^2)=0,$$
(10)
while $`B`$ can be found from the suggestive relation
$$\frac{B}{A}=\frac{1}{2}\left(\frac{3H^2}{A^2}1\right).$$
(11)
Equation (10) has two solutions, depending on the initial conditions. If $`A_i<H`$, then $`A`$ is the smaller of the two dimensions and the shape of spatial hyper-surfaces is similar to that of a rugby ball. Then the solution is
$$\frac{A}{H}=\mathrm{tanh}\left[\frac{3}{2}H(tt_i)+\mathrm{tanh}^1\left(\frac{A_i}{H}\right)\right],$$
(12)
with $`A_i=A(t_i)`$. On the other hand, if $`A_i>H`$, then $`A`$ is the larger of the dimensions and the shape of spatial hyper-surfaces is similar to that of a pumpkin. Then the solution is
$$\frac{A}{H}=\mathrm{coth}\left[\frac{3}{2}H(tt_i)+\mathrm{coth}^1\left(\frac{A_i}{H}\right)\right].$$
(13)
Note that in both cases the ratio $`A/H`$ tends to unity exponentially fast, and hence the same happens with the ratio $`B/A`$. In other words, inflation tends to make the universe more isotropic, as expected. An easy way to see this is to consider the ratio of the two different dimensions, $`D=B/A`$, and to study its evolution equation. One easily finds
$$\dot{D}=\sqrt{6}H\left(D+\frac{1}{2}\right)^{1/2}(1D),$$
(14)
which has an obvious attractor at $`D=1`$.
Note that even though we have so far assumed (for simplicity) that $`p=\rho `$, the same analysis can be carried out for an inflating universe with $`p=(\gamma 1)\rho `$ with $`\gamma 0`$ by numerically solving the conservation equation
$$\dot{\rho }+\theta (\rho +p)=0,$$
(15)
together with equations (8) and (9). Indeed, the more general case will be relevant for what follows.
In figure 1 we plot the evolution of the asymmetry parameter $`E=Y(t)/X(t)`$, according to eqns. (7) and (8), for several values of $`\alpha _i=\mathrm{log}_{10}(D_i)`$ assuming $`\gamma =0`$ and $`\gamma =2/3`$ (note that $`\gamma =2/3`$ is the maximum value of $`\gamma `$ which violates the strong energy condition). Note that we do not include the defect network in the simulation. (We assume that the network at the initial time $`t_i`$ is statistically isotropic.) We take $`X(t_i)=Y(t_i)`$. We can see that depending on the initial degree of anisotropy, specified by $`\alpha _i`$, the value of $`E`$ can grow to be very large, especially if $`\alpha _i`$ is large. Moreover, although for $`\gamma =0`$, the value of $`E`$ becomes approximately constant in one Hubble time that does not happen so rapidly for inflating universes with larger $`\gamma `$. This removes the necessity of producing the domain walls right at the onset of the inflationary era.
What about the evolution of the domain walls? Based on rather general grounds, we expect it to have a number of similarities with the much better studied case of cosmic strings . In particular, one can define a ‘characteristic length-scale’, which we shall denote by $`L`$, that can be roughly interpreted as a typical curvature radius or a correlation length of the wall network. It is also a length-scale that measures the total energy of the domain wall network per unit volume, since we can define
$$\rho \frac{\sigma }{L},$$
(16)
where $`\sigma `$ is the domain wall energy per unit area. Note that in a more rigorous treatment that allowed for the expected build-up of small-scale ‘wiggles’ on the walls (in analogy with what happens for the case of cosmic strings ) each of these three length scales would be different. However, for our present purposes it is adequate to suppose that they are all similar.
Then we can expect to find two different evolution regimes. While the network is non-relativistic, we expect it to be conformally stretched by the cosmological expansion, and hence
$$La,\rho _wa^1.$$
(17)
In this case there is essentially no dynamics. An extreme example of this regime happens during inflation We can see from eqn. (5) that due to the very rapid expansion which occurs in the inflationary regime the time derivatives of the field $`\varphi `$ rapidly approach zero so that the network of domain walls will simply be frozen in comoving coordinates.
On the other hand, once the network becomes relativistic, one expects it to evolve in a linear scaling regime where
$$Lt,\rho _wt^1.$$
(18)
This is the case of ‘maximal’ dynamics, in the sense that the network is evolving (in particular, losing energy by wall collisions and re-connections) as fast as allowed by causality. We note that previous work of Press, Ryder and Spergel suggests that there may be logarithmic corrections to this linear regime.
## III Numerical simulations
At late times (after the inflationary epoch) the universe is homogeneous and isotropic with $`A=B=C`$ with the average dynamics of the universe being specified by the evolution of the scale-factor $`a(t)`$. We now consider the evolution of isotropic and anisotropic defect networks in this background. In particular, we are interested in determining how the networks evolve as they re-enter the horizon, since if one finds differences in the dynamics of the two cases then this should translate into observational tests that will allow us to discriminate between then and hence probe pre-inflationary physics.
It is useful for numerical purposes to re-write equation (5) as a function of the conformal time $`\eta `$ defined by $`d\eta =dt/a`$. In this case equation (5) becomes
$$\frac{^2\varphi }{\eta ^2}+2\frac{\dot{a}}{a}\frac{\varphi }{\eta }^2\varphi =a^2\frac{V}{\varphi }.$$
(19)
with
$$^2=\frac{^2}{x^2}+\frac{^2}{y^2}+\frac{^2}{z^2}.$$
(20)
When making numerical simulations of the evolution of domain wall networks (or indeed other defects) it is also often convenient to modify the equation of motion for the scalar field $`\varphi `$ in such a way that the comoving thickness of the walls is fixed in comoving coordinates. This is known as the PRS algorithm , and it is generally believed not to significantly affect the large-scale dynamics of domain walls.
We note, however, that recent high-resolution simulations have revealed that the accuracy of this algorithm is not as good as has been claimed. This effect is expected to increase with increasing dynamic range. In particular, the PRS algorithm artificially prevents the build-up of small-scale features on the domain walls (or, for that matter, any other defect). This turns out to be crucial for a quantitatively accurate description of their evolution, and hence for a reliable analysis of their observational consequences. For our purposes in the present work, however, the PRS algorithm is enough as an approximation to the true wall dynamics. In a subsequent, more detailed publication we shall compare results obtained using this algorithm with those from the true wall dynamics.
Having clarified this point, we will modify the evolution equation for the scalar field $`\varphi `$ in the isotropic phase according to the PRS prescription:
$$\frac{^2\varphi }{\eta ^2}+\beta _1\frac{\dot{a}}{a}\frac{\varphi }{\eta }^2\varphi =a^{\beta _2}\frac{V}{\varphi }.$$
(21)
where $`\beta _1`$ and $`\beta _2`$ are constants. We choose $`\beta _2=0`$ in order for the walls to have constant comoving thickness and $`\beta _1=3`$ by requiring that the momentum conservation law for how a wall slows down in an expanding universe is maintained .
We perform two-dimensional simulations of domain wall evolution for which $`^2\varphi /z^2=0`$. These have the advantage of allowing a larger dynamic range and better resolution than tree-dimensional simulations.
We solve equation (21) numerically assuming a matter-dominated Einstein-de-Sitter cosmology with $`a\eta ^2`$. We used a standard difference scheme second-order accurate in space and time and periodic boundary conditions (see for a more detailed description of the algorithm and other related numerical issues).
The initial properties of the network of domain walls depend strongly on the details of the phase transition which originated them. It is conceivable that the initial network is already formed asymmetric with the walls being elongated along preferred directions. However, this is beyond the scope of the present paper. For our present purposes, we can ignore this possibility and assume that the initial domain wall network is statistically isotropic. This assumption will not modify the conclusions of the paper—if anything, any ab initio anisotropies would only enhance the effects we are describing.
Hence, we assume the initial value of $`\varphi `$ to be a random variable between $`\varphi _0`$ and $`\varphi _0`$ and the initial value of $`\dot{\varphi }`$ to be equal to zero everywhere. We normalise the numerical simulations so that $`\varphi _0=1`$. We set the conformal time at the start of the simulation and the comoving spacing between the mesh points to be respectively $`\eta _i=1`$ and $`\mathrm{\Delta }x=1`$.
The wall thickness, defined by
$$\omega _0=\frac{\pi \varphi _0}{\sqrt{2V_0}}$$
(22)
is set to be equal to $`5`$. The kinetic energy of the field $`\varphi `$ is calculated by
$$E_{kin}=\frac{1}{8\pi }\underset{i,j}{}\dot{\varphi }_{i,j}.$$
(23)
On the other hand, the rest energy of the walls is calculated by multiplying the comoving area of the walls, A, by the energy density per comoving area, which can be written as
$$\sigma =\frac{2\gamma _wV_0\omega _0}{3\pi ^2}.$$
(24)
with $`\gamma _w=(1v_w^2)^{1/2}`$, and $`v_w`$ is the value of the physical velocity of the domain walls. Finally the total area of the walls is defined as the area of the surfaces on which $`\varphi =0`$ and is computed using the method described in ref. .
## IV Results and discussion
As pointed out above, it will be of fundamental importance to study the dynamics of the wall network at late times, as the Hubble length becomes larger than the typical size of the major axis of a domain wall. A crucial issue will be the timescale required for the wall network to switch from the non-relativistic regime to the relativistic one. For our present purposes, the main difference between these two regimes is that a friction-dominated network can remain anisotropic if it froze out that way, whereas a relativistic network will rapidly become isotropic and erase any imprints from the earlier anisotropic phase. In our simulations we ignore the possibility that the network can be friction dominated due to particle scattering when the domain walls come back inside the horizon—again, this would only enhance the effects we are describing.
We consider three simulations with different initial conditions. In the first one (case I) we evolve the initial network generated in the manner specified in the previous section from the conformal time $`\eta _i=1`$. In the second one (case II) the initial conditions at the time $`\eta _i=1`$ were specified by the network configuration of the previous simulation at the conformal time $`\eta _{}=20`$, with the velocities reset to zero. Physically, this corresponds to starting with the network outside the horizon. Finally, the case III is similar to the second one but with the initial network of case II stretched in the $`y`$ direction by a factor of $`E=2`$ (see fig. 2), and corresponds to the anisotropic case. We have performed $`1024^2`$ simulations for each of the three cases, plus an additional $`2048^2`$ run of case I, in order to test for possible box effects.
For each run we plot (see fig. 3) the ratios $`A/V`$ and $`A\eta /V`$ (note that $`A`$ and $`V`$ are the comoving area and volume, respectively), as well as the ratio of the kinetic and rest energies, as in Press, Ryden and Spergel . These are plotted from the beginning of the simulation until the time when the horizon becomes one half (for the $`1024^2`$ runs) or one quarter (for the $`2048^2`$ run) of the box size. In addition to these (which we plot mainly for the purposes of comparison with previous work ) we plot a ‘scaling coefficient’ which will be our main analysis tool. We will define it by analogy with the cosmic string case , as follows. Assume that
$$\eta \frac{A}{V}\eta ^\lambda ;$$
(25)
then what we plot is the ‘instantaneous’ or ‘effective’ value of $`\lambda `$ as a function of conformal time. For a given $`\lambda `$, the physical network correlation $`L`$ length will be evolving as
$$Lt^{1\lambda /3},\frac{L}{a}\eta ^{1\lambda }.$$
(26)
Note that $`\lambda `$ can, in general, be a time-dependent quantity. However, for the two scaling regimes discussed above, we expect it to be a constant, namely
$$\lambda _{nr}=1$$
(27)
in the non-relativistic limit where the network is being conformally stretched, and
$$\lambda _r=0$$
(28)
in the linear scaling regime.
From these it is trivial to deduce the behaviour of $`A/V`$ and $`A\eta /V`$ in both scaling regime. One expects
$$\eta \frac{A}{V}\eta \frac{A}{V}const.$$
(29)
in the non-relativistc regime and
$$\eta \frac{A}{V}const\frac{A}{V}\eta ^1$$
(30)
in the linear scaling regime. Similarly, the ratio $`E_k/E_{rest}`$ should be a constant in the linear scaling regime (with its numerical value providing a measure of the characteristic network scaling speed), and it should approach zero in the non-relativistic limit.
Firstly, we note that the two case I runs produce very similar results: significant differences can only be seen at late times. This is an indication that the resolution we are using is adequate for our present purposes. As expected, the network in case I becomes relativistic very quickly, while those of cases II and III start in the extreme non-relativistic regime and only evolve away from it fairly slowly, after they re-enter the horizon.
More importantly, there are two non-trivial observations to be made. Firstly, we confirm that there is a correction to the linear scaling regime. We find
$$\lambda _{sc}0.12,$$
(31)
which corresponds to to evolve in a linear scaling regime where
$$Lt^{0.96},\rho _wt^{0.96},$$
(32)
in agreement with the previous result by Press, Ryden and Spergel . This means that the network is not straightening out as fast as allowed by causality. Secondly, the rates at which the networks in cases II and III approach the relativistic regime are different. One might expect this on physical grounds: if the network is stretched in one direction, then there are in fact different ‘network correlation lengths’ for each direction, and interactions between the domain walls will tend to occur faster along the directions with smaller correlation lengths, and more slowly in the others.
Another way of saying this is that the network will only start evolving towards the relativistic regime when its larger axis has re-entered the horizon. Note that this mechanism also tends to make the domain wall network more isotropic. So one can naively say that the approach to the linear scaling regime takes longer in an anisotropic universe because the dynamics of the walls must accomplish two tasks (make the wall network relativistic and isotropic) rather than just one.
## V Conclusion
In this paper we have discussed a simple example of what we believe to be a rather generic feature of topological defect models, namely that they can easily retain information about the properties of the very early universe. This information is encoded in the scaling (ie, ‘macroscopic’) and statistical (ie, ‘microscopic’) properties of the defect networks. This is even more relevant given the fact that defects can survive significant amounts of inflation. Hence, they can provide a unique probe of the pre-inflationary universe. The two crucial scales in the problem are the defect mass scale and the epoch when the defects come back inside the horizon.
Specifically, we have discussed the role of domain walls. We have highlighted the existence of two scaling regimes for the domain wall network, in agreement with previous work . Furthermore, we have shown that an anisotropic network re-entering the horizon will take longer to approach scaling than an isotropic one. Hence, if the very early universe had an anisotropic phase which was erased by an inflationary epoch, and if domain walls are present, then the walls can retain an imprint of the earlier phase, and this can have important observational consequences, eg for structure formation scenarios.
As is well known, there are quite strong constraints on the mass of domain walls formed after inflation. These are basically due to the fact that their density will decay more slowly than the radiation and matter densities. However, essentially all of these can be evaded (or at least significantly relaxed) by walls forming before or during inflation (and also by walls evolving in a friction-dominated regime). Having said this, how could these anisotropies be detected? The most naive answer would be through their imprint on CMB, but this is only true if their energy density is not too low, and such models are constrained in a variety of other ways (not only from the cosmology side, but also from the high-energy physics side). The case of ‘light’ walls is therefore more interesting: note that just like in the case of ‘light strings’ , these are expected to be friction-dominated throughout most of the cosmic history. Here the observational detection of the effects we have described becomes somewhat non-trivial. The best way of doing it should be through observations of numbers of objects as a function of redshift in different directions (assuming that one has a reliable understanding of other possible evolutionary effects). Two specific examples would be large-scale velocity flows and gravitational lensing statistics of extragalactic surveys .
Finally, there is also an important implication of our work if at least one of the minima of the scalar field potential has a non-zero energy density, which is an anisotropic non-zero vacuum density. In a subsequent, more detailed publication, we shall discuss this scenario in more detail, as well as the analogous one for cosmic strings.
To conclude, we have shown that the importance of topological defects as a probe of cosmological physics goes well beyond structure formation. Even if defects turn out to be unimportant for structure formation they can still (if detected) provide us with extremely valuable information about the physical conditions of the very early universe.
###### Acknowledgements.
We would like to thank Paulo Carvalho for enlightening discussions. C.M. is funded by JNICT (Portugal) under ‘Programa PRAXIS XXI’ (grant no. PRAXIS XXI/BPD/11769/97).
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# 1 Introduction
## 1 Introduction
An important set of classically integrable field theories in 1+1 dimensions is provided by non-linear sigma-models with target manifold a symmetric space $`G/H`$. These have been studied for many years, particularly in connection with a one-parameter ‘dual symmetry’ which allows the construction of both local and non-local conserved quantities. Little is known about the local charges arising from this approach, however, with comparatively few examples which can be written in closed form. A comprehensive discussion, with many references, can be found in . Examples of more recent work involving various models based on symmetric spaces are .
Principal chiral models (PCMs) can be regarded as special cases of symmetric space models (SSMs), with target manifold a Lie group $`G`$. In , it was shown that other well-known local conservation laws in the PCMs, apparently different from those arising from dual symmetry, exhibit striking and hitherto unexpected properties. In particular, it was possible to construct mutually-commuting sets of charges with a characteristic pattern of spins given by the exponents of each classical group $`G`$ modulo its Coxeter number. The same pattern of spins arises in affine Toda field theories, and proves central to understanding a number of their most important properties .
In this note we will show that results similar to those of can be obtained for any sigma-model based on a symmetric space $`G/H`$ with $`G`$ a classical group. Specifically, we will consider local conserved currents which can be written in a simple, closed form, and which lead to commuting sets of charges whose spins are related to the underlying symmetric space data. We begin by formulating the field theory and the relevant conservation laws.
Let $`g(x^\mu )`$ be a field on two-dimensional Minkowski space taking values in some compact Lie group $`G`$, with Lie algebra $`𝐠`$. Let $`HG`$ be some subgroup and $`𝐡𝐠`$ the corresponding Lie subalgebra. To formulate the sigma-model with target space $`G/H`$, we introduce a gauge field $`A_\mu `$ in $`𝐡`$ and define a covariant derivative
$$D_\mu g=_\mu ggA_\mu $$
(1.1)
with the property that
$$ggh,A_\mu h^1A_\mu h+h^1_\mu hD_\mu g(D_\mu g)h$$
(1.2)
for any function $`h(x^\mu )`$. It is also useful to introduce the $`𝐠`$-valued currents
$$j_\mu =g^1_\mu g,J_\mu =g^1D_\mu g=j_\mu A_\mu .$$
(1.3)
The latter current is covariant under gauge transformations, with
$$J_\mu h^1J_\mu h.$$
(1.4)
The $`G/H`$ sigma-model is defined by the lagrangian
$$=\frac{1}{2}\mathrm{Tr}(J^\mu J_\mu )=\frac{1}{2}\mathrm{Tr}(g^1D^\mu gg^1D_\mu g)$$
(1.5)
which has a global $`G`$ symmetry (acting from the left on $`g`$) and the local $`H`$ symmetry discussed above. The equation of motion for the field $`g`$ is
$$D_\mu J^\mu =_\mu J^\mu +[A_\mu ,J^\mu ]=0.$$
(1.6)
By combining this with the identity $`_\mu j_\nu _\nu j_\mu +[j_\mu ,j_\nu ]=0`$ we obtain
$$2_{}J_++2[A_{},J_+]=[J_+,J_{}]+F_+$$
(1.7)
where $`F_{\mu \nu }=_\mu A_\nu _\nu A_\mu +[A_\mu ,A_\nu ]`$, and we have made use of light-cone coordinates, $`V_\pm =V_0\pm V_1`$. The equations of motion for the $`A_\mu `$ fields are
$$J_\mu =0\mathrm{on}𝐡.$$
(1.8)
Everything we have said so far applies to an arbitrary homogeneous space $`G/H`$. The special nature of symmetric spaces emerges when we look for conserved quantities. It is helpful to recall what happens for the PCM based on $`G`$, which can be obtained by specializing the analysis above to the case in which $`H`$ is trivial, setting $`A_\mu =0`$, and hence $`J_\mu =j_\mu `$. The equation of motion then becomes $`_{}j_+=\frac{1}{2}[j_+,j_{}]`$ which implies conservation equations such as $`_{}\mathrm{Tr}(j_+^m)=0`$. For a general, non-trivial gauge group $`H`$, the additional term $`F_+`$ in (1.7) prevents one from carrying out a similar construction in any obvious way. But for $`G/H`$ a symmetric space there is an orthogonal decomposition of the Lie algebra, and a compatible $`𝐙_\mathrm{𝟐}`$ grading of the Lie bracket,
$$𝐠=𝐡+𝐤:[𝐡,𝐡]𝐡,[𝐡,𝐤]𝐤,[𝐤,𝐤]𝐡.$$
(1.9)
Since the $`A_\mu `$ equations of motion force $`J_\mu `$ to take values in $`𝐤`$, this grading then implies that the left- and right-hand sides of (1.7) must vanish separately:
$$_{}J_+=[A_{},J_+]\mathrm{in}𝐤,[J_+,J_{}]=F_+\mathrm{in}𝐡.$$
(1.10)
The first of these conditions allows the construction of conservation laws very similar to those of the PCM, e.g. $`_{}\mathrm{Tr}(J_+^m)=0`$. A conserved current of this type can be written down using any symmetric invariant tensor, and we now discuss the possibilities.
## 2 Currents and invariants on $`G/H`$
Let us introduce a basis of anti-hermitian generators $`\{t^a\}`$ for $`𝐠`$, obeying $`[t^a,t^b]=f^{abc}t^c`$ and $`\mathrm{Tr}(t^at^b)=\delta ^{ab}`$. Since $`𝐠`$ is compact, we need not distinguish upper and lower Lie algebra indices, and the structure constants $`f^{abc}`$ are real and totally antisymmetric. We can assume our basis is chosen so that it splits $`a(\widehat{\alpha },\alpha )`$, with the subsets $`\{t^{\widehat{\alpha }}\}`$ and $`\{t^\alpha \}`$ providing bases for $`𝐡`$ and $`𝐤`$ respectively in the decomposition in (1.9). The currents of the $`G/H`$ SSM model can now be written $`j_\mu ^a`$, while the non-trivial components of the gauge fields are $`A_\mu ^{\widehat{\alpha }}`$. The symmetric space condition, or $`𝐙_2`$ grading, means that the non-vanishing structure constants are $`f_{\widehat{\alpha }\widehat{\beta }\widehat{\gamma }}`$, and $`f_{\widehat{\alpha }\beta \gamma }`$, up to permutations of indices.
Consider some totally symmetric tensor $`d_{a_1\mathrm{}a_m}^{(m)}`$ of degree $`m`$ on $`𝐠`$ (we use ‘degree’ rather than ‘rank’ to avoid confusion with the rank of the algebra; we shall not always indicate the degree explicitly). By virtue of its symmetry, this tensor is completely determined by the associated function $`d(X)d_{a_1\mathrm{}a_m}X^{a_1}\mathrm{}X^{a_m}`$ where $`X=X^at^a`$ is an arbitrary element of the Lie algebra. If $`\tau `$ is any map from $`𝐠`$ to itself, then we define a new tensor $`\tau (d)`$ by $`\tau (d)(X)=d(\tau (X))`$. We shall call $`d`$ a $`G`$-invariant tensor, or simply an invariant tensor on $`𝐠`$, if $`d=\tau (d)`$ whenever $`\tau `$ is an inner automorphism of $`𝐠`$, so that $`\tau (X)=gXg^1`$ for some $`g`$ in $`G`$. This is equivalent to the condition
$$d_{c(a_1\mathrm{}a_{m1}}^{(m)}f_{a)bc}^{}=0.$$
Similar definitions apply in an obvious way to subgroups of $`G`$ acting on subspaces of $`𝐠`$.
A conserved charge of spin $`s`$ in the PCM based on $`G`$ can be constructed from each symmetric invariant tensor $`d_{a_1a_2\mathrm{}a_sa_{s+1}}^{(s+1)}`$ . A related conservation law in the $`G/H`$ SSM arises by restricting such a tensor to $`𝐤`$, thus considering just the components $`d_{\alpha _1\alpha _2\mathrm{}\alpha _s\alpha _{s+1}}^{(s+1)}`$. The invariance condition written above means that the restricted tensor obeys
$$d_{\gamma (\alpha _1\mathrm{}\alpha _s}^{(s+1)}f_{\alpha )\widehat{\beta }\gamma }=0.$$
(2.1)
It is easy to check directly that this implies
$$_{}(d_{\alpha _1\mathrm{}\alpha _{s+1}}J_+^{\alpha _1}\mathrm{}J_+^{\alpha _{s+1}})=0$$
(2.2)
on using the equation of motion in (1.10). From this current we obtain a conserved charge of spin $`s`$ in the usual way:
$$q_s=𝑑xd_{\alpha _1\mathrm{}\alpha _{s+1}}J_+^{\alpha _1}\mathrm{}J_+^{\alpha _{s+1}}.$$
(2.3)
Similar conserved charges can be constructed from $`J_{}^\alpha `$; their properties are directly analogous and we will not discuss them further.
It is important to check that the tensor $`d^{(s+1)}`$ does not vanish when restricted to $`𝐤`$ in order to have a non-trivial conserved quantity in the SSM. We shall return to this point in section 6. It is also clear that (2.1) and (2.2) rely only on the fact that the tensor is $`H`$-invariant on $`𝐤`$. It is not obvious, a priori, that any such tensor should arise as the restriction of some $`G`$-invariant tensor on $`𝐠`$, but this emerges from the analysis below. These and other matters can be understood in terms of the special role played by primitive invariants.
For each Lie algebra $`𝐠`$ there are exactly $`\mathrm{rank}(G)`$ primitive symmetric invariants, with the property that all others can be written as polynomial functions of them (see e.g. ). This happens essentially because any element in $`𝐠`$ is conjugate to some element in a fixed Cartan subalgebra (CSA), and so any invariant tensor is determined by its values on the $`\mathrm{rank}(G)`$ independent basis elements of this CSA. The degrees of the primitive invariant tensors for each classical group $`G`$ are given in the table, in terms of the exponents, $`s`$.
One convenient choice for the primitive invariants consists of symmetric traces, with $`d(X)=\mathrm{Tr}(X^m)`$ (of the appropriate degrees), together with the Pfaffian invariant $`d(X)=\mathrm{Pf}(X)ϵ_{i_1j_1\mathrm{}i_nj_n}X_{i_1j_1}\mathrm{}X_{i_nj_n}`$ for the special case of $`SO(2n)`$. Other choices are possible, and will be important later. However, any choices of the primitive invariants differ only by terms which are products of polynomials of lower degrees.
We need to determine how these facts generalize from a group $`G`$ to a symmetric space $`G/H`$. One can define a CSA for $`G/H`$ as a maximal set of mutually commuting generators in $`𝐤`$, and $`\mathrm{rank}(G/H)`$ is then the number of elements in such a set. It can also be shown that any element of $`𝐤`$ is conjugate, by an element of $`H`$, to a member of some chosen CSA . This means that, just as for groups, any invariant is determined by its values on the CSA, and there are precisely $`\mathrm{rank}(G/H)`$ primitive invariants, in terms of which all others can be expressed.
The degrees of the primitive invariants for each symmetric space $`G/H`$ with $`G`$ classical are given by the data in the table. Each primitive $`H`$-invariant tensor is obtained by restricting a primitive $`G`$-invariant tensor on $`𝐠`$ to the subspace $`𝐤`$. Other invariants which are primitive on $`𝐠`$ may vanish when restricted to $`𝐤`$, or else fail to be primitive on $`𝐤`$ in some more complicated fashion. For our purposes we may take the values $`s`$ given in the table as a definition of the exponents of $`G/H`$. The Pfaffian in $`SO(2n)`$ appears as something of a special case, lying outside the regular sequence formed by the other invariants. Because of this, we have chosen to separate the values of $`s`$ corresponding to the Pfaffian, or its restriction, by a semi-colon.
$$\begin{array}{cccc}G/H\text{symmetric space}& \text{rank}(G/H)& s:d^{(s+1)}\text{primitive}& h\\ & & & \\ & & & \\ & & & \\ SU(n)& n1& 1,2,\mathrm{},n1& n\\ & & & \\ & & & \\ & & & \\ SO(2n+1)& n& 1,3,\mathrm{},2n1& 2n\\ & & & \\ & & & \\ & & & \\ SO(2n)& n& 1,3,\mathrm{},2n3;n1& 2n2\\ & & & \\ & & & \\ & & & \\ Sp(2n)& n& 1,3,\mathrm{},2n1& 2n\\ & & & \\ & & & \\ & & & \\ SU(p+q)/S(U(p)\times U(q))(pq)& p& 1,3,\mathrm{},2p1& 2p\\ & & & \\ & & & \\ & & & \\ SO(p+q)/SO(p)\times SO(q)(p<q)& p& 1,3,\mathrm{},2p1& 2p\\ & & & \\ SO(2n)/SO(n)\times SO(n)& n& 1,3,\mathrm{},2n3;n1& 2n2\\ & & & \\ & & & \\ & & & \\ Sp(2p+2q)/Sp(2p)\times Sp(2q)(pq)& p& 1,3,\mathrm{},2p1& 2p\\ & & & \\ & & & \\ & & & \\ SU(n)/SO(n)& n1& 1,2,\mathrm{},n1& n\\ & & & \\ & & & \\ & & & \\ Sp(2n)/U(n)& n& 1,3,\mathrm{},2n1& 2n\\ & & & \\ & & & \\ & & & \\ SO(2n)/U(n)& [n/2]& 1,3,\mathrm{},2[n/2]1& 2[n/2]\\ & & & \\ & & & \\ & & & \\ SU(2n)/Sp(2n)& n1& 1,2,\mathrm{},n1& n\\ & & & \\ & & & \end{array}$$
We obtained the results in the table by making use of convenient canonical forms for the CSA generators in $`G/H`$. In view of the remarks above, it is the set of eigenvalues of the CSA generators, and how they behave under $`H`$, which determines the allowed invariants, and which of them are primitive. By comparing with the well-known data for classical groups, it is not difficult to arrive at the results for symmetric spaces. The method is best illustrated by some examples, but a fuller explanation of this sort requires some additional technical preparation, and so we consign these details to a separate section below. Our main task—understanding the conserved charges in the $`G/H`$ sigma-model—will then be resumed in section 4.
## 3 Some details and examples
Recall that the CSA of each classical Lie algebra can be parameterized by a set of real ‘eigenvalues’ $`\lambda _i`$ as follows
$`su(n):`$ $`\mathrm{diag}(i\lambda _1,\mathrm{},i\lambda _n)\mathrm{with}\lambda _1+\mathrm{}+\lambda _n=0`$
$`so(2n):`$ $`\mathrm{diag}(\left[\begin{array}{cc}0& \lambda _1\\ \lambda _1& 0\end{array}\right],\mathrm{},\left[\begin{array}{cc}0& \lambda _n\\ \lambda _n& 0\end{array}\right])`$
$`so(2n+1):`$ $`\mathrm{diag}(\left[\begin{array}{cc}0& \lambda _1\\ \lambda _1& 0\end{array}\right],\mathrm{},\left[\begin{array}{cc}0& \lambda _n\\ \lambda _n& 0\end{array}\right],\mathrm{\hspace{0.17em}0})`$
$`sp(2n):`$ $`\mathrm{diag}(i\lambda _1,\mathrm{},i\lambda _n,i\lambda _1,\mathrm{},i\lambda _n)`$
where we have used an obvious block notation for the orthogonal algebras. The function $`d(X)`$ defined by a symmetric invariant tensor $`d`$ on $`𝐠`$ is some polynomial in the eigenvalues $`\lambda _i`$. This polynomial must be totally symmetric, because in all cases a member of the CSA of $`𝐠`$ can be conjugated by specific elements of $`G`$ so as to permute the eigenvalues in any desired way. (To show this it actually suffices to use the simple result $`\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)\left(\begin{array}{cc}a& 0\\ 0& b\end{array}\right)\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)=\left(\begin{array}{cc}b& 0\\ 0& a\end{array}\right)`$ in appropriate block forms.)
Now, for $`𝐠=su(n)`$ there are invariant tensors corresponding to all symmetric polynomials in the eigenvalues, except for $`\lambda _1+\mathrm{}+\lambda _n=0`$. The symmetrized traces mentioned earlier clearly give rise to the power sums $`_i\lambda _i^m`$, and a basis of primitive invariants corresponds to the finite subset with $`m=s+1`$ and $`s`$ an exponent. On the other hand, for $`so(2n+1)`$ and $`sp(2n)`$, only polynomials in even powers of $`\lambda _i`$ are allowed, since the sign of any eigenvalue can be reversed by conjugating with a suitable element of $`G`$. For example, conjugating a CSA element of $`𝐠=so(2n+1)`$ with $`\mathrm{diag}(1,1,0,\mathrm{},0,1)`$ (which certainly belongs of $`G=SO(2n+1)`$) changes the sign of $`\lambda _1`$. For $`𝐠=so(2n)`$ it is also possible to reverse the signs of eigenvalues, but only in pairs. For instance, we can conjugate by $`\mathrm{diag}(1,1,1,1,0,\mathrm{},0)`$ which reverses the signs of $`\lambda _1`$ and $`\lambda _2`$, but we cannot conjugate by any element of $`G=SO(2n)`$ and change the sign of just one eigenvalue. In addition to the even symmetric powers, these symmetry properties allow precisely one more independent invariant, which is the Pfaffian, proportional to $`\lambda _1\mathrm{}\lambda _n`$.
In this manner the allowed invariants and primitive invariants for each Lie algebra are characterized as certain totally symmetric polynomials in the CSA eigenvalues. Moreover, we can distinguish three classes of polynomials, depending on their additional symmetry properties: A-type, like $`su(n)`$—no additional symmetries; B/C-type, like $`so(2n+1)`$ or $`sp(2n)`$—invariant under reversal of sign of each eigenvalue separately; D-type, like $`so(2n)`$—invariant under reversal of signs of pairs of eigenvalues.<sup>3</sup><sup>3</sup>3 These symmetry operations on the CSA eigenvalues are usually presented as actions of the Weyl group.
The results for each symmetric space $`G/H`$ can now be found by the following steps. First, find a convenient parameterization of $`𝐤`$ and its CSA, as in e.g. ; $`H`$-invariant tensors on $`𝐤`$ are polynomials in the ‘eigenvalues’ of the CSA generators. Next, examine the action (via conjugation) of specific elements of $`H`$ on these eigenvalues, and so determine that the polynomials have symmetry type A, B/C, or D using the terminology introduced above. Finally, check that all primitive invariants of this type indeed arise as restrictions of $`G`$-invariant tensors on $`𝐠`$, e.g. by considering traces or symmetric powers. We now sketch how this works for some examples; the remaining cases in the table can be handled similarly.
The Grassmannians $`SO(p+q)/SO(p)\times SO(q)`$ and $`SU(p+q)/S(U(p)\times U(q))`$ are conveniently treated together. For either family the subgroup $`HG`$ has the block structure $`\left(\begin{array}{cc}P& 0\\ 0& Q\end{array}\right)`$ and $`𝐤`$ consists of matrices $`\left(\begin{array}{cc}0& X\\ X^{}& 0\end{array}\right)`$ where $`P`$ is $`(p\times p)`$, $`Q`$ is $`(q\times q)`$, and $`X`$ is $`(p\times q)`$ with real or complex entries. For both the real and complex families, the CSA can be parameterized by
$$X=\left(\begin{array}{ccccccc}\lambda _1& 0& \mathrm{}& 0& 0& \mathrm{}& 0\\ 0& \lambda _2& \mathrm{}& 0& 0& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& & \mathrm{}\\ 0& 0& \mathrm{}& \lambda _p& 0& \mathrm{}& 0\end{array}\right)$$
with real ‘eigenvalues’ $`\lambda _i`$ (we assume $`pq`$). Notice that the effect on $`𝐤`$ of conjugating by a general element of $`H`$ is $`XPXQ^1`$. One can readily choose $`P`$ and $`Q`$ so as to permute the CSA eigenvalues in any desired way.
When $`p<q`$, it also easy to find elements of $`H`$ which change the sign of any given eigenvalue, just as in the previous discussion of $`so(2n+1)`$. This implies that the invariant polynomials are of B/C-type. They arise as restrictions of symmetric traces on $`𝐠`$, and in fact they can be written $`\mathrm{Tr}(XX^{}\mathrm{}XX^{})`$, which is manifestly invariant under $`H`$. When $`p=q`$, a new feature arises for the real Grassmannians: just as for the Lie algebra $`so(2p)`$, it is now only possible to change the signs of eigenvalues in pairs, so the pattern of invariants is type D. An additional primitive invariant on $`𝐤`$ arises as the restriction of the Pfaffian on $`𝐠=so(2p)`$, and it can be written $`ϵ_{i_1\mathrm{}i_p}ϵ_{j_1\mathrm{}j_p}X_{i_1j_1}\mathrm{}X_{i_pj_p}`$. Note that under the action of $`H`$ this expression changes by a factor $`\mathrm{det}(P)\mathrm{det}(Q^1)`$ which is indeed unity for these examples. For the complex Grassmannians, however, the invariants are still of type B/C, even when $`p=q`$. This is because there are elements in $`H`$ such as $`P=Q^1=\mathrm{diag}(i,1,\mathrm{},1)`$ which change the sign of a single eigenvalue, in this case $`\lambda _1`$. Consistent with this, such elements have $`\mathrm{det}(P)\mathrm{det}(Q^1)1`$.
The next example is $`G/H=SU(n)/SO(n)`$. The subspace $`𝐡`$ consists of real, antisymmetric matrices, while $`𝐤`$ consists of imaginary, symmetric, traceless matrices, and its natural CSA coincides with the standard choice for $`su(n)`$. Invariance under $`H=SO(n)`$ means precisely that the eigenvalues can be permuted in any desired way. The set of invariants is of type A, the same as those for $`𝐠=su(n)`$, and they clearly descend from these by restriction to $`𝐤`$.
Finally, consider $`G/H=SU(2n)/Sp(2n)`$. We can choose the subspaces $`𝐡`$ and $`𝐤`$ to consist of matrices with the block forms $`\left(\begin{array}{cc}A& B\\ B^{}& A^{}\end{array}\right)`$ and $`\left(\begin{array}{cc}C& D\\ D^{}& C^{}\end{array}\right)`$ respectively, where $`A`$ lives in $`u(n)`$, $`B`$ is complex and symmetric, $`C`$ lives in $`su(n)`$, and $`D`$ is complex and antisymmetric. The CSA elements are $`X=\left(\begin{array}{cc}Y& 0\\ 0& Y\end{array}\right)`$ where $`Y`$ is any generator in the standard CSA of $`su(n)`$. The set of primitive invariants is once again A-type, matching those of $`su(n)`$. Unlike in the previous example, however, they arise here as restrictions of tensors on $`𝐠=su(2n)`$.
## 4 Poisson brackets of currents and charges
Returning now to our treatment of the $`G/H`$ sigma-model, we are interested in the classical Poisson brackets (PBs) of the currents $`j_\mu ^a`$ and $`J_\mu ^a`$. A convenient way to calculate these directly, without first finding PBs for the underlying fields $`g`$, is to use the approach of , which is easily adapted to the present situation. We take $`j_1^a`$ as an independent canonical coordinate and view $`j_0^a`$ as a non-local function of it by means of the identity $`_0j_1_1j_0+[j_0,j_1]=0`$, which implies
$$j_0=𝒟^1(_0j_1)\mathrm{where}𝒟X_1X+[j_1,X].$$
(4.1)
(We must assume appropriate boundary conditions on the fields, to allow the definition and manipulation of inverse differential operators.) In addition, the field $`A_1^{\widehat{\alpha }}`$ can be eliminated from the lagrangian immediately by its equation of motion, $`J_1^{\widehat{\alpha }}=j_1^{\widehat{\alpha }}A_1^{\widehat{\alpha }}=0`$. On making these replacements, the lagrangian becomes
$$=\frac{1}{2}(𝒟^1_0j_1A_0)^a(𝒟^1_0j_1A_0)^a\frac{1}{2}j_1^\alpha j_1^\alpha .$$
(4.2)
The momentum conjugate to $`j_1^a`$ is
$$\pi ^a=(𝒟^2_0j_1𝒟^1A_0)^a$$
(4.3)
and we impose canonical, equal-time PBs
$$\{j_1^a(x),\pi ^b(y)\}=\delta ^{ab}\delta (xy).$$
(4.4)
The resulting Hamiltonian density is
$$=\frac{1}{2}J_0^aJ_0^a+A_0^{\widehat{\alpha }}J_0^{\widehat{\alpha }}+\frac{1}{2}j_1^\alpha j_1^\alpha $$
(4.5)
where
$$J_0^a=𝒟\pi ^a.$$
(4.6)
The remaining, time-like component of the gauge field $`A_0^{\widehat{\alpha }}`$ acts as a Lagrange multiplier, imposing the constraint
$$J_0^{\widehat{\alpha }}0.$$
(4.7)
(We have taken a well-known short-cut by applying Dirac’s procedure without introducing momenta conjugate to the gauge fields, which play the role of Lagrange multipliers.)
To summarize: the independent canonical variables are $`j_1^a`$ and $`\pi ^a`$, obeying (4.4). In terms of these, $`j_0^a`$ is defined by (4.1); $`J_1^{\widehat{\alpha }}=0`$; $`J_1^\alpha =j_1^\alpha `$; and $`J_0^a`$ is defined by (4.6). The system is governed by the Hamiltonian density in (4.5), together with the constraint (4.7). From (4.4) it is straightforward to calculate the equal time PBs
$`\{J_0^a(x),J_0^b(y)\}`$ $`=`$ $`f^{abc}J_0^c(x)\delta (xy)`$
$`\{J_0^a(x),j_1^b(y)\}`$ $`=`$ $`f^{abc}j_1^c(x)\delta (xy)+\delta ^{ab}\delta ^{}(xy)`$ (4.8)
$`\{j_1^a(x),j_1^b(y)\}`$ $`=`$ $`0`$
which are the objects of central importance for us. The first of these implies that the constraints (4.7) are first-class, corresponding to the original $`H`$ gauge invariance. (They are analogous to the Gauss Law constraint arising in canonical treatments of electromagnetism.) It is also a simple consequence of (4.8) that the constraints weakly commute with the Hamiltonian, so they are preserved in time, and the canonical formulation has therefore been completed in a consistent fashion.
Let us now consider the canonical brackets of two conserved charges of type (2.3). We need not concern ourselves with gauge fixing the remaining constraints (4.7), because each current (2.2) commutes with them (each current is invariant under $`H`$-gauge transformations) and so the Dirac bracket of two conserved charges, obtained after imposing some gauge choice, is always identical to their Poisson bracket
$$\{q_s,q_r\}=\{𝑑xd_{\alpha _1\mathrm{}\alpha _{s+1}}^{(s+1)}J_+^{\alpha _1}(x)\mathrm{}J_+^{\alpha _{s+1}}(x),𝑑yd_{\beta _1\mathrm{}\beta _{r+1}}^{(r+1)}J_+^{\beta _1}(y)\mathrm{}J_+^{\beta _{r+1}}(y)\}$$
which can be calculated from (4.8). The terms in (4.8) containing $`\delta (xy)`$ do not contribute, by invariance of the $`d`$-tensors (the arguments are exactly similar to those given in for the PCM). The terms involving $`\delta ^{}(xy)`$, on the other hand, result in an integrand proportional to
$$d_{\alpha _1\mathrm{}\alpha _s\gamma }^{(s+1)}d_{\beta _1\mathrm{}\beta _r\gamma }^{(r+1)}J_+^{\alpha _1}\mathrm{}J_+^{\alpha _s}J_+^{\beta _1}\mathrm{}J_+^{\beta _{r1}}_1J_+^{\beta _r}.$$
The charges $`q_s`$ and $`q_r`$ will commute if and only if this integrand is a total derivative, which is true if and only if
$$d_{(\alpha _1\mathrm{}\alpha _s}^{(s+1)}{}_{}{}^{\gamma }d_{\beta _1\mathrm{}\beta _{r1})\beta _r}^{(r+1)}{}_{}{}^{\gamma }=d_{(\alpha _1\mathrm{}\alpha _s}^{(s+1)}{}_{}{}^{\gamma }d_{\beta _1\mathrm{}\beta _r)}^{(r+1)}{}_{}{}^{\gamma }.$$
(4.9)
It was shown in (see also eqn. (2.39) of ) that there exist tensors $`k^{(s+1)}`$ for each classical Lie group $`G`$ with the property that
$$k_{(a_1\mathrm{}a_s}^{(s+1)}{}_{}{}^{c}k_{b_1\mathrm{}b_{r1})b_r}^{(r+1)}{}_{}{}^{c}=k_{(a_1\mathrm{}a_s}^{(s+1)}{}_{}{}^{c}k_{b_1\mathrm{}b_r)}^{(r+1)}{}_{}{}^{c}.$$
(4.10)
This is exactly what is required to ensure commuting charges in the PCM based on $`G`$, and these tensors exist whenever $`s`$ is an exponent of $`G`$ modulo $`h`$. An obvious possibility is to choose the same tensors $`k^{(s+1)}`$ for $`G/H`$ as for $`G`$. We can certainly restrict the free indices in (4.10) from $`𝐠`$ to the subspace $`𝐤`$, but it is not clear that we can restrict the repeated index from $`c`$ in (4.10) to $`\gamma `$, so as to arrive at (4.9). Such a restriction is allowed in all cases of interest, however, by virtue of the property
$$d_{\alpha _1\mathrm{}\alpha _s\alpha _{s+1}}^{(s+1)}0d_{\alpha _1\mathrm{}\alpha _s\widehat{\gamma }}^{(s+1)}=0.$$
(4.11)
We shall prove in the next section that this holds for any invariant tensor on a classical symmetric space.
## 5 More properties of invariant tensors on $`G/H`$
Let us return to the definition of $`G/H`$ via a $`𝐙_2`$ grading of $`𝐠`$. An equivalent statement is that there is a Lie algebra automorphism $`\sigma `$ of $`𝐠`$ with $`\sigma ^2=1`$. The subspaces $`𝐡`$ and $`𝐤`$ are the eigenspaces of $`\sigma `$ with eigenvalues $`+1`$ and $`1`$ respectively.
An invariant tensor $`d^{(m)}`$ on $`𝐠`$ is not necessarily invariant under $`\sigma `$ (unless $`\sigma `$ is an inner automorphism of $`𝐠`$), but let us assume for the moment that $`\sigma (d)=d`$. Since elements of $`𝐤`$ change sign under $`\sigma `$, it follows that if the degree $`m`$ is odd, then $`d^{(m)}`$ must vanish when restricted to $`𝐤`$. If $`m`$ is even, $`d^{(m)}`$ need not vanish on $`𝐤`$, but then it must satisfy $`d_{\alpha _1\mathrm{}\alpha _{m1}\widehat{\gamma }}^{(m)}=0`$, since $`m1`$ is odd.
Now suppose that $`d`$ is not invariant under $`\sigma `$. It turns out that in all such cases, $`d`$ is instead invariant under the map $`\stackrel{~}{\sigma }(X)\sigma (X)`$, as we shall see below. The map $`\stackrel{~}{\sigma }`$ is not an automorphism of $`𝐠`$, and it has eigenspaces $`𝐡`$ and $`𝐤`$ with the reversed eigenvalues $`1`$ and $`+1`$ respectively. Requiring that $`d^{(m)}`$ be invariant under $`\stackrel{~}{\sigma }`$ and that it be non-zero on $`𝐤`$ does not result in any restriction on the degree $`m`$. However, since $`\stackrel{~}{\sigma }`$ reverses the sign of each element in $`𝐡`$, invariance of $`d`$ under this map does imply that $`d_{\alpha _1\mathrm{}\alpha _{m1}\widehat{\gamma }}^{(m)}=0`$, as required.
We conclude that (4.11) holds for every tensor $`d`$ obeying either $`\sigma (d)=d`$, or $`\stackrel{~}{\sigma }(d)=d`$. To check that this exhausts all possibilities, we can consider the simple explicit forms for $`\sigma `$ that are available for each of the classical symmetric spaces .
For the three families of Grassmannian symmetric spaces as well as for $`SO(2n)/U(n)`$ and $`Sp(2n)/U(n)`$, we can take $`\sigma `$ to be a map on $`𝐠`$ of the form
$$\sigma (X)=MXM^1,$$
for some matrix $`M`$. (Once again, this is not necessarily an inner automorphism, because $`M`$ may not belong to $`G`$.) Clearly this implies $`\sigma (d)=d`$ whenever $`d`$ is a symmetrized trace. The only primitive invariant not of this type is the (restriction of) the Pfaffian for $`SO(2n)/SO(n)\times SO(n)`$. In this case $`\sigma (d)=(detM)d=(1)^nd`$, since $`M`$ is diagonal with $`n`$ pairs of eigenvalues $`\pm 1`$. When $`n`$ is even, $`\sigma (d)=d`$, but when $`n`$ is odd, the degree of $`d`$ is also odd, and then $`\stackrel{~}{\sigma }(d)=d`$, as claimed.<sup>4</sup><sup>4</sup>4 These remarks are consistent with the patterns of invariants given in earlier sections: some cosets always have $`\sigma (d)=d`$, and their invariants were already found to have even degrees.
Finally we consider the families $`SU(n)/SO(n)`$ and $`SU(2n)/Sp(2n)`$ for which the automorphism of $`𝐠`$ can be written in the form
$$\sigma (X)=MX^{}M^1\stackrel{~}{\sigma }(X)=MX^TM^1,$$
since $`X`$ is anti-hermitian. For the first family $`M`$ is the identity, while for the second it is some standard symplectic structure. In either case it is clear that $`\stackrel{~}{\sigma }(d)=d`$ whenever $`d`$ is a symmetrized trace, and it follows that the same is true for any tensor $`d`$.
## 6 The classical spectrum of commuting charges
Our results concerning commuting charges in the $`G/H`$ SSM can now be summarized as follows. From each invariant tensor $`k^{(s+1)}`$ on $`𝐠`$ we obtain a conserved current in the $`G/H`$ model by restricting the tensor to $`𝐤`$, leading to a conserved charge of spin $`s`$. These conserved charges always commute with one another, by virtue of the arguments given in sections 4 and 5 above. The final issue we must settle is precisely which tensors $`k^{(s+1)}`$ are non-zero when restricted to $`𝐤`$, so as to give non-trivial conserved charges for the $`G/H`$ SSM.
Let us first recall what happens for a group $`G`$ . The tensors $`k^{(s+1)}`$ provide a set of primitive invariants when $`s`$ runs over the exponents of $`G`$. These primitive invariants have $`s<h`$, the Coxeter number of $`G`$. For larger values of $`s`$, the tensors $`k^{(s+1)}`$ are not primitive, but they are still non-zero (yielding non-trivial commuting charges) precisely when $`s`$ is equal to an exponent of $`G`$ modulo $`h`$.
We have discussed in detail in sections 3 and 4 the patterns of invariants for each $`G/H`$, and how these are obtained from invariants on $`G`$. We can certainly choose our primitive tensors on $`G/H`$ from amongst the $`k^{(s+1)}`$ required for commuting charges. We have already defined the corresponding values of $`s`$ to be the exponents of the symmetric space $`G/H`$. Similarly, we now define an integer $`h`$ for each $`G/H`$, with values specified in the table; this will play the role of the Coxeter number.<sup>5</sup><sup>5</sup>5The discussion in section 3 effectively associates a Weyl group to each symmetric space, and our definition of $`h`$ coincides with the standard one for such groups .
With these definitions, our final result can be stated very simply: $`k^{(s+1)}`$ is non-zero on $`𝐤`$ only when $`s`$ is equal to an exponent of $`G/H`$ modulo $`h`$. Thus, just as for groups, there are commuting charges associated with primitive invariants and hence exponents $`s<h`$, and also with non-primitive, but non-vanishing invariants for values of $`s>h`$ which repeat in families modulo $`h`$.
In some cases this result is automatic, because the values of $`s`$ run over all the positive odd integers, which trivially repeat modulo an even value of $`h`$. In other instances the pattern of spins is more involved, but the result for $`G/H`$ still follows directly from the corresponding statement for the group $`G`$. However, there is one family, namely $`SU(2n)/Sp(2n)`$, which requires special care. For these spaces we have defined $`h=n`$, and so the validity of our result depends, in particular, upon the fact that the tensors $`k^{(pn+1)}`$ should vanish when restricted to $`𝐤`$ for any integer $`p`$. When $`p`$ is even we know that this tensor actually vanishes on $`𝐠`$, but when $`p`$ is odd it is non-zero on $`𝐠`$ and it is not immediately clear why it should vanish on $`𝐤`$.
To understand how this comes about we need to look more closely at the definition of the $`k`$-tensors for the algebras $`su(n)`$, and to make clear which algebra we are talking about we shall write $`k_{su(n)}^{(m)}`$. Although these objects depend upon the value of $`n`$, they do so in a rather simple way, as revealed by the formula :
$$k_{su(n)}^{(m)}(X)=𝒜^{(m)}(\frac{1}{n}\mathrm{Tr}(X^m),\frac{1}{n}\mathrm{Tr}(X^{m1}),\mathrm{},\frac{1}{n}\mathrm{Tr}(X^2)).$$
The polynomials $`𝒜^{(m)}`$ are in fact characterized by the property $`k_{su(n)}^{(n+1)}=0`$, and it can be shown that this implies $`k_{su(n)}^{(pn+1)}=0`$ for any integer $`p`$ .
The coset $`SU(2n)/Sp(2n)`$ was one of the examples discussed earlier in section 4. The CSA can be chosen to consist of matrices of the form $`X=\left(\begin{array}{cc}Y& 0\\ 0& Y\end{array}\right)`$ where $`Y`$ is diagonal, traceless, and pure-imaginary, and hence belongs to the CSA of $`su(n)`$. Now clearly
$`k_{su(2n)}^{(m)}(X)`$ $`=`$ $`𝒜^{(m)}({\displaystyle \frac{1}{2n}}\mathrm{Tr}(X^m),{\displaystyle \frac{1}{2n}}\mathrm{Tr}(X^{m1}),\mathrm{},{\displaystyle \frac{1}{2n}}\mathrm{Tr}(X^2))`$
$`=`$ $`𝒜^{(m)}({\displaystyle \frac{1}{n}}\mathrm{Tr}(Y^m),{\displaystyle \frac{1}{n}}\mathrm{Tr}(Y^{m1}),\mathrm{},{\displaystyle \frac{1}{n}}\mathrm{Tr}(Y^2))`$
$`=`$ $`k_{su(n)}^{(m)}(Y).`$
The last expression vanishes when $`m=pn+1`$ for any integer $`p`$, as mentioned above. Since any invariant tensor is determined by its values on a CSA, we conclude that $`k_{su(2n)}^{(pn+1)}`$ indeed vanishes when restricted to $`𝐤`$, as claimed.
## 7 Comments
The results of this paper reinforce the elegant mathematical structure underlying the classical integrability of sigma-models on symmetric spaces. It would be interesting to investigate whether the charges we have constructed might be related to those arising in , whose properties are rather more mysterious. There are also a number of other directions for future work. Although we have considered only classical groups, we expect similar results to hold for the exceptional groups and their symmetric spaces. At the quantum level, integrability is believed to depend upon $`H`$ being simple, and for these cases exact S-matrices have been proposed . It would be interesting to examine whether quantum-mechanical survival of our charges is consistent with these proposals, in light of the treatment of affine Toda theories in . Finally, one could extend the results of to supersymmetric $`G/H`$ models, with the added incentive that these are believed to be quantum-integrable for any symmetric space .
Acknowledgments. We thank Jose Azcárraga and Tony Sudbery for helpful discussions. JME is supported by NSF grant PHY98-02484 and by a PPARC Advanced Fellowship. AJM thanks the 1851 Royal Commission for a Research Fellowship. This work was also supported in part by PPARC under the SPG grant PPA/G/S/1998/00613.
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# Massive chiral random matrix ensembles at 𝛽 = 1 & 4 : QCD Dirac operator spectra
## I introduction
Random matrix theory of disordered Hamiltonians relies upon an ansatz that in the ergodic regime where the system size $`L`$ is much larger than the elastic mean free path but much smaller than the localization length, details of the Hamiltonian are lost except for its time-reversal and internal symmetries. This ansatz was materialized by Efetov who derived Wigner-Dyson statistics out of Anderson tight-binding model by retaining only the zero-momentum mode of his spectral non-linear $`\sigma `$ model (NL$`\sigma `$M).
By taking the flavor symmetry among quarks into consideration as an additional internal symmetry, Verbaarschot and collaborators have reinterpreted this wisdom in the context of quantum chromodynamics (QCD). In this context, the spectral NL$`\sigma `$M is a supersymmetric extension of the conventional NL$`\sigma `$M over the coset manifold associated with the chiral symmetry breaking :
$`Z^{(2)}(\theta ;M)={\displaystyle _{SU(N_f)}}𝑑U\mathrm{exp}(\mathrm{Re}\mathrm{tr}\mathrm{e}^{i\theta /N_f}MU^{}),`$ (2)
$`Z^{(1)}(\theta ;M)={\displaystyle _{SU(2N_f)/Sp(2N_f)}}𝑑U\mathrm{exp}(\mathrm{Re}\mathrm{tr}\mathrm{e}^{i\theta /N_f}MU𝐉U^T/2),`$ (3)
$`Z^{(4)}(\theta ;M)={\displaystyle _{SU(N_f)/SO(N_f)}}𝑑U\mathrm{exp}(\mathrm{Re}\mathrm{tr}\mathrm{e}^{i\theta /(N_cN_f)}MUU^T),`$ (4)
after retaining only the zero mode. Here the superscripts (2, 1, 4) of Dyson indices $`\beta `$ refer to the anti-unitary symmetry of Euclidean Dirac operators (which are considered as stochastic Hamiltonians) in three classes of QCD :
$`\beta =2`$ $`:`$ $`N_c3,N_f\text{fundamental Dirac fermions},`$ (5)
$`\beta =1`$ $`:`$ $`N_c=2,N_f\text{fundamental Dirac fermions},`$ (6)
$`\beta =4`$ $`:`$ $`N_c2,N_f\text{adjoint Majorana fermions}.`$ (7)
The rescaled quark mass matrices $`M`$ are defined as
$`M=\mathrm{diag}(\mu _1,\mathrm{},\mu _{N_f})(\beta =2,4),`$ (8)
$`M=\mathrm{diag}(\mu _1,\mathrm{},\mu _{N_f})J(\beta =1),`$ (9)
$`\mu _i\mathrm{\Sigma }L^4m_i,𝐉=𝟙__𝕗𝕁,𝕁=\left(\begin{array}{cc}\mathrm{𝟘}& \mathrm{𝟙}\\ \mathrm{𝟙}& \mathrm{𝟘}\end{array}\right),`$ (12)
with a limit
$$L\mathrm{},m_i0,\mu _i:\text{fixed}$$
(13)
being assumed, $`\mathrm{\Sigma }`$ stands for the quark condensate in the chiral limit, and $`\theta `$ stands for the vacuum angle. By the same token as the spectral NL$`\sigma `$M of the tight-binding model was derived from a conventional random matrix ensembles , these NL$`\sigma `$Ms have an alternative derivation from chiral random matrix ensembles ($`\chi `$RMEs) :
$`Z_\nu ^{(\beta )}(\{m\})={\displaystyle _0^{2\pi }}{\displaystyle \frac{d\theta }{2\pi }}\mathrm{e}^{i\nu \theta }Z^{(\beta )}(\theta ;\{m\})`$ (14)
$`={\displaystyle 𝑑W\mathrm{e}^{\beta \mathrm{tr}V(W^{}W)}\underset{i=1}{\overset{N_f}{}}det\left(\begin{array}{cc}m_i& W\\ W^{}& m_i\end{array}\right)},`$ (17)
under a limit
$$N\mathrm{},m_i0,\mu _i=\pi \overline{\rho }(0)m_i:\text{fixed}.$$
(18)
Here the integrals are over complex, real, and quaternion real $`(N+\nu )\times N`$ matrices $`W`$ for $`\beta =2,1,4`$, respectively, and $`\overline{\rho }(0)`$ stands for the large-$`N`$ spectral density of the random matrix $`𝒟=\left(\genfrac{}{}{0pt}{}{0W}{W^{}0}\right)`$:
$$\overline{\rho }(\lambda )=\underset{N\mathrm{}}{lim}\mathrm{tr}\delta (\lambda i𝒟),$$
(19)
at the origin. It is understood for $`\beta =4`$ that twofold degenerated eigenvalues in the determinant are only counted once, and the topological charge $`\nu `$ is substituted by $`\nu N_c`$. These $`\chi `$RMEs are motivated by the microscopic theories (Euclidean QCD) on a lattice, with a crude simplification of replacing matrix elements of the anti-Hermitian Dirac operator $`/D=(_\mu +iA_\mu )\gamma _\mu `$ by random numbers $`𝒟`$ generated according to the weight $`\mathrm{e}^{\beta \mathrm{tr}V(W^{}W)}`$. Under this correspondence, the microscopic limit (18) is equivalent to Leutwyler-Smilga limit (13), since the size $`N`$ of the matrix $`W`$ is interpreted as the number of cites $`L^4`$ of the lattice on which QCD is discretized, and the Dirac spectral density at zero virtuality $`\overline{\rho }(0)`$ is related to the quark condensate by Banks-Casher relation $`\mathrm{\Sigma }=\pi \overline{\rho }(0)/L^4`$ . Using the $`\chi `$RME representation, various correlation functions of microscopically rescaled Dirac eigenvalues,
$$N\mathrm{},\lambda 0,\zeta =\pi \overline{\rho }(0)\lambda :\text{fixed},$$
(20)
have been computed for all three values of $`\beta `$ in the massless case $`\mu 0`$ . On the other hand, in a presence of finite $`\mu `$’s, Dirac eigenvalues comparable to $`\mu `$’s are expected to obey statistics that interpolate the chiral ($`\mu 0`$) and quenched limits ($`\mu \mathrm{}`$ or $`N_f=0`$). These novel spectral correlation functions have been analytically computed, until recently, solely for the chiral unitary ($`\beta =2`$) ensemble . Therefore we aim to treat the remaining cases, chiral orthogonal ($`\beta =1`$) and symplectic ($`\beta =4`$) ensembles with finite mass parameters, and compute Dirac eigenvalue correlation functions for these ensembles. We anticipate that advances in numerical simulations of lattice QCD with dynamical quarks will confirm our analytical results, Eqs.(II) and (III) below, in a foreseeable future.
This Article is organized as follows. In Sect. 2 we compute the correlation functions for the chiral orthogonal ensemble, by utilizing the quaternion determinant method developed in Ref.. In Sect. 3 we exhibit an explicit form of the correlation functions for the chiral symplectic ensemble, which was obtained by the Authors as a corollary to the computation of the partition function. In Appendix A we collect definitions related to a quaternion determinant. In Appendices B and C we present alternative expressions of the quaternion kernels for the orthogonal ensemble, and the symplectic ensembles with quadruply degenerated masses, respectively. It enables us to identify our results with those in a paper by Akemann and Kanzieper , which appeared after the Letter by the present Authors and computed a 1-level correlator with a single mass (for $`\beta =1`$) and $`p`$-level correlators with quadruply degenerate masses (for $`\beta =4`$).
## II orthogonal ensemble
We start by expressing the partition function (17) of the $`\chi `$RME in terms of eigenvalues $`x_i=\lambda _i^2`$ of the Wishart matrix $`W^{}W`$ (up to an overall constant):
$`Z_\nu ^{(\beta )}(\{m\})=({\displaystyle \underset{i=1}{\overset{\alpha }{}}}m_i^\nu )\mathrm{\Xi }_0(\{m\}),`$ (21)
$`\mathrm{\Xi }_0(\{m\})={\displaystyle \frac{1}{N!}}{\displaystyle _0^{\mathrm{}}}\mathrm{}{\displaystyle _0^{\mathrm{}}}{\displaystyle \underset{j=1}{\overset{N}{}}}dx_jP(\{x\};\{m\}),`$ (22)
$`P(\{x\};\{m\})={\displaystyle \underset{j=1}{\overset{N}{}}}\left(\mathrm{e}^{\beta V(x_j)}x_j^{\beta \frac{\nu +1}{2}1}{\displaystyle \underset{i=1}{\overset{\alpha }{}}}(x_j+m_i^2)\right)`$ (23)
$`\times {\displaystyle \underset{j>k}{\overset{N}{}}}|x_jx_k|^\beta .`$ (24)
The indices $`\beta `$ and $`\nu `$ are suppressed for simplicity. Since the partition function (17) is even under $`\nu \nu `$, we have set $`\nu `$ non-negative integer, without loss of generality. The $`p`$-level correlation function of the Wishart matrix $`W^{}W`$ is defined as
$`\sigma (x_1,\mathrm{},x_p;\{m\})={\displaystyle \frac{\mathrm{\Xi }_p(x_1,\mathrm{},x_p;\{m\})}{\mathrm{\Xi }_0(\{m\})}},`$ (25)
$`\mathrm{\Xi }_p(x_1,\mathrm{},x_p;\{m\})=`$ (26)
$`{\displaystyle \frac{1}{(Np)!}}{\displaystyle _0^{\mathrm{}}}\mathrm{}{\displaystyle _0^{\mathrm{}}}{\displaystyle \underset{j=p+1}{\overset{N}{}}}dx_jP(\{x\};\{m\}).`$ (27)
Then the $`p`$-level correlation function of the block off-diagonal Hermitian matrix $`i𝒟`$,
$$\rho (\lambda _1,\mathrm{},\lambda _p;\{m\})=\underset{k=1}{\overset{p}{}}\mathrm{tr}\delta (\lambda _ki𝒟),$$
(28)
is expressed in terms of $`\sigma `$ multiplied by the Jacobian of the transformation $`\lambda x=\lambda ^2`$ :
$$\rho (\lambda _1,\mathrm{},\lambda _p;\{m\})=2^p\underset{j=1}{\overset{p}{}}|\lambda _j|\sigma (\lambda _1^2,\mathrm{},\lambda _p^2;\{m\}).$$
(29)
As the universality of correlation functions of the unitary ensemble in the microscopic limit (20) is known to inherit to those of orthogonal and symplectic ensembles , it suffices to concentrate on Laguerre ensembles, $`V(x)=x`$. This leads to Wigner’s semi-circle law
$$\overline{\rho }(\lambda )=\frac{2}{\pi }\sqrt{2N\lambda ^2}.$$
(30)
Now we concentrate on $`\beta =1`$, and define new variables $`z_j`$ as
$$z_j=\{\begin{array}{cc}\hfill m_j^2(0),& j=1,\mathrm{},\alpha ,\hfill \\ \hfill x_{j\alpha }(0),& j=\alpha +1,\mathrm{},\alpha +N.\hfill \end{array}$$
(31)
Then the multiple integral (27) is expressed as
$`\mathrm{\Xi }_p(z_1,\mathrm{},z_{\alpha +p})={\displaystyle \frac{1}{_{j=1}^\alpha \sqrt{w(z_j)}_{j>k}^\alpha z_jz_k}}\times `$ (32)
$`{\displaystyle \frac{1}{(Np)!}}{\displaystyle _0^{\mathrm{}}}\mathrm{}{\displaystyle _0^{\mathrm{}}}{\displaystyle \underset{j=\alpha +p+1}{\overset{\alpha +N}{}}}dz_j{\displaystyle \underset{j=1}{\overset{\alpha +N}{}}}\sqrt{w(z_j)}{\displaystyle \underset{j>k}{\overset{\alpha +N}{}}}z_jz_k,`$ (33)
where
$$w(z)=|z|^{\nu 1}\mathrm{e}^{2z}.$$
(34)
Eq.(32) resembles an $`(\alpha +p)`$-level correlation function of the conventional (massless) Laguerre ensemble with $`\alpha +N`$ levels. However, conventionally the levels $`z_1,\mathrm{},z_{\alpha +p}`$ are all positive, while in the present case some of them ($`z_1,\mathrm{},z_\alpha `$) are negative. We carefully incorporate this fact into the following evaluation.
Let us denote the integrand in Eq.(32) as
$$p(z_1,\mathrm{},z_\gamma )=\underset{j=1}{\overset{\gamma }{}}\sqrt{w(z_j)}\underset{j>k}{\overset{\gamma }{}}z_jz_k,$$
(35)
with $`\gamma \alpha +N`$. It can be readily seen that
$`p(z_1,\mathrm{},z_\gamma )={\displaystyle \underset{j=1}{\overset{\gamma }{}}}\sqrt{w(z_j)}{\displaystyle \underset{j>k}{\overset{\gamma }{}}}(z_jz_k)\times `$ (36)
$`\{\begin{array}{cc}\mathrm{Pf}[\mathrm{sgn}(z_kz_j)]_{j,k=1,\mathrm{},\gamma }\hfill & (\gamma :\text{even})\hfill \\ \mathrm{Pf}\left[\begin{array}{cc}[\mathrm{sgn}(z_kz_j)]_{j,k=1,\mathrm{},\gamma }\hfill & [g_j]_{j=1,\mathrm{},\gamma }\hfill \\ \left[g_k\right]_{k=1,\mathrm{},\gamma }\hfill & 0\hfill \end{array}\right]\hfill & (\gamma :\text{odd}),\hfill \end{array}`$ (41)
with $`g_j=g_k=1`$. The Pfaffians in the above can be represented as quaternion determinants . In doing so, we need to introduce monic skew-orthogonal polynomials $`R_n(z)=z^n+\mathrm{}`$, which satisfy the skew-orthogonality relation:
$$R_{2n},R_{2m+1}_R=R_{2m+1},R_{2n}_R=r_n\delta _{nm},\text{others}=0,$$
(42)
where
$$f,g_R=_0^{\mathrm{}}dz\sqrt{w(z)}g(z)_0^zdz^{}\sqrt{w(z^{})}f(z^{})(fg),$$
(43)
and (integrated) ‘wave functions’,
$`\mathrm{\Psi }_n(z)=\sqrt{w(z)}R_n(z),`$ (45)
$`\mathrm{\Phi }_n(z)={\displaystyle _0^{\mathrm{}}}𝑑z^{}\mathrm{sgn}(zz^{})\sqrt{w(z^{})}R_n(z^{}).`$ (46)
Note that for negative $`z`$, $`\mathrm{\Phi }_n(z)`$ is a constant:
$$\mathrm{\Phi }_n(z<0)=\mathrm{\Phi }_n(0)s_n.$$
(47)
Now we present the following theorems:
Theorem 1
For even $`\gamma `$, we can express $`p(z_1,\mathrm{},z_\gamma )`$ as
$$p(z_1,\mathrm{},z_\gamma )=\left(\underset{j=0}{\overset{\gamma /21}{}}r_j\right)\mathrm{Tdet}[f(z_j,z_k)]_{j,k=1,\mathrm{},\gamma }.$$
(48)
The quaternion elements $`f(z,z^{})`$ are represented as
$$f(z,z^{})=\left[\begin{array}{cc}S(z,z^{})& I(z,z^{})\\ D(z,z^{})& S(z^{},z)\end{array}\right].$$
(49)
The functions $`S(z,z^{})`$, $`D(z,z^{})`$ and $`I(z,z^{})`$ are given by
$`S(z,z^{})`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\gamma /21}{}}}{\displaystyle \frac{\mathrm{\Phi }_{2n}(z)\mathrm{\Psi }_{2n+1}(z^{})\mathrm{\Phi }_{2n+1}(z)\mathrm{\Psi }_{2n}(z^{})}{r_n}},`$ (51)
$`D(z,z^{})`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\gamma /21}{}}}{\displaystyle \frac{\mathrm{\Psi }_{2n}(z)\mathrm{\Psi }_{2n+1}(z^{})\mathrm{\Psi }_{2n+1}(z)\mathrm{\Psi }_{2n}(z^{})}{r_n}},`$ (52)
$`I(z,z^{})`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\gamma /21}{}}}{\displaystyle \frac{\mathrm{\Phi }_{2n}(z)\mathrm{\Phi }_{2n+1}(z^{})\mathrm{\Phi }_{2n+1}(z)\mathrm{\Phi }_{2n}(z^{})}{r_n}}.`$ (53)
Theorem 2
For odd $`\gamma `$, we can express $`p(z_1,\mathrm{},z_\gamma )`$ as
$`p(z_1,\mathrm{},z_\gamma )=`$ (54)
$`\left({\displaystyle \underset{j=0}{\overset{[\gamma /2]1}{}}}r_j\right)s_{\gamma 1}\mathrm{Tdet}[f^{\mathrm{odd}}(z_j,z_k)]_{j,k=1,\mathrm{},\gamma }.`$ (55)
The quaternion elements are represented as
$$f^{\mathrm{odd}}(z,z^{})=\left[\begin{array}{cc}S^{\mathrm{odd}}(z,z^{})& I^{\mathrm{odd}}(z,z^{})\\ D^{\mathrm{odd}}(z,z^{})& S^{\mathrm{odd}}(z^{},z)\end{array}\right]$$
(56)
and $`s_n`$ is defined in Eq.(47). The functions $`S^{\mathrm{odd}}`$, $`D^{\mathrm{odd}}`$, and $`I^{\mathrm{odd}}`$ are given in terms of $`S`$, $`D`$, and $`I`$ in Theorem 1 according to
$`S^{\mathrm{odd}}(z,z^{})`$ $`=`$ $`S(z,z^{})|_{}+{\displaystyle \frac{\mathrm{\Psi }_{\gamma 1}(z^{})}{s_{\gamma 1}}},`$ (58)
$`D^{\mathrm{odd}}(z,z^{})`$ $`=`$ $`D(z,z^{})|_{},`$ (59)
$`I^{\mathrm{odd}}(z,z^{})`$ $`=`$ $`I(z,z^{})|_{}+{\displaystyle \frac{\mathrm{\Phi }_{\gamma 1}(z)\mathrm{\Phi }_{\gamma 1}(z^{})}{s_{\gamma 1}}}.`$ (60)
Here $``$ stands for a substitution
$$R_n(z)R_n(z)\frac{s_n}{s_{\gamma 1}}R_{\gamma 1}(z)$$
(61)
for $`n=0,\mathrm{},\gamma 2`$, associated with a change in the upper limit of the sum
$$\gamma /21[\gamma /2]1.$$
(62)
Theorem 3
Let the quaternion elements $`q_{jk}`$ of a selfdual $`n\times n`$ matrix $`Q_n`$ depend on $`n`$ real or complex variables $`z_1,\mathrm{},z_n`$ as
$$q_{jk}=f(z_j,z_k).$$
(63)
We assume that $`f(z,z^{})`$ satisfies the following conditions.
$`{\displaystyle f(z,z)𝑑\mu (z)}=c,`$ (64)
$`{\displaystyle f(z,z^{\prime \prime })f(z^{\prime \prime },z^{})𝑑\mu (z^{\prime \prime })}=f(z,z^{})+\lambda f(z,z^{})f(z,z^{})\lambda .`$ (65)
Here $`d\mu (z)`$ is a suitable measure, $`c`$ is a constant scalar, and $`\lambda `$ is a constant quaternion. Then we have
$$\mathrm{Tdet}Q_n𝑑\mu (z_n)=(cn+1)\mathrm{Tdet}Q_{n1},$$
(66)
where $`Q_{n1}`$ is the $`(n1)\times (n1)`$ matrix obtained by removing the row and the column which contain $`z_n`$. It is straightforward to show that the quaternion element $`f(z,z^{})`$ and $`f^{\mathrm{odd}}(z,z^{})`$ in Theorem 1 and Theorem 2 both satisfy the conditions imposed on $`f(z,z^{})`$ in Theorem 3 with $`d\mu (z)=w(z)dz`$. This means that we can write
$`\mathrm{\Xi }_p(z_1,\mathrm{},z_{\alpha +p})={\displaystyle \frac{_{j=0}^{[(\alpha +N)/2]1}r_j}{_{j=1}^\alpha \sqrt{w(z_j)}_{j>k}^\alpha z_jz_k}}\times `$ (67)
$`\{\begin{array}{ccc}& \mathrm{Tdet}[f(z_j,z_k)]_{j,k=1,\mathrm{},\alpha +p}\hfill & (\alpha +N:\text{even})\hfill \\ s_{\alpha +N1}\hfill & \mathrm{Tdet}[f^{\mathrm{odd}}(z_j,z_k)]_{j,k=1,\mathrm{},\alpha +p}\hfill & (\alpha +N:\text{odd})\hfill \end{array}.`$ (70)
Since the final result in the asymptotic limit $`N\mathrm{}`$ should be insensitive to the parity of $`N`$, we consider only even $`\alpha +N`$ henceforth. Then the $`p`$-level correlation function (25) is finally written as
$$\sigma (x_1,\mathrm{},x_p;m_1,\mathrm{},m_\alpha )=\frac{\mathrm{Tdet}[f(z_j,z_k)]_{j,k=1,\mathrm{},\alpha +p}}{\mathrm{Tdet}[f(z_j,z_k)]_{j,k=1,\mathrm{},\alpha }}.$$
(71)
Now we proceed to take the asymptotic limit of the correlation function, by making use of explicit forms for the skew-orthogonal polynomials and their norms associated with the weight (34) obtained by Nagao and Wadati :
$`R_{2n}(z)`$ $`=`$ $`{\displaystyle \frac{(2n)!}{2^{2n+1}}}{\displaystyle \frac{d}{dz}}L_{2n+1}^{(\nu 1)}(2z),`$ (72)
$`R_{2n+1}(z)`$ $`=`$ $`{\displaystyle \frac{(2n+1)!}{2^{2n+1}}}L_{2n+1}^{(\nu 1)}(2z)`$ (74)
$`{\displaystyle \frac{(2n)!}{2^{2n+2}}}(2n+\nu ){\displaystyle \frac{d}{dz}}L_{2n}^{(\nu 1)}(2z),`$
$`r_n`$ $`=`$ $`2^{4n\nu }(2n)!(2n+\nu )!,`$ (75)
expressed in terms of the Laguerre polynomials
$$L_n^{(a)}(z)=\frac{z^a\mathrm{e}^z}{n!}\frac{d^n}{dz^n}(\mathrm{e}^zz^{n+a}).$$
(76)
We need to evaluate the local asymptotics of the quaternion function $`f(z,z^{})`$, whose constituent Laguerre polynomials tend to
$$L_n^{(a)}(z)\{\begin{array}{cc}\left(n/z\right)^{a/2}J_a(2\sqrt{nz})\hfill & (z>0)\hfill \\ \left(n/|z|\right)^{a/2}I_a(2\sqrt{n|z|})\hfill & (z<0)\hfill \end{array},$$
(77)
as $`n\mathrm{}`$, $`z0`$, with $`nz:`$ fixed. Three cases should be considered separately:
$`(a)z,z^{}>0,(b)z<0,z^{}>0,(c)z,z^{}<0`$
(the case $`z>0`$, $`z^{}<0`$ is unnecessary because of the selfduality $`f(z,z^{})=\widehat{f}(z^{},z)`$).
$`(a)`$ $`z,z^{}>0`$
We define microscopic variables $`\zeta `$ and $`\zeta ^{}`$ as
$$z=\frac{\zeta ^2}{8N},z^{}=\frac{\zeta _{}^{}{}_{}{}^{2}}{8N},$$
(78)
according to Eq.(20) with $`\pi \overline{\rho }(0)=2\sqrt{2N}`$. If both arguments are positive, the asymptotic limit is known to be
$`S_{++}(\zeta ,\zeta ^{}){\displaystyle \frac{1}{8N}}S(z,z^{}){\displaystyle \frac{1}{4}}{\displaystyle _0^1}𝑑tt^2{\displaystyle _0^\zeta }𝑑s\left(sJ_{\nu 1}(ts){\displaystyle \frac{J_\nu (t\zeta ^{})}{\zeta ^{}}}J_\nu (ts)J_{\nu 1}(t\zeta ^{})\right)+{\displaystyle \frac{J_\nu (\zeta ^{})}{4\zeta ^{}}},`$ (80)
$`D_{++}(\zeta ,\zeta ^{}){\displaystyle \frac{1}{(8N)^2}}D(z,z^{}){\displaystyle \frac{1}{16}}{\displaystyle _0^1}𝑑tt^2\left(J_{\nu 1}(t\zeta ){\displaystyle \frac{J_\nu (t\zeta ^{})}{\zeta ^{}}}{\displaystyle \frac{J_\nu (t\zeta )}{\zeta }}J_{\nu 1}(t\zeta ^{})\right),`$ (81)
$`I_{++}(\zeta ,\zeta ^{})I(z,z^{}){\displaystyle _0^1}𝑑tt^2{\displaystyle _0^\zeta }𝑑u{\displaystyle _0^\zeta ^{}}𝑑v\left(uJ_{\nu 1}(tu)J_\nu (tv)J_\nu (tu)vJ_{\nu 1}(tv)\right){\displaystyle _\zeta ^\zeta ^{}}J_\nu (u)𝑑u\mathrm{sgn}(\zeta \zeta ^{}).`$ (82)
(83)
$`(b)`$ $`z<0,z^{}>0`$
We define microscopic variables $`\mu `$ and $`\zeta `$ as
$$z=\frac{\mu ^2}{8N},z^{}=\frac{\zeta ^2}{8N},$$
(84)
according to Eqs.(18) and (20). When one of the arguments is negative, the identity (47) should be taken into consideration. We find
$`S_+(\mu ,\zeta )=S_+(\zeta ){\displaystyle \frac{1}{8N}}S(z,z^{}){\displaystyle \frac{J_\nu (\zeta )}{4\zeta }},`$ (86)
$`S_+(\zeta ,\mu ){\displaystyle \frac{1}{8N}}S(z^{},z){\displaystyle \frac{1}{4}}{\displaystyle _0^1}𝑑tt^2{\displaystyle _0^\zeta }𝑑s\left(sJ_{\nu 1}(ts){\displaystyle \frac{I_\nu (t\mu )}{\mu }}J_\nu (ts)I_{\nu 1}(t\mu )\right)+{\displaystyle \frac{I_\nu (\mu )}{4\mu }},`$ (87)
$`D_+(\mu ,\zeta ){\displaystyle \frac{1}{(8N)^2}}D(z,z^{}){\displaystyle \frac{1}{16}}{\displaystyle _0^1}𝑑tt^2\left(I_{\nu 1}(t\mu ){\displaystyle \frac{J_\nu (t\zeta )}{\zeta }}{\displaystyle \frac{I_\nu (t\mu )}{\mu }}J_{\nu 1}(t\zeta )\right),`$ (88)
$`I_+(\mu ,\zeta )=I_+(\zeta )I(z,z^{}){\displaystyle _0^\zeta }J_\nu (u)𝑑u+1.`$ (89)
$`(c)`$ $`z,z^{}<0`$
We define microscopic variables $`\mu `$ and $`\zeta `$ as
$$z=\frac{\mu ^2}{8N},z^{}=\frac{\mu _{}^{}{}_{}{}^{2}}{8N}.$$
(90)
We can readily derive
$`S_{}(\mu ,\mu ^{})=S_{}(\mu ^{}){\displaystyle \frac{1}{8N}}S(z,z^{}){\displaystyle \frac{I_\nu (\mu ^{})}{4\mu ^{}}},`$ (92)
$`D_{}(\mu ,\mu ^{}){\displaystyle \frac{1}{(8N)^2}}D(z,z^{}){\displaystyle \frac{1}{16}}{\displaystyle _0^1}𝑑tt^2\left(I_{\nu 1}(t\mu ){\displaystyle \frac{I_\nu (t\mu ^{})}{\mu ^{}}}{\displaystyle \frac{I_\nu (t\mu )}{\mu }}I_{\nu 1}(t\mu ^{})\right),`$ (93)
$`I_{}(\mu ,\mu ^{})I(z,z^{})\mathrm{sgn}(\mu \mu ^{}).`$ (94)
In Eqs.(86), (89), and (92), we have introduced symbols with one sign subscript (e.g. $`S_+`$) in order to indicate that they depend only on the second arguments of those with two sign subscripts (e.g. $`S_+`$).
The scaled correlation function $`\rho _s`$ is defined as
$$\rho _s(\zeta _1,\mathrm{},\zeta _p;\mu _1,\mathrm{},\mu _\alpha )=\frac{1}{(8N)^p}\rho (\frac{\zeta _1}{\sqrt{8N}},\mathrm{},\frac{\zeta _p}{\sqrt{8N}};\frac{\mu _1}{\sqrt{8N}},\mathrm{},\frac{\mu _\alpha }{\sqrt{8N}}).$$
(95)
By making use of Eqs.(29), (71), and Dyson’s equality (A16), we finally obtain
$`\rho _s(\zeta _1,\mathrm{},\zeta _p;\mu _1,\mathrm{},\mu _\alpha )`$ $`=(1)^{p(p1)/2}2^p{\displaystyle \underset{k=1}{\overset{p}{}}}|\zeta _k|{\displaystyle \frac{\mathrm{Pf}\left[\begin{array}{cccc}\hfill I_{}& \hfill S_{}& \hfill I_+& \hfill S_+\\ \hfill S_{}^T& \hfill D_{}& \hfill S_+^T& \hfill D_+\\ \hfill I_+^T& \hfill S_+& \hfill I_{++}& \hfill S_{++}\\ \hfill S_+^T& \hfill D_+^T& \hfill S_{++}^T& \hfill D_{++}\end{array}\right]}{\mathrm{Pf}\left[\begin{array}{cc}\hfill I_{}& \hfill S_{}\\ \hfill S_{}^T& \hfill D_{}\end{array}\right]}}`$ (114)
$`=(1)^{p(p1)/2}2^p{\displaystyle \underset{k=1}{\overset{p}{}}}|\zeta _k|{\displaystyle \frac{\mathrm{Pf}\left[\begin{array}{ccc}\hfill D_{}& \hfill S_+^T& \hfill D_+\\ \hfill S_+& \hfill I_{++}& \hfill S_{++}\\ \hfill D_+^T& \hfill S_{++}^T& \hfill D_{++}\end{array}\right]}{\mathrm{Pf}\left[D_{}\right]}}(\alpha :\text{even})`$
$`=(1)^{p(p1)/2}2^p{\displaystyle \underset{k=1}{\overset{p}{}}}|\zeta _k|{\displaystyle \frac{\mathrm{Pf}\left[\begin{array}{cccc}\hfill 0& \hfill S_{}& \hfill I_+& \hfill S_+\\ \hfill S_{}^T& \hfill D_{}& \hfill S_+^T& \hfill D_+\\ \hfill I_+^T& \hfill S_+& \hfill I_{++}& \hfill S_{++}\\ \hfill S_+^T& \hfill D_+^T& \hfill S_{++}^T& \hfill D_{++}\end{array}\right]}{\mathrm{Pf}\left[\begin{array}{cc}\hfill 0& \hfill S_{}\\ \hfill S_{}^T& \hfill D_{}\end{array}\right]}}(\alpha :\text{odd}).`$
The elements of the matrices $`S_{ϵϵ^{}}`$, $`D_{ϵϵ^{}}`$, $`I_{ϵϵ^{}}`$ and the row vectors $`S_ϵ`$, $`I_ϵ`$ ($`ϵ,ϵ^{}=+,`$) in the above are defined as
$$(S_{++})_k\mathrm{}=S_{++}(\zeta _k,\zeta _{\mathrm{}}),(S_+)_i\mathrm{}=S_+(\mu _i,\zeta _{\mathrm{}}),(S_{})_{ij}=S_{}(\mu _i,\mu _j),(S_+)_{\mathrm{}}=S_+(\zeta _{\mathrm{}}),(S_{})_j=S_{}(\mu _j),\text{etc.},$$
(115)
where the subscripts take their values in $`i,j=1,\mathrm{},\alpha `$ and $`k,\mathrm{}=1,\mathrm{},p`$. In the last two lines we have exploited a Pfaffian identity that holds for antisymmetric matrices $`A`$, $`B`$, and a row vector $`v`$:
$`\mathrm{Pf}\left[\begin{array}{cc}A& \begin{array}{c}v\\ \mathrm{}\\ v\end{array}\\ & \\ v^T\mathrm{}v^T& B\end{array}\right]`$ $`=`$ $`\mathrm{Pf}[A]\mathrm{Pf}[B](\mathrm{rank}(A),\mathrm{rank}(B):\text{even})`$ (122)
$`=`$ $`\mathrm{Pf}\left[\begin{array}{cc}A& \begin{array}{c}1\\ \mathrm{}\\ 1\end{array}\\ & \\ 1\mathrm{}1& 0\end{array}\right]\mathrm{Pf}\left[\begin{array}{cc}0& v\\ v^T& B\end{array}\right](\mathrm{rank}(A),\mathrm{rank}(B):\text{odd}).`$ (130)
In a special case $`p=\alpha =1`$, the expression (114) reduces to Akemann and Kanzieper’s recent result (see Appendix B).
In the quenched limit $`\mu _1,\mathrm{},\mu _\alpha \mathrm{}`$ when the ratio of two Pfaffians is replaced by a minor $`\mathrm{Pf}\left[\genfrac{}{}{0pt}{}{I_{++}S_{++}}{S_{++}^TD_{++}}\right]`$, the correlation function tends to that of Laguerre orthogonal ensemble computed by Nagao and Forrester, Eqs.(2.21), (2.18), (2.19), and (3.20) in Ref., with $`\nu =2a+1`$. By the same token, it satisfies a sequence
$$\rho _s(\{\zeta \};\mu _1,\mathrm{},\mu _\alpha )\stackrel{\mu _\alpha \mathrm{}}{}\rho _s(\{\zeta \};\mu _1,\mathrm{},\mu _{\alpha 1})\stackrel{\mu _{\alpha 1}\mathrm{}}{}\rho _s(\{\zeta \};\mu _1,\mathrm{},\mu _{\alpha 2})\stackrel{\mu _{\alpha 2}\mathrm{}}{}\mathrm{},$$
(131)
as each of the masses are decoupled by sending to infinity. To illustrate this decoupling, we exhibit in Fig.1 a plot of the spectral density $`\rho _s(\zeta ;\mu )`$ ($`\nu =0,p=1,\alpha =1`$) that interpolates between known results for the chiral and quenched limits.
## III symplectic ensemble
Although correlation functions of the massive chiral symplectic ensemble have previously been computed by the Authors , we nevertheless present their explicit expressions for the sake of completeness. We concentrate on the case with an even $`N_f(2\alpha )`$ number of flavors and pairwise degenerated mass parameters, corresponding to adjoint Dirac fermions in the QCD context. The scaled $`p`$-level correlation functions, defined in Eq.(95), is expressed by construction as a ratio of partition functions with $`2\alpha `$ and $`2\alpha +4p`$ flavors ,
$$\rho _s(\zeta _1,\mathrm{},\zeta _p;\{\mu \})=C_{\alpha ,\nu }^{(p)}\underset{k>\mathrm{}}{\overset{p}{}}(\zeta _k^2\zeta _{\mathrm{}}^2)^4\underset{k=1}{\overset{p}{}}\left(|\zeta _k|^3\underset{i=1}{\overset{\alpha }{}}(\zeta _k^2+\mu _i^2)^2\right)\frac{Z_\nu ^{(4)}(\mu _1,\mu _1,\mathrm{},\mu _\alpha ,\mu _\alpha ,\stackrel{4}{\stackrel{}{i\zeta _1,\mathrm{},i\zeta _1}},\mathrm{},\stackrel{4}{\stackrel{}{i\zeta _p,\mathrm{},i\zeta _p}})}{Z_\nu ^{(4)}(\mu _1,\mu _1,\mathrm{},\mu _\alpha ,\mu _\alpha )}.$$
(132)
Here $`C_{\alpha ,\nu }^{(p)}`$ stands for a constant to be fixed below. Using an explicit form of the partition function, Eq.(31) of Ref., and taking confluent limits in $`\zeta _k`$’s, we obtain
$`\rho _s(\zeta _1,\mathrm{},\zeta _p;\{\mu \})`$ $`=`$ $`(1)^{p(p+1)/2}2^p{\displaystyle \underset{k=1}{\overset{p}{}}}|\zeta _k|{\displaystyle \frac{\mathrm{Pf}\left[\begin{array}{ccc}\hfill I_{}& \hfill I_+& \hfill S_+\\ \hfill I_+^T& \hfill I_{++}& \hfill S_{++}\\ \hfill S_+^T& \hfill S_{++}^T& \hfill D_{++}\end{array}\right]}{\mathrm{Pf}[I_{}]}}(\alpha :\text{even})`$ (137)
$`=`$ $`(1)^{p(p+1)/2}2^p{\displaystyle \underset{k=1}{\overset{p}{}}}|\zeta _k|{\displaystyle \frac{\mathrm{Pf}\left[\begin{array}{cccc}\hfill I_{}& \hfill Q_{}& \hfill I_+& \hfill S_+\\ \hfill Q_{}^T& \hfill 0& \hfill Q_+^T& \hfill P_+^T\\ \hfill I_+^T& \hfill Q_+& \hfill I_{++}& \hfill S_{++}\\ \hfill S_+^T& \hfill P_+& \hfill S_{++}^T& \hfill D_{++}\end{array}\right]}{\mathrm{Pf}\left[\begin{array}{cc}\hfill I_{}& \hfill Q_{}\\ \hfill Q_{}^T& \hfill 0\end{array}\right]}}(\alpha :\text{odd}).`$ (144)
The elements of the matrices $`S_{ϵϵ^{}}`$, $`D_{ϵϵ^{}}`$, $`I_{ϵϵ^{}}`$ and the column vectors $`Q_ϵ`$, $`P_ϵ`$ ($`ϵ,ϵ^{}=+,`$) in the above are defined as
$`(I_{})_{ij}=I_{}(\mu _i,\mu _j)\mu _i\mu _j{\displaystyle _0^1}𝑑tt{\displaystyle _0^1}𝑑u\left(I_{2\nu }(2t\mu _i)I_{2\nu }(2tu\mu _j)I_{2\nu }(2tu\mu _i)I_{2\nu }(2t\mu _j)\right),`$ (145)
$`(I_+)_i\mathrm{}=I_+(\mu _i,\zeta _{\mathrm{}})\mu _i\zeta _{\mathrm{}}{\displaystyle _0^1}𝑑tt{\displaystyle _0^1}𝑑u\left(I_{2\nu }(2t\mu _i)J_{2\nu }(2tu\zeta _{\mathrm{}})I_{2\nu }(2tu\mu _i)J_{2\nu }(2t\zeta _{\mathrm{}})\right),`$ (146)
$`(S_+)_i\mathrm{}=S_+(\mu _i,\zeta _{\mathrm{}})\mu _i{\displaystyle _0^1}𝑑tt^2{\displaystyle _0^1}𝑑u\left(I_{2\nu }(2t\mu _i)uJ_{2\nu +1}(2tu\zeta _{\mathrm{}})I_{2\nu }(2tu\mu _i)J_{2\nu +1}(2t\zeta _{\mathrm{}})\right),`$ (147)
$`(I_{++})_k\mathrm{}=I_{++}(\zeta _k,\zeta _{\mathrm{}})\zeta _k\zeta _{\mathrm{}}{\displaystyle _0^1}𝑑tt{\displaystyle _0^1}𝑑u\left(J_{2\nu }(2t\zeta _k)J_{2\nu }(2tu\zeta _{\mathrm{}})J_{2\nu }(2tu\zeta _k)J_{2\nu }(2t\zeta _{\mathrm{}})\right),`$ (148)
$`(S_{++})_k\mathrm{}=S_{++}(\zeta _k,\zeta _{\mathrm{}})\zeta _k{\displaystyle _0^1}𝑑tt^2{\displaystyle _0^1}𝑑u\left(J_{2\nu }(2t\zeta _k)uJ_{2\nu +1}(2tu\zeta _{\mathrm{}})J_{2\nu }(2tu\zeta _k)J_{2\nu +1}(2t\zeta _{\mathrm{}})\right),`$ (149)
$`(D_{++})_k\mathrm{}=D_{++}(\zeta _k,\zeta _{\mathrm{}}){\displaystyle _0^1}𝑑tt^3{\displaystyle _0^1}𝑑uu\left(J_{2\nu +1}(2t\zeta _k)J_{2\nu +1}(2tu\zeta _{\mathrm{}})J_{2\nu +1}(2tu\zeta _k)J_{2\nu +1}(2t\zeta _{\mathrm{}})\right),`$ (150)
$`(Q_{})_j=\mu _j{\displaystyle _0^1}𝑑tI_{2\nu }(2t\mu _j),(Q_+)_{\mathrm{}}=\zeta _{\mathrm{}}{\displaystyle _0^1}𝑑tJ_{2\nu }(2t\zeta _{\mathrm{}}),(P_+)_{\mathrm{}}={\displaystyle _0^1}𝑑ttJ_{2\nu +1}(2t\zeta _{\mathrm{}}),`$ (151)
where the subscripts take their values in $`i,j=1,\mathrm{},\alpha `$ and $`k,\mathrm{}=1,\mathrm{},p`$. The overall constant is determined as the above by requiring that in the quenched limit $`\mu _1,\mathrm{},\mu _\alpha \mathrm{}`$ when the ratio of two Pfaffians is replaced by a minor $`\mathrm{Pf}\left[\genfrac{}{}{0pt}{}{I_{++}S_{++}}{S_{++}^TD_{++}}\right]`$, the correlation function tends to that of Laguerre symplectic ensemble computed by Nagao and Forrester, Eqs.(2.27), (2.25), and (4.7$``$9) in Ref. (whose notations are related to ours via an unfolding change of variables
$`I_{++}(\zeta ,\zeta ^{})=I_4({\displaystyle \frac{\zeta ^2}{8N}},{\displaystyle \frac{\zeta _{}^{}{}_{}{}^{2}}{8N}}),`$ (152)
$`S_{++}(\zeta ,\zeta ^{})={\displaystyle \frac{1}{8N}}S_4({\displaystyle \frac{\zeta ^2}{8N}},{\displaystyle \frac{\zeta _{}^{}{}_{}{}^{2}}{8N}}),`$ (153)
$`D_{++}(\zeta ,\zeta ^{})={\displaystyle \frac{1}{(8N)^2}}D_4({\displaystyle \frac{\zeta ^2}{8N}},{\displaystyle \frac{\zeta _{}^{}{}_{}{}^{2}}{8N}}),`$ (154)
and $`\nu =2a1/2`$). It is easy to confirm that the correlation functions satisfies the decoupling sequence (131) as each of the masses are sent to infinity. To illustrate this decoupling, we exhibit in Fig.2 a plot of the spectral density $`\rho _s(\zeta ;\mu ,\mu )`$ ($`\nu =0,p=1,\alpha =1`$) that interpolates between known results for the chiral and quenched limits.
After the Authors announced the above formula in , Akemann and Kanzieper presented another form of the asymptotic correlations in a special case of quadruply degenerate masses. In Appendix C, we shall reproduce their result as a confluent limit of our formula.
###### Acknowledgements.
This work was supported in part (SMN) by JSPS Research Fellowships for Young Scientists, and by Grant-in-Aid No. 411044 from the Ministry of Education, Science, and Culture, Japan.
## A quaternion determinant
A quaternion is defined as a linear combination of four basic units $`\{1,e_1,e_2,e_3\}`$:
$$q=q_0+𝐪𝐞=q_0+q_1e_1+q_2e_2+q_3e_3.$$
(A1)
Here the coefficients $`q_0,q_1,q_2`$ and $`q_3`$ are real or complex numbers. The first part $`q_0`$ is called the scalar part of $`q`$. The quaternion basic units satisfy the multiplication laws
$`11=1,1e_j=e_j1=e_j,j=1,2,3,`$ (A2)
$`e_1^2=e_2^2=e_3^2=e_1e_2e_3=1.`$ (A3)
The multiplication is associative and in general not commutative. The dual $`\widehat{q}`$ of a quaternion $`q`$ is defined as
$$\widehat{q}=q_0𝐪𝐞.$$
(A4)
For a selfdual $`N\times N`$ matrix $`Q`$ with quaternion elements $`q_{jk}`$ has a dual matrix $`\widehat{Q}=[\widehat{q}_{kj}]`$. The quaternion units can be represented as $`2\times 2`$ matrices
$`1\left[\begin{array}{cc}1& 0\\ 0& 1\end{array}\right],e_1\left[\begin{array}{cc}0& 1\\ 1& 0\end{array}\right],`$ (A9)
$`e_2\left[\begin{array}{cc}0& i\\ i& 0\end{array}\right],e_3\left[\begin{array}{cc}i& 0\\ 0& i\end{array}\right].`$ (A14)
We define a quaternion determinant Tdet of a selfdual $`Q`$ (i.e., $`Q=\widehat{Q}`$) as
$$\mathrm{Tdet}Q=\underset{P}{}(1)^{Nl}\underset{1}{\overset{l}{}}(q_{ab}q_{bc}\mathrm{}q_{da})_0,$$
(A15)
where $`P`$ denotes any permutation of the indices $`(1,2,\mathrm{},N)`$ consisting of $`l`$ exclusive cycles of the form $`(abc\mathrm{}da)`$ and $`(1)^{Nl}`$ is the parity of $`P`$. The subscript $`0`$ means that the scalar part of the product is taken over each cycle. Note that a quaternion determinant of a selfdual quaternion matrix is always a scalar. The quaternion determinant can as well be represented by the $`2N\times 2N`$ representation $`C(Q)`$ :
$$\mathrm{Tdet}Q=\mathrm{Pf}[𝐉C(Q)],𝐉=𝟙_{}\left[\begin{array}{cc}\mathrm{𝟘}& \mathrm{𝟙}\\ \mathrm{𝟙}& \mathrm{𝟘}\end{array}\right].$$
(A16)
## B quaternion kernel for the orthogonal ensemble
Recently Akemann and Kanzieper derived the asymptotic correlation function in a special case $`p=\alpha =1`$. Though their result is clearly in agreement with that special case of ours, there is a difference in the appearance because they adopted alternative asymptotic formulas. Their formulas are based on an asymptotic relation for $`z,z^{}>0`$ and $`\nu 1`$ (see below):
$`S(z,z^{})`$ $``$ $`2N[2{\displaystyle _0^1}dttJ_{\nu 1}(t\zeta )J_{\nu 1}(t\zeta ^{})`$ (B2)
$`{\displaystyle \frac{J_\nu (\zeta ^{})}{\zeta ^{}}}({\displaystyle _0^\zeta }dsJ_{\nu 2}(s)1)],`$
where we have adopted the microscopic variables (78). This asymptotic relation was derived by Forrester, Nagao, and Honner in a study of parametric random matrix ensembles. Note that the first integral in the above is equal to the Bessel kernel,
$$\frac{\zeta J_\nu (\zeta )J_{\nu 1}(\zeta ^{})J_{\nu 1}(\zeta )\zeta ^{}J_\nu (\zeta ^{})}{\zeta ^2\zeta _{}^{}{}_{}{}^{2}}.$$
(B3)
From the derivation the equivalence of (80) and (B2) was well established. However, it is worth directly proving it here for an unambiguous identification.
Using the Bessel function identities, we can readily see that
$`t{\displaystyle _0^\zeta }𝑑ssJ_{\nu 1}(ts)=\zeta J_\nu (t\zeta )+(\nu 1){\displaystyle _0^\zeta }𝑑sJ_\nu (ts),`$ (B4)
$`t{\displaystyle _0^\zeta }𝑑sJ_\nu (ts)=J_{\nu 1}(t\zeta )+J_{\nu 1}(0)+(\nu 1){\displaystyle _0^\zeta }{\displaystyle \frac{ds}{s}}J_{\nu 1}(ts).`$ (B5)
Substitution of Eq.(B5) into (80) yields
$`S(z,z^{})`$ $``$ $`2N[{\displaystyle \frac{\zeta }{\zeta ^{}}}{\displaystyle _0^1}dttJ_\nu (t\zeta )J_\nu (t\zeta ^{})+{\displaystyle _0^1}dtt(J_{\nu 1}(t\zeta )J_{\nu 1}(0))J_{\nu 1}(t\zeta ^{})+{\displaystyle \frac{J_\nu (\zeta ^{})}{\zeta ^{}}}`$ (B7)
$`+{\displaystyle \frac{\nu 1}{\zeta ^{}}}{\displaystyle _0^1}dtt{\displaystyle _0^\zeta }dsJ_\nu (ts)J_\nu (t\zeta ^{})(\nu 1){\displaystyle _0^1}dtt{\displaystyle _0^\zeta }{\displaystyle \frac{ds}{s}}J_{\nu 1}(ts)J_{\nu 1}(t\zeta ^{})].`$
By partial integrations, we find
$$\frac{\zeta }{\zeta ^{}}_0^1𝑑ttJ_\nu (t\zeta )J_\nu (t\zeta ^{})=_0^1𝑑ttJ_{\nu 1}(t\zeta )J_{\nu 1}(t\zeta ^{})J_{\nu 1}(\zeta )\frac{J_\nu (\zeta ^{})}{\zeta ^{}}$$
(B8)
and
$$\frac{1}{\zeta ^{}}_0^1𝑑tt_0^\zeta 𝑑sJ_\nu (ts)J_\nu (t\zeta ^{})=_0^1𝑑tt_0^\zeta \frac{ds}{s}J_{\nu 1}(ts)J_{\nu 1}(t\zeta ^{})_0^\zeta \frac{ds}{s}J_{\nu 1}(s)\frac{J_\nu (\zeta ^{})}{\zeta ^{}}.$$
(B9)
We substitute (B8) and (B9) into (B7) and obtain
$`S(z,z^{})`$ (B10)
$`2N\left[2{\displaystyle _0^1}𝑑ttJ_{\nu 1}(t\zeta )J_{\nu 1}(t\zeta ^{})J_{\nu 1}(0){\displaystyle _0^1}𝑑ttJ_{\nu 1}(t\zeta ^{})\left(J_{\nu 1}(\zeta )+(\nu 1){\displaystyle _0^\zeta }{\displaystyle \frac{ds}{s}}J_{\nu 1}(s)1\right){\displaystyle \frac{J_\nu (\zeta ^{})}{\zeta ^{}}}\right].`$ (B11)
Again by a partial integration, we find
$$J_{\nu 1}(\zeta )+(\nu 1)_0^\zeta \frac{ds}{s}J_{\nu 1}(s)=J_{\nu 1}(0)+_0^\zeta 𝑑sJ_{\nu 2}(s)$$
(B12)
and thus arrive at the desired result (B2) provided that $`\nu `$ is not equal to $`1`$. Since $`J_{\nu 1}(0)=\delta _{\nu ,1}`$, the case $`\nu =1`$ is exceptional.
Similarly, for $`z<0,z^{}>0`$ and $`\nu 1`$, an alternative expression
$$S(z^{},z)2N\left[2_0^1𝑑ttJ_{\nu 1}(t\zeta )I_{\nu 1}(t\mu )\frac{I_\nu (\mu )}{\mu }\left(_0^\zeta 𝑑sJ_{\nu 2}(s)1\right)\right]$$
(B13)
in terms of the microscopic variables (84) is available.
## C quaternion kernel for the symplectic ensemble
In previous and this works, the Authors evaluated Dirac eigenvalue correlation functions for the symplectic ensemble with doubly degenerate masses. If masses are quadruply degenerate, the evaluation of the correlation functions is easier because the conventional integration method for the massless Laguerre ensemble works without any modification. In that case, the multiple integral we need to evaluate is
$`\mathrm{\Xi }_p(x_1,\mathrm{},x_p;\{m\})`$ (C1)
$`=`$ $`{\displaystyle \frac{1}{(Np)!}}{\displaystyle _0^{\mathrm{}}}\mathrm{}{\displaystyle _0^{\mathrm{}}}{\displaystyle \underset{j=p+1}{\overset{N}{}}}dz_j{\displaystyle \underset{j=1}{\overset{N}{}}}x_j^{2\nu +1}\mathrm{e}^{4x_j}`$ (C3)
$`\times {\displaystyle \underset{j=1}{\overset{N}{}}}{\displaystyle \underset{i=1}{\overset{\alpha }{}}}(x_j+m_i^2)^4{\displaystyle \underset{j>k}{\overset{N}{}}}x_jx_k^4,`$
where we set $`N_f=4\alpha `$ and
$`\{m\}=(\stackrel{4}{\stackrel{}{m_1,\mathrm{},m_1}},\mathrm{},\stackrel{4}{\stackrel{}{m_\alpha ,\mathrm{},m_\alpha }}).`$
The conventional ‘massless’ theory tells us that the correlation functions are written in terms of quaternion determinants:
$`\sigma (x_1,\mathrm{},x_p;\{m\})={\displaystyle \frac{\mathrm{\Xi }_p(z_1,\mathrm{},z_{\alpha +p})}{\mathrm{\Xi }_0(z_1,\mathrm{},z_\alpha )}}`$ (C4)
$`={\displaystyle \frac{\mathrm{Tdet}[f(z_j,z_k)]_{j,k=1,\mathrm{},\alpha +p}}{\mathrm{Tdet}[f(z_j,z_k)]_{j,k=1,\mathrm{},\alpha }}},`$ (C5)
where we have adopted the notation (31).
The quaternion function $`f(z,z^{})`$ is represented as
$$f(z,z^{})=\left[\begin{array}{cc}S(z,z^{})& I(z,z^{})\\ D(z,z^{})& S(z^{},z)\end{array}\right].$$
(C6)
We evaluate the asymptotic limit of the quaternion function $`f(z,z^{})`$ in each of the three cases
$`(a)z,z^{}>0,(b)z<0,z^{}>0,(c)z,z^{}<0,`$
as in Sect.2. in terms of the microscopic variables $`\zeta ,\zeta ^{},\mu ,\mu ^{}`$ defined by Eqs.(78), (84), (90), respectively.
$`(a)`$ $`z,z^{}>0`$
In the case $`z,z^{}>0`$, Nagao and Forrester derived the asymptotic limit (153), that is
$`{\displaystyle \frac{1}{8N}}S(z,z^{})`$ $``$ $`S_{++}(\zeta ,\zeta ^{}),`$ (C7)
$`{\displaystyle \frac{1}{(8N)^2}}D(z,z^{})`$ $``$ $`D_{++}(\zeta ,\zeta ^{}),`$ (C8)
$`I(z,z^{})`$ $``$ $`I_{++}(\zeta ,\zeta ^{}),`$ (C9)
where $`S_{++}`$, $`D_{++}`$, and $`I_{++}`$ are defined in Eq.(148).
$`(b)`$ $`z<0,z^{}>0`$
Using the asymptotic formula for the Bessel function (77), we can similarly treat negative argument cases to obtain
$`{\displaystyle \frac{1}{8N}}S(z,z^{})`$ $``$ $`S_+(\mu ,\zeta ),`$ (C10)
$`{\displaystyle \frac{1}{8N}}S(z^{},z)`$ $``$ $`\zeta {\displaystyle _0^1}𝑑tt^2{\displaystyle _0^1}𝑑u\left(J_{2\nu }(2tu\zeta )I_{2\nu +1}(2t\mu )J_{2\nu }(2t\zeta )uI_{2\nu +1}(2tu\mu )\right),`$ (C11)
$`{\displaystyle \frac{1}{(8N)^2}}D(z,z^{})`$ $``$ $`{\displaystyle _0^1}𝑑tt^3{\displaystyle _0^1}𝑑uu\left(I_{2\nu +1}(2tu\mu )J_{2\nu +1}(2t\zeta )I_{2\nu +1}(2t\mu )J_{2\nu +1}(2tu\zeta )\right),`$ (C12)
$`I(z,z^{})`$ $``$ $`I_+(\mu ,\zeta ),`$ (C13)
where $`S_+`$ and $`I_+`$ are defined in Eq.(148).
$`(c)`$ $`z,z^{}<0`$
$`{\displaystyle \frac{1}{8N}}S(z,z^{})`$ $``$ $`\mu {\displaystyle _0^1}𝑑tt^2{\displaystyle _0^1}𝑑u\left(I_{2\nu }(2tu\mu )I_{2\nu +1}(2t\mu ^{})I_{2\nu }(2t\mu )uI_{2\nu +1}(2tu\mu ^{})\right),`$ (C14)
$`{\displaystyle \frac{1}{(8N)^2}}D(z,z^{})`$ $``$ $`{\displaystyle _0^1}𝑑tt^3{\displaystyle _0^1}𝑑uu\left(I_{2\nu +1}(2tu\mu )I_{2\nu +1}(2t\mu ^{})I_{2\nu +1}(2t\mu )I_{2\nu +1}(2tu\mu ^{})\right),`$ (C15)
$`I(z,z^{})`$ $``$ $`I_{}(\mu ,\mu ^{}),`$ (C16)
where $`I_{++}`$ is defined in Eq.(148).
We can use Dyson’s equality (A16) to see that the above quaternion determinant expression is identical to the limit of quadruple mass degeneracy of the general formula (137) employing Pfaffians. In this case, yet another equivalent asymptotic formula was recently presented by Akemann and Kanzieper . Now we shall directly demonstrate the equivalence. We should firstly note that, under the change of the quaternion elements
$`S(z,z^{})`$ $``$ $`\stackrel{~}{S}(z,z^{})S(z,z^{}){\displaystyle \frac{1}{W(z^{})}}{\displaystyle \frac{dW(z^{})}{dz^{}}}I(z,z^{}),`$ (C17)
$`I(z,z^{})`$ $``$ $`\stackrel{~}{I}(z,z^{}){\displaystyle _z^z^{}}\stackrel{~}{S}(z,z^{\prime \prime })𝑑z^{\prime \prime },`$ (C18)
$`D(z,z^{})`$ $``$ $`\stackrel{~}{D}(z,z^{}){\displaystyle \frac{}{z}}\stackrel{~}{S}(z,z^{}),`$ (C19)
where
$$W(z)=|z|^{\nu +1/2}\mathrm{e}^{2z}.$$
(C20)
the quaternion determinant is unchanged. This transformation was introduced in Ref. .
For $`z,z^{}>0`$, we find an identity
$$\zeta \zeta ^{}_0^1𝑑tt_0^1𝑑uJ_{2\nu }(2tu\zeta )J_{2\nu }(2t\zeta ^{})=\zeta _{}^{}{}_{}{}^{2}_0^1𝑑tt_0^1\frac{du}{u}J_{2\nu +1}(2tu\zeta )J_{2\nu +1}(2t\zeta ^{})+\frac{\zeta ^{}}{2}J_{2\nu }(2\zeta ^{})_0^1\frac{du}{u}J_{2\nu +1}(2u\zeta ),$$
(C21)
by a partial integration. The asymptotic formulas (C8), together with the identity (C21), yield
$`\stackrel{~}{S}(z,z^{})`$ $``$ $`2N[2{\displaystyle _0^1}dttJ_{2\nu +1}(2t\zeta ^{})(2t\zeta {\displaystyle _0^1}duJ_{2\nu }(2tu\zeta )(2\nu +1){\displaystyle _0^1}{\displaystyle \frac{du}{u}}J_{2\nu +1}(2tu\zeta ))`$ (C24)
$`+2{\displaystyle \frac{\zeta }{\zeta ^{}}}{\displaystyle _0^1}𝑑ttJ_{2\nu }(2t\zeta )\left(2t\zeta ^{}{\displaystyle _0^1}𝑑uuJ_{2\nu +1}(2tu\zeta ^{})+(2\nu +1){\displaystyle _0^1}𝑑uJ_{2\nu }(2tu\zeta ^{})\right)`$
$`(2\nu +1){\displaystyle \frac{J_{2\nu }(2\zeta ^{})}{\zeta ^{}}}{\displaystyle _0^1}{\displaystyle \frac{du}{u}}J_{2\nu +1}(2u\zeta )].`$
Partial integrations give rise to the Bessel function equalities
$`(2\nu +1){\displaystyle _0^1}{\displaystyle \frac{du}{u}}J_{2\nu +1}(2tu\zeta )=J_{2\nu +1}(2t\zeta )+2t\zeta {\displaystyle _0^1}𝑑uJ_{2\nu }(2tu\zeta ),`$ (C25)
$`(2\nu +1){\displaystyle _0^1}𝑑uJ_{2\nu }(2tu\zeta ^{})=J_{2\nu }(2t\zeta ^{})+2t\zeta ^{}{\displaystyle _0^1}𝑑uuJ_{2\nu +1}(2tu\zeta ^{}).`$ (C26)
Substituting (C26) into (C24) and using the formula (B8), we obtain
$$\stackrel{~}{S}(z,z^{})2N\left[4_0^1𝑑ttJ_{2\nu +1}(2t\zeta )J_{2\nu +1}(2t\zeta ^{})+\frac{J_{2\nu }(2\zeta ^{})}{\zeta ^{}}J_{2\nu +1}(2\zeta )(2\nu +1)\frac{J_{2\nu }(2\zeta ^{})}{\zeta ^{}}_0^1\frac{du}{u}J_{2\nu +1}(2u\zeta )\right].$$
(C27)
By a partial integration, we can rewrite it as
$$\stackrel{~}{S}(z,z^{})2N\left[4_0^1𝑑ttJ_{2\nu +1}(2t\zeta )J_{2\nu +1}(2t\zeta ^{})\frac{J_{2\nu }(2\zeta ^{})}{\zeta ^{}}_0^{2\zeta }𝑑sJ_{2\nu +2}(s)\right].$$
(C28)
This is the asymptotic formula derived by Forrester, Nagao, and Honner in a study of parametric random matrices and then applied by Akemann and Kanzieper to the massive Dirac operator problem. Thus we established the equivalence in the case of quadruply degenerate masses. We can straightforwardly extend it to the formulas with $`z`$ and/or $`z^{}`$ negative.
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# The Evolution of the MLLA Parton Spectra.
## 1 Introduction
Quantum Chromodynamics (QCD) currently does not allow direct calculations of the hadronic final state observables. To make predictions of the final state it is necessary to model the transition from partons to hadrons. One approach that has been successful in describing the general features of the inclusive energy spectra, in both $`\mathrm{e}^+\mathrm{e}^{}`$ annihilation and deep inelastic scattering experiments , is the modified leading log approximation (MLLA) using local parton hadron duality (LPHD) as the method for relating the MLLA partonic predictions to the hadronic observables.
The production of hadrons in hard scattering processes is controlled by the underlying partonic behaviour. At small values of momentum fraction (of the outgoing parton from the original hard scatter) this parton behaviour, often referred to as a parton shower, is dominated by gluon bremsstrahlung. The branching processes $`qqg`$ and $`ggg`$ (double logarithmic processes) in addition to $`gq\overline{q}`$ (single logarithmic process) give rise to this parton shower. MLLA accounts for both these double and single logarithmic effects in the evolution equations . The perturbative properties of partonic distributions have been calculated in the framework of the MLLA. They are governed by two free parameters: a running strong coupling, defined by a QCD scale $`\mathrm{\Lambda },`$ and an energy cut-off, $`Q_0,`$ below which the parton evolution is truncated.
Using LPHD, the non-perturbative effects of particle distributions are reduced to a simple factor of normalisation that relates the hadronic distributions to the partonic ones. At large enough energies, away from the influence of $`Q_0,`$ this ‘hadronisation’ factor should be independent of the energy scale at which the spectra are being calculated.
Various approaches, detailed below, have been taken to calculate the single particle spectra and their moments within the MLLA framework. In this paper these theoretical approaches are compared as a function of energy scale, $`Q.`$ Substantial difference are found in the predictions of the higher moments. The theoretical results, in particular the higher order moments, are also compared with $`\mathrm{e}^+\mathrm{e}^{}`$ data over a range of centre of mass energies, $`Q=E_{CM}/2.`$ It is found that the general characteristics are well described but that the theory is not in accord with all features of the data if it is assummed that $`\mathrm{\Lambda },Q_0`$ and the LPHD normalisation are energy independent. The studies reported here avoid the need to extrapolate into regions not experimentally measured.
## 2 Single Particle Spectra
Given a high energy parton which fragments via secondary partons into a jet of hadrons, the MLLA evolution equation allows the secondary parton spectra for the logarithmic scaled energy, $`\xi ,`$ to be calculated . The variable $`\xi `$ is defined as $`\mathrm{ln}(Q/E)\mathrm{ln}(1/x),`$ where $`Q`$ is the energy of the original parton and $`E`$ is the energy of the secondary parton. The cut-off, $`Q_0,`$ limits the parton energy to $`Ek_TQ_0,`$ where $`k_T`$ is the transverse energy of the decay products in the jet evolution. In order to reconstruct the $`\xi `$ distributions an inverse Mellin transformation is performed
$$\frac{1}{N}\frac{dn_h}{d\xi }\overline{D}(\xi ,Y,\lambda )=_{ϵı\mathrm{}}^{ϵ+ı\mathrm{}}\frac{d\omega }{2\pi ı}x^\omega D(\omega ,Y,\lambda )$$
(1)
where the integral runs parallel to the imaginary axis on the right of all singularities in the complex $`\omega `$plane, $`Y=\mathrm{ln}(Q/Q_0)`$ and $`\lambda =\mathrm{ln}(Q_0/\mathrm{\Lambda }).`$
The Mellin-transformed distributions, $`D(\omega ,Y,\lambda ),`$ can be expressed in terms of confluent hypergeometric functions, $`\mathrm{\Phi },`$ as
$$\begin{array}{cc}D(\omega ,Y,\lambda )\hfill & =\frac{t_1A}{B(B+1)}\mathrm{\Phi }(A+B+1,B+2;t_1)\mathrm{\Phi }(AB,1B;t_2)\\ & +\left(\frac{t_2}{t_1}\right)^B\mathrm{\Phi }(A,B;t_1)\mathrm{\Phi }(A,B+1;t_2),\end{array}$$
(2)
where
$$\begin{array}{cc}t_1=\omega (Y+\lambda ),\hfill & \hfill t_2=\omega \lambda \end{array}$$
(3)
and $`A`$ and $`B`$ are defined as
$$\begin{array}{cc}A=4N_c/b\omega ,\hfill & \hfill B=a/b,\end{array}$$
(4)
where $`N_c`$ is the number of colours, $`a=11N_c/3+2n_f/3N_c^2`$ , $`n_f`$ is the number of flavours and $`b=11N_c/3+2n_f/3.`$ Throughout this paper the number of flavours is assumed to be 3. Equation (1) is then calculated using a numerical integration in the complex $`\omega `$plane.
It is convenient to investigate the MLLA spectra in terms of moments. The cumulant moments of the $`\xi `$ distribution can be written as:
$$K_m=\left(\frac{}{\omega }\right)^mD(\omega ,Y,\lambda )|_{\omega =0}.$$
(5)
The analytic form of these cumulant moments have been calculated in for the first four moments. This allows the normalised moment, $`\xi ^m,`$ to be calculated and hence the dispersion ($`\sigma `$), skewness ($`s`$) and the kurtosis ($`k`$) to be constructed, see for example
$`\xi `$ $`=`$ $`K_1`$ (6)
$`\sigma ^2`$ $`=`$ $`K_2=\xi ^2\xi ^2`$ (7)
$`s`$ $`=`$ $`K_3/\sigma ^3={\displaystyle \frac{\xi ^33\xi ^2\xi +2\xi ^3}{\sigma ^3}}`$ (8)
$`k`$ $`=`$ $`K_4/\sigma ^4={\displaystyle \frac{\xi ^44\xi ^3\xi 3\xi ^2^2+12\xi ^2\xi ^26\xi ^4}{\sigma ^4}}.`$ (9)
The so-called limiting spectrum is the case when $`Q_0=\mathrm{\Lambda }`$ i.e. $`\lambda =0`$. For this case, Fong and Webber have also calculated the behaviour of the moments of the $`\xi `$ spectra with energy scale, $`Q.`$ They point out that the spectra can be represented close to the maximum of the distribution by a distorted Gaussian of the form:
$$\overline{D}(\xi ,Y)\mathrm{exp}\left[\frac{1}{8}k\frac{1}{2}s\delta \frac{1}{4}(2+k)\delta ^2+\frac{1}{6}s\delta ^3+\frac{1}{24}k\delta ^4\right]$$
(10)
where $`\delta =(\xi l)/\sigma `$ and $`l`$ is the mean. This expression, when fitted to the spectrum, allows the moments to be determined when the full spectrum is either unmeasured or uncalculable.
The changes to the parton spectra for quark, $`q,`$ and gluon, $`g,`$ jets are Next-to-MLLA effects which result in additional terms in the integral shown in Eq. (1). These corrected spectra (at $`Q_0=\mathrm{\Lambda }`$) for the quark and gluon jet can be related to the limiting spectra as follows
$$\overline{D}_{q,g}=\left[1+\mathrm{\Delta }_{q,g}\left(\frac{}{l}+\frac{}{Y}\right)\right]\overline{D}^{lim}(l,Y)\overline{D}^{lim}(l+\mathrm{\Delta }_{q,g},Y+\mathrm{\Delta }_{q,g}),$$
(11)
with
$$\mathrm{\Delta }_g=\frac{1}{3}n_f\frac{N_c^21}{2N_c^3},\mathrm{\Delta }_q=\mathrm{\Delta }_0+\mathrm{\Delta }_g,\mathrm{\Delta }_0=\frac{a3N_c}{4N_c}.$$
The effect on the moments is that the $`\sigma ,`$ skewness and kurtosis for the quark/gluon distribution would be approximately the same as the limiting distribution at an effective energy $`Y+\mathrm{\Delta }_{q,g}.`$ The mean is the same as that of the limiting spectra at this effective energy but shifted in value by $`\mathrm{\Delta }_{q,g}.`$ These shifts are small with the limiting spectra being between the distribution for the quark and gluon jets. It should be noted that the predictions of Fong and Webber for the relative shifts between quark and gluon jets are slightly different from those discussed above.
In the MLLA approach, the partons are assumed massless so the scaled energy and momentum spectra are identical. Experimentally the scaled momentum distribution is usually measured and as the observed hadrons are massive the equivalence of the two spectra no longer holds. In the assumption is made that the cut-off $`Q_0`$ can be related to the masses of hadrons. This allows the logarithmic scaled momentum distribution, $`\xi _p,`$ to be written as
$$\frac{1}{N}\frac{dn_h}{d\xi _p}\frac{p_h}{E_h}\overline{D}(\xi ,Y),$$
(12)
where
$$\xi =\mathrm{log}\frac{Q}{\sqrt{Q^2e^{2\xi _p}+Q_0^2}},$$
the energy of a hadron with a momentum $`p_h`$ is $`E_h=\sqrt{p_h^2+Q_0^2}.`$ Limiting momentum spectra based on massless partons and massive partons will be referred to as MLLA-0 and MLLA-M spectra, respectively.
## 3 Behaviour of Theoretical Spectra
In this section the general characteristics of the spectra are discussed in order to illustrate their behaviour as a function of the three variables $`Q`$, $`\lambda `$, and $`Q_0`$.
Figure 1(a) shows the limiting energy spectra for three energies $`Q=183./2`$ (LEP II), $`91.2/2`$ (LEP) and $`14.8/2\mathrm{GeV}`$ (TASSO) with a $`\mathrm{\Lambda }`$ value of $`250\mathrm{MeV}.`$ As the energy scale increases, the parton multiplicity (the area under the curve) grows and the peak position shifts to the right i.e. the parton scaled energy spectra is softer. The cut-off on the right hand side of the plot corresponds to the truncation in the parton energy spectra at $`Q_0=\mathrm{\Lambda }.`$
Figure 1(b) shows the scaled energy spectra at fixed energy scale $`Q=91.2\mathrm{GeV}`$ and $`\mathrm{\Lambda }=250\mathrm{MeV}`$ but at two different values for the cut-off of the parton evolution, $`Q_0=\mathrm{\Lambda }\mathrm{and}2\mathrm{\Lambda }.`$ (The curves are cropped at $`\xi =1.0`$ as below this value the numerical integration of the truncated spectra is unstable.) Truncating the cascade at higher values of $`Q_0`$ leads to a lower parton multiplicity with a harder energy spectra.
Figure 1(c), again, shows the limiting energy spectra but at a fixed energy scale, $`Q=91.2/2\mathrm{GeV},`$ with values of $`\mathrm{\Lambda }=50,250\mathrm{and}400\mathrm{MeV}.`$ As $`\mathrm{\Lambda }`$ decreases, more partons are produced and their energy spectra is softer. In the case of the limiting spectra, this behaviour is dominated by the fact that an increase in $`\mathrm{\Lambda }`$ is accompanied by an increase in the cut-off $`Q_0.`$
Figure 1(d) shows the limiting scaled momentum spectra for the MLLA-0 $`(E_h=p_h)`$ and MLLA-M case $`(E_hp_h)`$ at an $`E_{CM}=91.2\mathrm{GeV}.`$ The MLLA-0 spectra displays the usual truncation at large $`\xi `$ associated with the cut-off $`Q_0,`$ but the MLLA-M case does not. This is due to the fact the calculation is now regulated by the cut-off entering as a mass term in the expression of $`\xi _p`$ (see Eq. 12.) As the momentum $`Q,`$ i.e. $`\xi _p0.0,`$ the introduction of a mass term has no effect on the MLLA calculation. As $`\xi _p`$ increases (the momentum becomes smaller) the mass term begins to play an increasingly important rôle. From these arguments it can be seen that as $`Q`$ decreases mass term has a more significant influence over a larger (fractional) range of the $`\xi _p`$ spectra.
Figure 2 shows the evolution of the mean, $`\sigma ,`$ skewness and kurtosis of the limiting energy spectra with $`\mathrm{\Lambda }=250\mathrm{MeV}`$ as a function of $`Q.`$ Three different approaches to calculating the moments have been investigated: (i) the analytic expression of Fong and Webber for gluon jets (dashed line) and quark jets(dash-dotted line); (ii) the analytic calculation of Dokshitzer et al. (dotted line); (iii) fitting a distorted Gaussian (Eq. 10) over $`\pm 1\sigma `$ around the mean value (full line). The range of the fit is motivated by the phenomenological scope of validity. All approaches exhibit the same trends, in that all the cumulants increase as $`Q`$ increase. The values of skewness and kurtosis are negative, tending towards zero as $`Q`$ increases, i.e. the spectra are becoming more like a pure Gaussian. The mean and the $`\sigma `$ exhibit a very similar dependence on $`Q`$ for all three approaches, though with different offsets. The skewness and kurtosis exhibit a different dependence on $`Q,`$ but as $`Q`$ approaches the asymptotic limit the predictions are converging.
The differences between quark and gluon jets for the predictions of Fong and Webber, in Fig. 2, are small. A feature of the data is the marked difference between the analytic calculations of Dokshitzer et al. and the fit to the limiting spectra. Calculating the moments over the full range of the spectra gives the same result as the analytic calculations, not too surprisingly since both calculations use the same assumptions. Investigating a limited range of the spectra, as is done during the fit, produces marked differences, highlighting the sensitivities of these higher-order moments to the behaviour of the tails as well as the incorrect form of the distorted Gaussian away from the peak position. The difference between the analytic calculation of Fong and Webber and that of Dokshitzer et al. could well be due to the differences in the predictions at large $`x_p,`$ where the theoretical assumption made in both calculations are no longer valid.
Figures 1 and 2 imply that $`\mathrm{\Lambda }`$ influences the position of the maximum and the width of the distribution of $`\xi `$. The effect of the relation between $`Q_0`$ and $`\mathrm{\Lambda }`$ is more complex; setting $`Q_0=2\mathrm{\Lambda }`$ lowers the position of the maximum for a given $`\mathrm{\Lambda }`$. The maximum can be returned to the original position by reducing $`\mathrm{\Lambda }`$ but with a consequential increase of the width of the distribution relative to that found for the limiting spectrum. The effect of introducing mass terms into the momentum spectra is to increase the tail of the spectra at high $`\xi ,`$ thus increasing skewness. As $`Q`$ increases, the influence of the mass term becomes less important.
## 4 Comparison with data
The differential $`\xi _p`$ distributions for $`\mathrm{e}^+\mathrm{e}^{}`$ data and the models are shown in Figure 3. The limiting spectra were calculated with $`\mathrm{\Lambda }=250\mathrm{MeV}`$; this value was chosen to give the correct peak position of the spectra for an $`E_{CM}(=2Q)`$ at the mass of the $`\mathrm{Z}^0,\mathrm{M}_\mathrm{Z}.`$ For each energy, the theoretical spectra are normalised to the same maximum height as the data.
The theoretical spectra follow the general trend of the data for both for both MLLA-0 and MLLA-M spectra; in detail there are discrepancies. At low energies, for the MLLA-M case, the skewness is larger (i.e. more positive ) than that observed for the data; at high energies, the limiting spectra have a smaller width than the data. Also shown are the MLLA energy spectra for $`Q_0=2\mathrm{\Lambda }`$. Here $`\mathrm{\Lambda }=50\mathrm{MeV}`$ is required to get the correct position of the maximum for $`E_{CM}=\mathrm{M}_\mathrm{Z}.`$ The position of the maximum has approximately the correct energy dependence; however, the width of the theoretical spectrum is consistently greater than that of the data.
To study the energy evolution of the MLLA theoretical spectra and the $`\mathrm{e}^+\mathrm{e}^{}`$ data, fits around the peak position, using the distorted Gaussian (Eq. 10) have been made to determine the moments of the $`\xi _p`$ distributions. In view of the statistical limitations, the range of the fit for the data was about three units around the peak (see Figure 3). The distorted Gaussian describes the data well over the fitted region with a $`\chi ^2/dof`$ of 1.0 or better. To permit direct comparison with the data, a similar fit range was used for the theoretical spectra. The distorted Gaussian gives a good description of the MLLA-M spectra. The description is less good for the MLLA-0 spectra but for all energies the Gaussian represents the model to better than $`1\%`$.
The moments from the fits are shown in Figure 4 as solid lines for fits to MLLA-0 spectra and as dashed lines for fits to the MLLA-M spectra. The data points are from our fit to $`\mathrm{e}^+\mathrm{e}^{}`$ data .
It may be concluded from Figure 4 that the MLLA-0 model gives a good description of the $`\mathrm{e}^+\mathrm{e}^{}`$ data for the mean and skewness. There are however discrepancies in $`\sigma `$ and kurtosis as is also evident in Figure 3; $`\sigma `$ is smaller, and the model is more platykurtic than the data. The MLLA-M model gives a poorer description of all variables at low $`E_{CM}`$ but approaches both the MLLA-0 predictions and the data at high $`E_{CM}`$.
This is contrary to the conclusions reported by Lupia and Ochs . In that paper the moments of the experimental distributions corrected for the mass effects were compared to the predictions for the limiting spectra and found to be in good agreement. It should be noted that in an extrapolation was made into regions where no experimental measurement exists to allow the moments to be calculated. In addition, the agreement of the MLLA limiting spectra with the mass corrected experimental $`\xi `$ distribution is poor. The spectra at low $`\xi `$ are well described but the region around the peak and at higher $`\xi `$ is generally not well reproduced. The studies reported here avoid the need to extrapolate into regions not experimentally measured.
LPHD postulates that a constant factor relates the MLLA predicted spectrum to the experimental spectrum, independent of the $`E_{CM}.`$ In the spectra discussed above, the maxima of the MLLA predictions have been scaled to the data. This scale factor, for the MLLA-0 model, is shown in Figure 5 as a function of $`E_{CM}.`$ It falls from $`1.4`$ at the lowest $`E_{CM}`$ to $`1.2`$ at the highest. A similar result is obtained for an area normalisation over approximately three units of $`\xi _p`$ for the MLLA-M case. Thus our observations disagree with the LPHD postulate of a constant scale factor. This suggests that at the currently accessible experimental energies the theoretical spectra are still influenced by the $`Q_0`$ cut-off.
Similar comparisons of $`\mathrm{e}^+\mathrm{e}^{}`$ data with MLLA predictions have been made by the DEPHI and OPAL collaborations. CDF has examined jets produced in $`p\overline{p}`$ interactions . In contrast to the analysis presented here, no explicit analysis of the higher moments of the $`\xi `$ distribution has been made. However, where the analyses can be compared, the results are in agreement with those presented here. In particular, $`\mathrm{\Lambda }`$ is close to 250MeV, the shape of the MLLA predicted spectrum is close to that of the data but disagrees in detail, and the value of the LPHD constant changes with energy.
## 5 Conclusions
The MLLA predictions of the logarithmic scaled energy spectra, and in particular their moments, have been investigated as function of energy scale. Various theoretical approaches have been compared and differences discussed. We note that care should be taken in the comparisons to ensure that a consistent approach is maintained; in particular attention should be paid to the range of application in $`\xi .`$ There is a large discrepancy in the skewness and kurtosis at low $`Q`$ but the various predictions converge as the asymptotic limit is approached.
The theoretical results are compared to measurements taken in $`\mathrm{e}^+\mathrm{e}^{}`$ annihilation experiments. The calculations are generally in good agreement with the data. It is observed that the limiting spectra is preferred over the predictions of the truncated cascade ($`Q_0\mathrm{\Lambda }`$.) The introduction of the mass term has a large effect at all but the highest $`E_{CM}.`$ The data is consistently broader than the limiting spectra over the energy range studied here.
The normalisation factor of LPHD between the data and the MLLA theoretical predictions is not constant. Contrary to expectation it displays a dependence on $`E_{CM},`$ decreasing as $`E_{CM}`$ increases, thus suggesting a residual influence of the $`Q_0`$ cut-off.
## Acknowledgements
The authors would like to thank Yuri Dokshitzer for useful discussion and clarification of issues.
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# Measurements of Nondegenerate Discrete Observables
## I Introduction
Every measurement on a quantum system causes a quantum state reduction, a state transformation $`\rho \rho _x`$ from the system state $`\rho `$ just before the measurement to the system state $`\rho _x`$ just after the measurement conditional upon the outcome $`x`$ of measurement. In order to determine the quantum state reduction caused by a measurement of a given observable, von Neumann posed the repeatability hypothesis \[1, p. 335\]: If an observable is measured twice in succession in a system, then we get the same value each time. Then, a measurement of a discrete observable $`A`$ satisfies the repeatability hypothesis if and only if $`\rho _x`$ is a mixture of eigenstates of $`A`$ corresponding to the eigenvalue $`x`$. Thus, under this hypothesis the measurement of a nondegenerate discrete observable $`A`$ causes the unique quantum state reduction such that $`\rho _x`$ is the unique eigenstate corresponding to the eigenvalue $`x`$. However, von Neumann admitted that there are many quantum state reductions caused by measuring the degenerate observable $`A`$ even when the repeatability hypothesis holds \[1, p. 348\]. Thus, for degenerate observables further hypothesis were demanded. In order to characterize the least disturbing measurement, Lüders posed the projection postulate: The measurement of a discrete observable $`A`$ in the state $`\rho `$ leaves the system in the state $`\rho _x=E_x\rho E_x/\text{Tr}[E_x\rho ]`$, where $`E_x`$ is the eigenprojection corresponding to the eigenvalue $`x`$. Obviously, the projection postulate implies the repeatability hypothesis and determines the output state uniquely, even for the degenerate discrete observable $`A`$.
Despite the above attempts, Davies and Lewis conjectured that no measurements of continuous observables satisfy the repeatability hypotheses and proposed abandoning the repeatability hypothesis. Actually, their conjecture was proved later in ; the essential part of the proof given in shows that even in the measurement of a continuous observable the output state conditional upon the outcome can be still described by a density operator. Moreover, it can be readily seen that there are many ways of measuring the same observable without satisfying the repeatability hypothesis such as photon counting , which arises every optical experiment, and contractive state measurement , which beats the standard quantum limit for monitoring the free-mass position claimed in .
Once we abandon the repeatability hypothesis or the projection postulate, the problem of determining all the possible quantum state reductions caused by measurements of a given observable has a primary importance in quantum mechanics. Especially, the problem receives increasing interests recently not only from the foundational point of view but also from the technological point of view, since the measurement is used for preparing the state of the system in such processes as purification procedures in the field of quantum information .
The purpose of this paper is to give the complete solution to the above problem for the measurements of nondegenerate discrete observables. It will be shown that a surprisingly general condition for the measurement statistics suffices to determine all the possible quantum state reductions realized by indirect measurement models. It is shown that for measurements of nondegenerate discrete observables the output state is independent of the input state in any measurement and that the family of output states can be arbitrarily chosen by the choice of the apparatus. Moreover, all of them are shown to have indirect measurement models.
In order to obtain a mathematical description of quantum state reductions for the most general class of measurements we consider the two requirements: one is necessary and the other is sufficient.
The necessary one is the mixing law of the joint probability that requires that the joint probability distribution of the outcomes of the successive measurement depends affinely on the input state. We require this condition as a necessary condition for every apparatus to satisfy. It will be shown that this is equivalent to the requirement that every apparatus has a normalized positive superoperator valued measure that satisfies the Davies-Lewis description of conditional state transformations . The notion of normalized positive superoperator valued measures was first introduced by Davies and Lewis to obtain a general description of conditional state transformation by unifying the notions of operations , effects , and probability operator valued measures . Thus, the problem of determining possible quantum state reductions is reduced to the problem as to which normalized positive superoperator valued measure corresponds to an apparatus.
The sufficient condition is the unitary realizability condition that requires the existence of an indirect measurement model comprising of the probe preparation, the measuring interaction with unitary time evolution, and the probe detection. We require this condition as a sufficient condition so that if a normalized positive superoperator valued measure has an indirect measurement model then the corresponding apparatus exists. It was proved in that this condition is equivalent to the condition that the normalized positive superoperator valued measure is completely positive.
According to the above approach, the class of possible quantum state reductions is included in the class of conditional state transformations satisfying the mixing law, i.e., the normalized positive superoperator valued measures, and includes the one satisfying the realizability condition, i.e., the normalized completely positive superoperator valued measures. These two classes are generally different.
Nevertheless, for the case where $`A`$ is nondegenerate, this paper shows, the above two conditions are actually equivalent. Thus, both of them are necessary and sufficient and we reach a clear-cut conclusion. According to the analysis developed in this paper, for any apparatus $`𝐀`$ measuring a nondegenerate discrete observable $`A=_na_n|\varphi _n\varphi _n|`$ there is a sequence $`\{\mathit{\varrho }_n\}`$ of density operators independent of the input state $`\rho `$ such that the measurement leaves the system in the state $`\mathit{\varrho }_n`$ with the probability $`\varphi _n|\rho |\varphi _n`$, and conversely for any sequence $`\{\mathit{\varrho }_n\}`$ of density operators such an apparatus exists.
## II Measuring Apparatuses
Let us consider the conventional quantum-mechanical description of the measurement of an observable represented by a self-adjoint operator $`A`$ with purely discrete spectrum on a separable Hilbert space $``$. For any real number $`x`$ we shall denote by $`E^A(x)`$ the projection of $``$ onto the subspace $`\{\psi |A\psi =x\psi \}`$. If $`A`$ has eigenvalues $`a_1,a_2,\mathrm{}`$ then $`E^A(a_n)`$ is the spectral projection corresponding to $`a_n`$ and $`E^A(x)=0`$ if $`x`$ is not an eigenvalue of $`A`$. If the state of the system at the instant before the measurement is given by the density operator $`\rho `$ on $``$, then the measurement yields the outcome $`a_n`$ with the probability $`\text{Tr}[E^A(a_n)\rho ]`$. If this measurement satisfies the projection postulate , then the state at the instant after the measurement is
$$\rho _n=\frac{E^A(a_n)\rho E^A(a_n)}{\text{Tr}[E^A(a_n)\rho ]}$$
(1)
provided that the measurement leads to the outcome $`a_n`$.
As it can be seen from the above description, every measuring apparatus $`𝐀`$ has the output variable $`𝐱`$ that takes the outcome in each measurement carried out by $`𝐀`$. Thus, the output variable is a random variable the probability distribution of which depends only on the input state, the state of the system at the instant just before the measurement. Throughout this paper, we assume that the output variable takes the values in a countable subset of the real line $`𝐑`$. The probability distribution $`\mathrm{Pr}\{𝐱=x\rho \}`$ of $`𝐱`$ in the input state $`\rho `$ is called the output distribution of $`𝐀`$. The change from the unknown input state to the output distribution is called the objective state reduction. Depending on the input state $`\rho `$ and the outcome $`𝐱=x`$, the state $`\rho _{\{𝐱=x\}}`$ just after the measurement is determined uniquely. The state $`\rho _{\{𝐱=x\}}`$ is called the output state relative to the input state $`\rho `$ and the outcome $`𝐱=x`$. If the output probability of $`𝐱=x`$ is 0, the output state $`\rho _{\{𝐱=x\}}`$ is taken to be indefinite. The change from the unknown input state to the output state is called the quantum state reduction. The above two mathematical objects, the objective state reduction and the quantum state reduction, are called the statistical property of $`𝐀`$. Two apparatuses are called statistically equivalent if they have the same statistical property. In what follows, every apparatus is supposed to have its own distinctive output variable and we denote by $`𝐀(𝐱)`$ the apparatus having the output variable $`𝐱`$.
In the above measurement of $`A`$ satisfying the projection postulate, let us denote the measuring apparatus by $`𝐀(𝐚)`$ where $`𝐚`$ stands for the output variable. The statistical property of $`𝐀(𝐚)`$ is represented as follows.
output distribution: $`\mathrm{Pr}\{𝐚=x\rho \}=\text{Tr}[E^A(x)\rho ]`$ (2)
output state: $`\rho _{\{𝐚=a_n\}}={\displaystyle \frac{E^A(a_n)\rho E^A(a_n)}{\text{Tr}[E^A(a_n)\rho ]}}`$ (3)
In the above, $`a_n`$ is an eigenvalue such that $`\text{Tr}[E^A(a_n)\rho ]>0`$.
Now, the following problem arises: Does every measuring apparatus for the observable $`A`$ necessarily have the above statistical property? It is postulated by the Born statistical formula that the output distribution of the measurement of the observable $`A`$ satisfies (2). Hence, every measuring apparatus for the observable $`A`$ satisfies (2) by definition. The following argument will show, however, that the existence of an apparatus satisfying the projection postulate implies the existence of another apparatus which does not satisfy the projection postulate. Therefore, we cannot postulate that every measurement satisfies the projection postulate.
Suppose that the observable $`Y`$ has degenerate eigenvalues and can be represented by
$$Y=\underset{n,m}{}y_n|n,mn,m|$$
(4)
for some orthonormal basis $`\{|n,m\}`$. Consider the following process of measuring $`Y`$: (i) One measures the nondegenerate discrete observable
$$X=\underset{n,m}{}x_{n,m}|n,mn,m|$$
(5)
where $`x_{n,m}`$ are all different. (ii) If the outcome $`𝐱`$ of the $`X`$ measurement leads to the value $`x_{n,m}`$ then the outcome $`𝐲`$ of the $`Y`$ measurement is determined as $`y_n`$. Then, even if the $`X`$ measurement satisfies the projection postulate, the $`Y`$ measurement does not satisfy it. In fact, with the probability $`n,m|\rho |n,m`$ the $`X`$ measurement leads to the outcome $`x_{n,m}`$ and leaves the system in the state $`|n,mn,m|`$ by the projection postulate. It follows that if the outcome is $`y_n`$ then the state at the instant after the $`Y`$ measurement is given by
$`\rho _{\{𝐲=y_n\}}={\displaystyle \frac{_mn,m|\rho |n,m|n,mn,m|}{_mn,m|\rho |n,m}}.`$ (6)
The above state depends on the choice of the orthonormal basis $`\{|n,m\}`$. If $`Y`$ is degenerate, there are infinitely many essentially different choices of $`\{|n,m\}`$ and each choice gives a process of $`Y`$ measurement which does not satisfy the projection postulate.
Generalizing the above, if two observables $`X,Y`$ has the relation $`Y=f(X)`$, then for any apparatus $`𝐀(𝐱)`$ measuring $`X`$ we have the apparatus $`𝐀(f(𝐱))`$ measuring $`Y`$ that outputs the outcome $`f(𝐱)=f(x)`$ whenever $`𝐀(𝐱)`$ outputs the outcome $`𝐱=x`$. In this case, even if $`𝐀(𝐱)`$ satisfies the projection postulate, $`𝐀(f(𝐱))`$ does not necessarily satisfies the projection postulate. Therefore, the output distribution of $`Y`$ measurement is unique but the quantum states reduction depends on the way of measuring the same observable $`Y`$. More general construction of measuring apparatuses that do not satisfy the projection postulate will be discussed in Section VIII.
Can one determine all the possible quantum state reductions arising in measuring $`A`$ that are allowed by the basic principles of quantum mechanics? This problem will be considered in the following sections.
## III Successive measurements
In order to clarify the operational meaning of the quantum state reduction, we shall generalize von Neumann’s idea on repeated measurements of the same observable \[1, pp. 211–223\] to arbitrary pair of measuring apparatuses and consider the joint probability distribution of the outcomes of the two measurements carried out in succession.
We suppose that the $`A`$ measurement described in the preceding section is immediately followed by a measurement of a discrete observable $`B`$ with eigenvalues $`b_m`$. Then, the conditional probability of obtaining the outcome $`b_m`$ at the $`B`$ measurement is $`\text{Tr}[E^B(b_m)\rho _n]`$ conditional upon having obtained $`a_n`$ at the $`A`$-measurement. From (1), the joint probability of obtaining $`a_n`$ at the $`A`$ measurement and $`b_m`$ at the $`B`$ measurement is therefore
$`p_{n,m}`$ $`=`$ $`\text{Tr}[E^B(b_m)\rho _n]\text{Tr}[E^A(a_n)\rho ]`$ (7)
$`=`$ $`\text{Tr}[E^B(b_m)E^A(a_n)\rho E^A(a_n)].`$ (8)
Generally speaking, if a measurement by the apparatus $`𝐀(𝐱)`$ in the input state $`\rho `$ is immediately followed by a measurement by the apparatus $`𝐀(𝐲)`$, the joint probability distribution $`\mathrm{Pr}\{𝐱=x,𝐲=y\rho \}`$ of the output variables $`𝐱`$ and $`𝐲`$ depends only on the input state $`\rho `$ of the first measurement and is given by
$$\mathrm{Pr}\{𝐱=x,𝐲=y\rho \}=\mathrm{Pr}\{𝐲=y\rho _{\{𝐱=x\}}\}\mathrm{Pr}\{𝐱=x\rho \}.$$
(9)
This joint probability distribution has the following significant property.
Mixing law of the joint probability: For any measuring apparatuses $`𝐀(𝐱)`$ and $`𝐀(𝐲)`$, if the input state $`\rho `$ is the mixture of $`\rho _1`$ and $`\rho _2`$ such that $`\rho =\alpha \rho _1+(1\alpha )\rho _2`$ with $`0<\alpha <1`$ then the joint probability distribution of the outcomes of the successive measurement satisfies
$`\mathrm{Pr}\{𝐱=x,𝐲=y\rho \}`$ $`=`$ $`\alpha \mathrm{Pr}\{𝐱=x,𝐲=y\rho _1\}`$ (11)
$`+(1\alpha )\mathrm{Pr}\{𝐱=x,𝐲=y\rho _2\}.`$
This is justified as follows. If the system is in the state $`\rho _1`$ with the probability $`\alpha `$ and in the state $`\rho _2`$ with the probability $`1\alpha `$ then the joint probability is their mixture in the right hand side. On the other hand, in this case the state of the system is described by the density operator $`\rho `$ and hence the above equality holds.
In the previous example, if the observable $`B`$ is measured by an apparatus $`𝐀(𝐛)`$ then from (8) we have
$$\mathrm{Pr}\{𝐚=a_n,𝐛=b_m\rho \}=\text{Tr}[E^B(b_m)E^A(a_n)\rho E^A(a_n)].$$
(12)
Obviously, this joint probability satisfies the above mixing law.
In what follows, we require the mixing law of the joint probability. For an arbitrary apparatus $`𝐀(𝐱)`$ with the output distribution $`\mathrm{Pr}\{𝐱=x\rho \}`$ and the output state $`\rho _{\{𝐱=x\}}`$, we define the output operator $`𝐗(x,\rho )`$ by
$$𝐗(x,\rho )=\mathrm{Pr}\{𝐱=x\rho \}\rho _{\{𝐱=x\}}.$$
(13)
Then, $`𝐗(x,\rho )`$ is a trace class operator determined by the statistical property of the apparatus $`𝐀(𝐱)`$, the input state $`\rho `$, and the outcome $`𝐱=x`$.
For the measuring apparatus $`𝐀(𝐚)`$, the output operator $`𝐗_𝐚(x,\rho )`$ is given by
$$𝐗_𝐚(x,\rho )=E^A(x)\rho E^A(x).$$
The above expression extends the definition of $`𝐗_𝐚(x,\rho )`$ to arbitrary trace class operators $`\rho `$. Then, it is easy to see that $`𝐗_𝐚(x,\rho )`$ has the following properties: (i) $`𝐗_𝐚(x,\rho )`$ is a positive operator if $`\rho `$ is positive, (ii) the correspondence $`\rho 𝐗_𝐚(x,\rho )`$ is linear, (iii) for any $`\rho `$ we have
$$\text{Tr}[\underset{x}{}𝐗_𝐚(x,\rho )]=\text{Tr}[\rho ].$$
In the following, we shall show that the output operator $`𝐗(x,\rho )`$ of every apparatus $`𝐀(𝐱)`$ has the above properties.
Return to the joint probability distribution $`\mathrm{Pr}\{𝐱=x,𝐲=y\rho \}`$. If one measures the observable $`B`$ by the apparatus $`𝐀(𝐛)`$ instead of $`𝐀(𝐲)`$, from (9) and (13) we have
$`\mathrm{Pr}\{𝐱=x,𝐛=b_m\rho \}`$ $`=`$ $`\text{Tr}[E^B(b_m)\rho _{\{𝐱=x\}}]\mathrm{Pr}\{𝐱=x\rho \}`$ (14)
$`=`$ $`\text{Tr}[E^B(b_m)𝐗(x,\rho )].`$ (15)
Suppose that $`\rho `$ is the mixture $`\rho =\alpha \rho _1+(1\alpha )\rho _2`$. From (11) we have
$`\text{Tr}[E^B(b_m)𝐗(x,\rho )]`$
$`=`$ $`\alpha \text{Tr}[E^B(b_m)𝐗(x,\rho _1)]+(1\alpha )\text{Tr}[E^B(b_m)𝐗(x,\rho _2)]`$
$`=`$ $`\text{Tr}[E^B(b_m)[\alpha 𝐗(x,\rho _1)+(1\alpha )𝐗(x,\rho _2)]].`$
Since $`B`$ is arbitrary, we have
$$𝐗(x,\rho )=\alpha 𝐗(x,\rho _1)+(1\alpha )𝐗(x,\rho _2).$$
(16)
In what follows, for any $`x`$ let $`𝐗(x)`$ be the mapping that maps a density operator $`\rho `$ to the trace class operator $`𝐗(x,\rho )`$. Since every trace class operator $`\sigma `$ can be represented as the linear combination
$$\sigma =\lambda _1\sigma _1\lambda _2\sigma _2+i\lambda _3\sigma _3i\lambda _4\sigma _4$$
(17)
with four density operators $`\sigma _1,\mathrm{},\sigma _4`$ and four positive numbers $`\lambda _1,\mathrm{},\lambda _4`$, we can extend the mapping $`𝐗(x)`$ to a linear transformation on the space $`\tau c()`$ of trace class operators on $``$ by
$`𝐗(x)\sigma `$ $`=`$ $`\lambda _1𝐗(x)\sigma _1\lambda _2𝐗(x)\sigma _2`$ (19)
$`+i\lambda _3𝐗(x)\sigma _3i\lambda _4𝐗(x)\sigma _4.`$
Since the decomposition $`(\text{17})`$ is not unique, in order for the extension (19) to be well-defined we need to show that the left hand side of (19) is uniquely determined independent of the decomposition of $`\sigma `$. This can be proved from (16) and the proof will be shown in Appendix A.
We have, therefore, shown that for every apparatus $`𝐀(𝐱)`$ there exists a family $`\{𝐗(x)|x𝐑\}`$ of linear transformations on $`\tau c()`$ such that for every density operator $`\rho `$, we have
$$𝐗(x)\rho =\mathrm{Pr}\{𝐱=x\rho \}\rho _{\{𝐱=x\}}.$$
(20)
The linear transformation $`𝐗(x)`$ defined above is called the operation of the apparatus $`𝐀(𝐱)`$ for the outcome $`𝐱=x`$. The family $`\{𝐗(x)|x𝐑\}`$ is called the operational distribution of the apparatus $`𝐀(𝐱)`$. It is obvious from (20) that by taking advantage of the operational distribution, the output distribution is represented by
$$\mathrm{Pr}\{𝐱=x\rho \}=\text{Tr}[𝐗(x)\rho ]$$
(21)
and the output state by
$$\rho _{\{𝐱=x\}}=\frac{𝐗(x)\rho }{\text{Tr}[𝐗(x)\rho ]},$$
(22)
where the outcome $`𝐱=x`$ is supposed to have positive probability.
## IV Operational distributions
In order to explore mathematical properties of the operational distribution $`\{𝐗(x)|x𝐑\}`$ of the apparatus $`𝐀(𝐱)`$, we shall provide relevant mathematical terminology. A linear transformation $`𝐋`$ on the space $`\tau c()`$ of trace class operators on $``$ is said to be bounded if there is a constant $`K>0`$ such that
$$𝐋\rho _{tr}K\rho _{tr}$$
for all $`\rho \tau c()`$, where $`_{tr}`$ stands for the trace norm. Then, the norm of $`𝐋`$ is defined by
$$𝐋_{tr}=\underset{\rho _{tr}1}{sup}𝐋\rho _{tr}.$$
(23)
A linear transformation $`𝐌`$ on the space $`()`$ of bounded operators on $``$ is said to be bounded if there is a constant $`K>0`$ such that
$$𝐌AKA$$
for all $`A()`$, where $``$ stands for the operator norm. Then, the norm of $`𝐌`$ is defined by
$$𝐌=\underset{A1}{sup}𝐌A.$$
(24)
A bounded linear transformation on $`\tau c()`$ is called a superoperator. For any superoperator $`𝐋`$ on $`\tau c()`$, its dual superoperator $`𝐋^{}`$ is the bounded linear transformation on $`()`$ defined by
$$\text{Tr}[A(𝐋\rho )]=\text{Tr}[(𝐋^{}A)\rho ]$$
(25)
for all $`A()`$ and $`\rho \tau c()`$. In this case, we have $`𝐋_{tr}=𝐋^{}`$. A superoperator or a dual superoperator is said to be positive iff it maps positive operators to positive operators. Then, a superoperator $`𝐋`$ is positive if and only if so is its dual. A super operator or a dual superoperator is said to be contractive iff it has the norm less than or equal to one. Then, a superoperator $`𝐋`$ is a contractive if and only if so is its dual. We have the following characterizations of positive contractive superoperators \[23, p. 216\], \[16, p. 18\].
###### Theorem 1
For a positive superoperator $`𝐋`$ the following conditions are all equivalent:
(i) $`𝐋`$ is a contractive superoperator.
(ii) $`𝐋^{}`$ is a contractive dual superoperator.
(iii) $`0\text{Tr}[𝐋\rho ]1`$ for all density operators $`\rho `$.
(vi) $`0𝐋^{}(I)I.`$
Moreover, a superoperator $`𝐋`$ is trace preserving, i.e.,
$$\text{Tr}[𝐋(\rho )]=\text{Tr}[\rho ]$$
for all $`\rho \tau c()`$ if and only if $`𝐋^{}`$ is unital, i.e.,
$$𝐋^{}(I)=I.$$
Let us return to the operational distribution $`\{𝐗(x)|x𝐑\}`$ of the apparatus $`𝐀(𝐱)`$. Let $`𝐀(𝐛)`$ be an apparatus measuring a discrete observable $`B`$ with eigenvalues $`b_m`$ and let $`\rho `$ be an arbitary density operator. By the property of joint probability, we have
$$0\mathrm{Pr}\{𝐱=x,𝐛=b_m\rho \}1.$$
From (15) we have
$$0\text{Tr}[E^B(b_m)𝐗(x)\rho ]1.$$
Since $`B`$ and $`\rho `$ are arbitrary, the operation $`𝐗(x)`$ is a positive superoperator. Taking $`B=I`$ and $`b_m=1`$, we have
$$0\text{Tr}[𝐗(x)\rho ]1.$$
It follows from Thorem 1 that the operation $`𝐗(x)`$ is a positive contractive superoperator. By the unicity of total probability, we have
$$\underset{x𝐑}{}\mathrm{Pr}\{𝐱=x\rho \}=1.$$
Hence, we have
$$\text{Tr}[\underset{x𝐑}{}𝐗(x)\rho ]=1.$$
for all density operator $`\rho `$. Let $`𝐗(x)^{}`$ be the dual of the operation $`𝐗(x)`$. It follows that
$$\underset{x𝐑}{}𝐗(x)^{}I=I.$$
(26)
and that
$$\text{Tr}[\underset{x𝐑}{}𝐗(x)\rho ]=\text{Tr}[\rho ]$$
(27)
for all $`\rho \tau c(cH)`$. For any $`x𝐑`$, define the operator $`X(x)`$ by
$$X(x)=𝐗(x)^{}I.$$
(28)
We call $`X(x)`$ the effect of $`𝐀(𝐱)`$ for the outcome $`𝐱=x`$. The family $`\{X(x)|x𝐑\}`$ of the effects of $`𝐀(𝐱)`$ is called the effect distribution of the apparatus $`𝐀(𝐱)`$. From (21) the output distribution of the apparatus $`𝐀(𝐱)`$ is determined by the effect as
$$\mathrm{Pr}\{𝐱=x\rho \}=\text{Tr}[X(x)\rho ].$$
(29)
By the positivity of probability, we have $`\text{Tr}[X(x)\rho ]0`$. Since the density operator $`\rho `$ is arbitrary, $`X(x)`$ is a positive operator. From (26), we have
$$\underset{x𝐑}{}X(x)=I.$$
(30)
In this case, $`X(x)=0`$ except for countable number of $`x`$s. From (29), $`𝐀(𝐱)`$ measures an observable $`A`$ if and only if
$$X(x)=E^A(x).$$
(31)
Thus, $`𝐀(𝐱)`$ measures an observable if and only if the effect distribution coincides with its spectral projections. Otherwise, the apparatus $`𝐀(𝐱)`$ is interpreted to carry out a more general measurement such as an approximate measurement of an observable.
We define the positive superoperator $`𝐓`$ by
$$𝐓\rho =\underset{x𝐑}{}𝐗(x)\rho ,$$
(32)
where the sum is a countable sum since $`𝐗(x)=0`$ except for countable number of $`x`$s. This $`𝐓`$ is called the nonselective operation of the apparatus $`𝐀(𝐱)`$. From (30) we have $`𝐓^{}I=I`$ and hence $`𝐓`$ is a trace preserving positive superoperator.
A family $`\{𝐖(x)|x𝐑\}`$ of positive superoperators is called a superoperator distribution iff
$$\underset{x𝐑}{}𝐖(x)^{}I=I.$$
A family $`\{F(x)|x𝐑\}`$ of positive operators is called a operator distribution iff
$$\underset{x𝐑}{}F(x)=I.$$
The family $`\{E^A(x)|x\}`$ of spectral projections of a discrete self-adjoint operator $`A`$ is an operator distribution. The family $`\{W(x)|x𝐑\}`$ of positive operators defined by
$$W(x)=𝐖(x)^{}I$$
is an operator distribution and is called the operator distribution of $`\{𝐖(x)|x𝐑\}`$. The superoperator $`𝐓`$ definied by
$$𝐒=\underset{x𝐑}{}𝐖(x)$$
is a positive trace preserving superoperator and is called the total superoperator of $`\{𝐖(x)|x𝐑\}`$.
We have shown under the mixing law of the joint probability that the operational distribution $`\{𝐗(x)|x𝐑\}`$ of an apparatus $`𝐀(𝐱)`$ is a superoperator distribution, the effect distribution of $`𝐀(x)`$ is the operator distribution $`\{𝐗(x)|\rho \}`$, and the nonselective superoperator of $`𝐀(𝐱)`$ is the total superoperator of $`\{𝐗(x)|x𝐑\}`$. Conversely, if for given apparatuses $`𝐀(𝐱)`$ and $`𝐀(𝐲)`$ there are superoperator distributions $`\{𝐗(x)|x𝐑\}`$ and $`\{𝐘(y)|y𝐑\}`$ satisfy (20) respectively, then the joint probability distribution of the outcomes of the successive measurements carried out by $`𝐀(𝐱)`$ and $`𝐀(𝐲)`$ satisfies
$$\mathrm{Pr}\{𝐱=x,𝐲=y\rho \}=\text{Tr}[𝐘(y)𝐗(x)\rho ],$$
and hence the mixing law of the joint probability holds.
From the arguments so far, we conclude that the mixiing law of the joint porbability is equivalent with the following requirement: For any measuring apparatus $`𝐀(𝐱)`$, there is a superoperator distribution $`\{𝐗(x)|x𝐑\}`$ such that the statistical property of $`𝐀(𝐱)`$ is represented as follows.
output distribution: $`\mathrm{Pr}\{𝐱=x\rho \}=\text{Tr}[𝐗(x)\rho ]`$ (33)
output state: $`\rho _{\{𝐱=x\}}={\displaystyle \frac{𝐗(x)\rho }{\text{Tr}[𝐗(x)\rho ]}}`$ (34)
In (34) the outcome $`𝐱=x`$ is supposed to have positive probability; henceforce, the analogous assumption will be required implicitly in the similar expressions on the output state.
It follows that the problem as to what statistical property is possible is reduced to the problem as to what superoperator distributions are the operational distributions of apparatuses.
## V Davies-Lewis postulate
For the case of the discrete output variables, the notion of superoperator distributions is equivalent to the notion of normalized positive superoperator valued measures introduced by Davies and Lewis . A positive superoperator valued (PSV) measure is a mapping $``$ which maps every Borel set $`\mathrm{\Delta }`$ to a positive superoperator $`(\mathrm{\Delta })`$ such that if $`\mathrm{\Delta }_1,\mathrm{\Delta }_2,\mathrm{}`$ is a countable Borel partition of $`\mathrm{\Delta }`$, then we have
$$(\mathrm{\Delta })\rho =\underset{n}{}(\mathrm{\Delta }_n)\rho $$
for any $`\rho \tau c()`$, where the sum is convergent in the trace norm. The PSV measure $``$ is said to be normalized if it satisfies the further condition
$$\text{Tr}[(𝐑)\rho ]=\text{Tr}[\rho ]$$
for any $`\rho \tau c()`$. The equivalence is given below analogous to the case of discrete probability measures. If $``$ is a normalized PSV measure, then the corresponding superoperator distribution $`\{𝐗(x)|x𝐑\}`$ is given by
$$𝐗(x)=(\{x\}),$$
(35)
where $`\{x\}`$ is the singleton set containing the point $`x`$. Conversely, if $`\{𝐗(x)|x𝐑\}`$ is a superoperator distribution, then the corresponding normalized PSV measure is given by
$$(\mathrm{\Delta })=\underset{x\mathrm{\Delta }}{}𝐗(x).$$
(36)
For the apparatus $`𝐀(𝐱)`$, the probability $`\mathrm{Pr}\{𝐱\mathrm{\Delta }\rho \}`$ of obtaining the outcome in the Borel set $`\mathrm{\Delta }`$ is given by
$$\mathrm{Pr}\{𝐱\mathrm{\Delta }\rho \}=\underset{x\mathrm{\Delta }}{}\mathrm{Pr}\{𝐱=x\rho \}$$
(37)
and the output state of the ensemble of the samples with the outcome in the Borel set $`\mathrm{\Delta }`$ is given by
$$\rho _{\{𝐱\mathrm{\Delta }\}}=\frac{_{x\mathrm{\Delta }}\mathrm{Pr}\{𝐱=x\rho \}\rho _{\{𝐱=x\}}}{\mathrm{Pr}\{𝐱\mathrm{\Delta }\rho \}}.$$
(38)
Davies and Lewis proposed the following description of measurement statistics:
Davies-Lewis postulate: For any measuring apparatus $`𝐀(𝐱)`$, there is a normalized PSV measure $``$ satisfying the following relations for any density operator $`\rho `$ and Borel set $`\mathrm{\Delta }`$:
(DL1) $`\mathrm{Pr}\{𝐱\mathrm{\Delta }\rho \}=\text{Tr}[(\mathrm{\Delta })\rho ]`$
(DL2) $`\rho _{\{𝐱\mathrm{\Delta }\}}={\displaystyle \frac{(\mathrm{\Delta })\rho }{\text{Tr}[(\mathrm{\Delta })\rho ]}}`$
Although the Davies-Lewis description of measurement is quite general, it is not clear by itself whether it is general enough to exhaust all the possible measurements. Our arguments are about to complete proving the following theorem that shows indeed it is the case.
###### Theorem 2
The Davies-Lewis postulate is equivalent to the mixing law of the joint probability.
In fact, under the Davies-Lewis postulate, we have the normalized PSV measures $`_𝐱`$ and $`_𝐲`$ for any apparatuses $`𝐀(𝐱)`$ and $`𝐀(𝐲)`$. By substituting (DL1) and (DL2) in (9), the joint probability is given by
$$\mathrm{Pr}\{𝐱=x,𝐲=y\rho \}=\text{Tr}[_𝐲(\{y\})_𝐱(\{x\})\rho ].$$
From the linearity of $`_𝐱(\{x\})`$ and $`_𝐲(\{x\})`$, the mixing law follows. Conversely, under the mixing law, we have shown that there is a superoperator distribution $`\{𝐗(x)|x𝐑\}`$ satisfying (21) and (22). Now, it is easy to check that relations (35)–(38) leads to the Davies-Lewis description (DL1)–(DL2) and the proof is completed.
## VI Measurements of discrete observables
For a given discrete self-adjoint operator $`A`$, a superoperator distribution $`\{𝐗(x)|x𝐑\}`$ is called $`A`$-compatible iff $`𝐗(x)^{}I=E^A(x)`$ for all $`x𝐑`$. The operational distribution of an apparatus measuring the observable $`A`$ is an $`A`$-compatible superoperator distribution .
We have the following theorem ; a simplified proof will be given in Appendix B.
###### Theorem 3
Let $`A`$ be a discrete self-adjoint operator. Let $`\{𝐗(x)|x𝐑\}`$ be an $`A`$-compatible superoperator distribution and $`𝐓`$ its total superoperator. For any real number $`x`$ and trace class operator $`\rho `$, we have
$`𝐗(x)\rho `$ $`=`$ $`𝐓[E^A(x)\rho ]=𝐓[\rho E^A(x)]`$ (39)
$`=`$ $`𝐓[E^A(x)\rho E^A(x)].`$ (40)
For any real number $`x`$ and bounded operator $`B`$, we have
$`𝐗(x)^{}B`$ $`=`$ $`E^A(x)𝐓^{}(B)=𝐓^{}(B)E^A(x)`$ (41)
$`=`$ $`E^A(x)𝐓^{}(B)E^A(x).`$ (42)
From the above theorem, the operational distribution of an apparatus measuring an observable $`A`$ is determined uniquely by the nonselective operation. It follows from (41) that the range of $`𝐓^{}`$ consists of operators commuting with $`A`$. Let us define the commutant of $`A`$, denoted by $`\{A\}^{}`$, as the set of all bounded operators commuting with $`A`$. A trace preserving positive superoperator $`𝐋`$ on $`\tau c()`$ is called $`A`$-compatible iff the range of its duel $`𝐋^{}`$ is included in the commutant $`\{A\}^{}`$ of $`A`$.
For any trace preserving positive superoperator $`𝐋`$, let
$$𝐋^{}\rho =\underset{x𝐑}{}𝐋[E^A(x)\rho E^A(x)].$$
Then $`𝐋^{}`$ is an $`A`$-compatible positive superoperator. Obviously, $`𝐋`$ itself is $`A`$-compatible if and only if $`𝐋^{}=𝐋`$.
From (41), the total superoperator of an $`A`$-compatible superoperator distribution is an $`A`$-compatible positive superoperator. Conversely, for any $`A`$-compatible positive superoperator $`𝐓`$, let $`𝐗(x)\rho =𝐓[E^A(x)\rho ]`$ for all $`\rho \tau c()`$. Then $`\{𝐗(x)|x𝐑\}`$ is an $`A`$-compatible superoperator distribution and $`𝐓`$ is its total superoperator. From the above argument, we have obtained the following theorem.
###### Theorem 4
Let $`A`$ be a discrete self-adjoint operator on $``$. The relation
$$𝐗(x)\rho =𝐓[E^A(x)\rho ]$$
(43)
for all real number $`x`$ and trace class operator $`\rho `$ sets up a one-to-one correspondence between the $`A`$-compatible superoperator distribution $`\{𝐗(x)|x𝐑\}`$ and the $`A`$-compatible positive superoperators $`𝐓`$.
From the above theorem, we conclude the following: For any apparatus $`𝐀(𝐱)`$ measuring a discrete observable $`A`$, there is an $`A`$-compatible positive superoperator $`𝐓`$ such that the statistical property of $`𝐀(𝐱)`$ is represented as follows.
output distribution: $`\mathrm{Pr}\{𝐱=x\rho \}=\text{Tr}[E^A\rho ]`$ (44)
output state: $`\rho _{\{𝐱=x\}}={\displaystyle \frac{𝐓[E^A(x)\rho ]}{\text{Tr}[E^A(x)\rho ]}}`$ (45)
It follows that the problem of determining all the possible quantum state reductions $`\rho \rho _{\{𝐱=x\}}`$ arising in the apparatus measuring $`A`$ is reduced to the following problems: (i) Does every $`A`$-compatible positive superoperator have the corresponding measuring apparatus? (ii) If not, what condition does ensure the existence of the corresponding measuring apparatus?
## VII Measurements of nondegenerate discrete observables
In this section, we confine our attention to the observables with nondegenerate eigenvalues. In this case, the projection $`E^A(a_n)`$ is of rank 1 and is the density operator representing the eigenstate, so that we have
$$E^A(a_n)\rho E^A(a_n)=\text{Tr}[E^A(a_n)\rho ]E^A(a_n).$$
Let $`𝐓`$ be an $`A`$-compatible positive superoperator. From (39), we have
$$𝐓[E^A(a_n)\rho ]=\text{Tr}[E^A(a_n)\rho ]𝐓[E^A(a_n)].$$
(46)
We define a sequence $`\{\mathit{\varrho }_n\}`$ of density operators by
$$\mathit{\varrho }_n=𝐓[E^A(a_n)].$$
(47)
Then, we have
$$𝐓(\rho )=\underset{n}{}\text{Tr}[E^A(a_n)\rho ]\mathit{\varrho }_n.$$
(48)
Conversely, for any sequence $`\{\mathit{\varrho }_n\}`$ of density operators, we define the positive superoperator $`𝐓`$ on $`\tau c()`$ by (48). Then, $`𝐓`$ is an $`A`$-compatible positive superoperator satisfying (47). Thus, we have proved the following theorem.
###### Theorem 5
Let $`A`$ be a nondegenerate discrete self-adjoint operator. The relation
$$𝐓(\rho )=\underset{n}{}\text{Tr}[E^A(a_n)\rho ]\mathit{\varrho }_n,$$
where $`\rho \tau c()`$, sets up a one-to-one correspondence between the families $`\{\mathbf{\varrho }_x|x𝐑\}`$ of density operators and the $`A`$-compatible positive superoperators $`𝐓`$ on $`\tau c()`$.
Let $`\rho _{\{𝐱=x\}}`$ be the output state of an apparatus measuring $`A`$ for the input state $`\rho `$. Then, there is an $`A`$-compatible positive superoperator $`𝐓`$ satisfying (45) and there is a sequence $`\{\mathit{\varrho }_n\}`$ of density operators satisfying (47), so that we have
$$\rho _{\{𝐱=a_n\}}=\frac{𝐓[E^A(a_n)\rho ]}{\text{Tr}[E^A(a_n)\rho ]}=𝐓[E^A(a_n)]=\mathit{\varrho }_n.$$
It follows that the output state for the output $`𝐱=a_n`$ is given by
$$\rho _{\{𝐱=a_n\}}=\mathit{\varrho }_n.$$
(49)
From the above argument we conclude the following: For any apparatus $`𝐀(𝐱)`$ measuring a nondegenerate discrete observable $`A=_na_n|\varphi _n\varphi _n|`$, there is a sequence $`\{\mathbf{\varrho }_𝐧\}`$ of density operators such that the statisitcal property of $`𝐀(𝐱)`$ is represented as follows.
output distribution: $`\mathrm{Pr}\{𝐱=a_n\rho \}=\varphi _n|\rho |\varphi _n`$ (50)
output state: $`\rho _{\{𝐱=a_n\}}=\mathit{\varrho }_n`$ (51)
It follows that the problem of determining all the possible quantum state reductions arising in the measurement of a nondegenerate discrete observable $`A`$ is reduced to the problem as to what sequence $`\{\mathit{\varrho }_n\}`$ of states can be obtained from the measurement of $`A`$. In order to obtain the answer to this question, in the next section we shall consider indirect measurement models and ask what sequences can be obtained from those models.
It should be noted here that the apparatus satisfies the projection postulate if and only if we have
$$\mathit{\varrho }_𝒏=E^A(a_n)$$
for all $`n`$. Von Neumann \[1, pp. 439–442\] showed that this case can be obtained from an indirect measurement model.
## VIII Indirect measurement models
In general, if a measurement on the object in the input state $`\rho `$ by the apparatus $`𝐀(𝐱)`$ is immediately followed by a measurement of the observable $`B`$ by the apparatus $`𝐀(𝐛)`$, the joint probability distribution of their output variables is given by (15). Now, consider the marginal probability
$$\mathrm{Pr}\{𝐱𝐑,𝐛=b_m\rho \}=\underset{x𝐑}{}\mathrm{Pr}\{𝐱=x,𝐛=b_m\rho \}.$$
(52)
Then, this represents the probability of obtaining the outcome $`𝐛=b_m`$ after interacting the apparatus $`𝐀(𝐱)`$ with the object without reading out the outcome of the $`𝐱`$ measurement. Such a process of is called a nonselective measurement. Let $`𝐓`$ be the nonselective operation of the apparatus $`𝐀(𝐱)`$. Then, by (15), we have
$$\mathrm{Pr}\{𝐱𝐑,𝐛=b_m\rho \}=\text{Tr}[E^B(b_m)𝐓\rho ].$$
(53)
Thus, the nonselective measurement transforms the input state $`\rho `$ to the output state $`\rho _{\{𝐱𝐑\}}=𝐓\rho `$.
Let us call any interaction between the object and the apparatus caused by a measurement as the measuring interaction. Then, the superoperator $`𝐓`$ is determined by the measuring interaction. In what follows we shall examine the properties of the measuring interaction.
Since the nonselective measurement transforms the input state $`\rho `$ to the output state $`𝐓\rho `$, there should be an interaction during finite time interval when the object changes from $`\rho `$ to $`𝐓\rho `$. Moreover, the object should be free from the apparatus before and after the interaction. Thus, we suppose that the measuring interaction turns on from the time $`t`$ just before the measurement to the time $`t+\mathrm{\Delta }t`$ just after the measurement where $`\mathrm{\Delta }t>0`$, and that the object is free from the apparatus before the the time $`t`$ and after the time $`t+\mathrm{\Delta }t`$. It follows that if the second measurement on the same object follows immediately after the above measurement, the time just before the second measurement coincides with the time $`t+\mathrm{\Delta }t`$ just after the first measurement. In this way, the temporal boundary of the measuring interaction is determined as a fixed domain from time $`t`$ to $`t+\mathrm{\Delta }t`$.
Next, in order to determine the spatial boundary of the measuring interaction, we consider the smallest subsystem of the measuring apparatus such that the composite system of the object and the subsystem is isolated from the time $`t`$ to the time $`t+\mathrm{\Delta }t`$. We call the above subsystem as the prove.
The effect of the measuring interaction is given by the change of an observable $`M`$, called the probe observable, from $`t`$ to $`t+\mathrm{\Delta }t`$. From the minimality of the probe, it is natural to assume that the interaction Hamiltonian excludes any macroscopic part of the measuring apparatus such as the macroscopic pointer position. It follows that the measuring interaction is a quantum mechanical interaction and the state change can be described by the unitary time evolution of the composite system of the object and the probe.
On the other hand, in order to transduce the microscopic change in the probe observable $`M`$ to the macroscopic change such as the change of the position of the pointer, we need an amplification process in the apparatus after $`t+\mathrm{\Delta }t`$. This transduction from a microscopic observable to a macroscopic observable corresponds to the direct measurement of the probe observable $`M`$ at the time $`t+\mathrm{\Delta }t`$. The problem of describing this process as a dynamical process belongs to the so-called measurement problem. Within quantum mechanics, the Born statistical formula gives the the probability distribution of the outcome of the $`M`$ measurement. Let $`t+\mathrm{\Delta }t+\tau `$ be the time just after this amplification process where $`\tau >0`$. This time is called the time of read-out.
According to the above description, the process from the time just before the measurement to the time of read-out is divided into the measuring interaction and the amplification. It should be noted that just after the measuring interaction, the object is free from the apparatus so that it is possible to start the interaction with the second apparatus. It follows that in the successive measurement experiment the time just before the second measurement is considered to be the time just after the measuring interaction rather than the time of read-out. The above description of measuring process is called an indirect measurement description.
Let $``$ be the state space of the object $`𝐒`$, and $`𝒦`$ the state space of the probe $`𝐏`$. The state of the object at the time $`t`$ of measurement is the input state $`\rho `$. The probe $`𝐏`$ is supposed to be prepared in the fixed state $`\sigma `$ at the time of measurement. Thus, the state of the composite system at the time $`t`$ is
$$𝝆_{𝐒+𝐏}(t)=\rho \sigma .$$
If the time evolution of the composite system $`𝐒+𝐏`$ from $`t`$ to $`t+\mathrm{\Delta }t`$ is represented by the unitary operator $`U`$, the composite system is in the state
$$𝝆_{𝐒+𝐏}(t+\mathrm{\Delta }t)=U(\rho \sigma )U^{}$$
(54)
at $`t+\mathrm{\Delta }t`$. Suppose that the $`𝐀(𝐱)`$ measurement in $`\rho `$ is followed immediately by a measurement of an observable $`B`$ carried out by $`𝐀(𝐛)`$. Then, the observable $`B`$ is measured at the time $`t+\mathrm{\Delta }t`$ and the outcome is recorded by $`𝐛`$. On the other hand, the probe observable $`M`$ is also measured actually at the time $`t+\mathrm{\Delta }t`$ and the outcome is recorded by $`𝐱`$. Since the two measurements are carried out locally, it follows from the local measurement theorem that the joint probability distribution of the outcomes of the above two measurements satisfies
$`\mathrm{Pr}\{𝐱=x,𝐛=b_m\rho \}`$ (55)
$`=`$ $`\text{Tr}[(E^B(b_m)E^M(x))U(\rho \sigma )U^{}]`$ (56)
$`=`$ $`\text{Tr}[E^B(b_m)\text{Tr}_𝒦[(IE^M(x))U(\rho \sigma )U^{}]],`$ (57)
where $`\text{Tr}_𝒦`$ is the partial trace over the Hilbert space $`𝒦`$. Thus, from (15) we have
$$𝐗(x)\rho =\text{Tr}_𝒦[(IE^M(x))U(\rho \sigma )U^{}].$$
(58)
Hence, the statistical property of the apparatus $`𝐀(𝐱)`$ is given as follows.
output distribution: (60)
$`\mathrm{Pr}\{𝐱=x\rho \}=\text{Tr}[(IE^M(x))U(\rho \sigma )U^{}]`$
output state: (62)
$`\rho _{\{𝐱=x\}}={\displaystyle \frac{\text{Tr}_𝒦[(IE^A(x))U(\rho \sigma )U^{}]}{\text{Tr}[(IE^A(x))U(\rho \sigma )U^{}]}}`$
From (32) and (58), the nonselective operation of $`𝐀(𝐱)`$ is given by
$$𝐓\rho =\text{Tr}_𝒦[U(\rho \sigma )U^{}].$$
(63)
From (29) and (60), the effect distribution of $`𝐀(𝐱)`$ is given by
$$X(x)=\text{Tr}_𝒦[U^{}(IE^M(x))U(I\sigma )].$$
(64)
In general, a four tuple $`(𝒦,\sigma ,U,M)`$ is called an indirect measurement model iff it consists of a separable Hilbert space $`𝒦`$œ$`B!`$œ(B a density operator $`\sigma `$ on $`𝒦`$, a unitary operator $`U`$ on $`𝒦`$, and a self-adjoint operator $`M`$ on $`𝒦`$. So far we have not posed any sufficient condition for the existence of an apparatus except that every observable has at least one apparatus to measure it. Here, we pose the following hypothesis.
Unitary realizability hypothesis: For any indirect measurement model $`(𝒦,\sigma ,U,M)`$, there is an apparatus $`𝐀(𝐱)`$ with the following statistical property:
output distribution:
$`\mathrm{Pr}\{𝐱=x\rho \}=\text{Tr}[(IE^M(x))U(\rho \sigma )U^{}]`$
output state:
$`\rho _{\{𝐱=x\}}={\displaystyle \frac{\text{Tr}_𝒦[(IE^A(x))U(\rho \sigma )U^{}]}{\text{Tr}[(IE^A(x))U(\rho \sigma )U^{}]}}`$
A superoperator distribution $`\{𝐗(x)|x𝐑\}`$ is said to be realized by an indirect measurement model $`(𝒦,\sigma ,U,M)`$ iff (58) holds for any $`\rho \tau c()`$, and in this case it is called unitarily realizable. Under the unitary realizability hypothesis, unitarily realizable superoperator distributions are operational distributions of some apparatuses. In the next section, we shall give an intrinsic characterization of the unitarily realizable superoperator distributions.
## IX Complete positivity
Let $`𝒟=\tau c()`$ or $`𝒟=()`$. A linear transformation $`𝐋`$ on $`𝒟`$ is called completely positive (CP) iff for any finite sequences of bounded operators $`A_1,\mathrm{},A_n𝒟`$ and vectors $`\xi _1,\mathrm{},\xi _n`$ we have
$$\underset{ij}{}\xi _i|𝐋(A_i^{}A_j)|\xi _j0.$$
The above condition is equivalent to that $`𝐋I`$ maps positive operators in the algebraic tensor product $`𝒟(𝒦)`$ to positive operators in $`𝒟(𝒦)`$ for any Hilbert space $`𝒦`$. Obviously, every CP superoperators are positive. A superoperator is CP if and only if its dual superoperator is CP. A superoperator distribution $`\{𝐗(x)|x𝐑\}`$ is called completely positive iff every $`𝐗(x)`$ is CP. It can be seen easily from (58) that unitarily realizable superoperator distributions are CP. Conversely, the following theorem, proved in for an even more general formulation, asserts that every CP superoperator distribution is unitarily realizable.
###### Theorem 6
For any CP superoperator distribution $`\{𝐗(x)|x𝐑\}`$, there is a separable Hilbert space $`𝒦`$, a unit vector $`\mathrm{\Phi }`$ in $`𝒦`$, a unitary operator $`U`$ on $`𝒦`$, and a discrete self-adjoint operator $`M`$ on $`𝒦`$ satisfying the relation
$$𝐗(x)\rho =\text{Tr}_𝒦[(IE^M(x))U(\rho \sigma )U^{}].$$
for all $`\rho \tau c(cH)`$.
For any trace preserving CP superoperator $`𝐓`$, we have a CP superoperator distribution $`\{𝐗(x)|x𝐑\}`$ such that $`𝐗(0)=𝐓`$ and that $`𝐗(x)=0`$ for all $`x0`$. Applying the above theorem to this family, we obtain the following representation theorem of trace preserving CP superoperators, which was proved independently by Kraus and the present author .
###### Theorem 7
For any trace preserving CP superoperator $`𝐓`$, there is a separable Hilbert space $`𝒦`$, a unit vector $`\mathrm{\Phi }`$ in $`𝒦`$, a unitary operator $`U`$ on $`𝒦`$, such that $`𝐓`$ satisfies the relation
$$𝐓\rho =\text{Tr}_𝒦[U(\rho |\mathrm{\Phi }\mathrm{\Phi }|)U^{}].$$
for all $`\rho \tau c()`$.
From Theorem 3, every $`A`$-compatible superoperator distribution $`\{𝐗(x)|𝐑\}`$ satisfies the relation $`𝐗(x)\rho =𝐓[E^A(x)\rho E^A(x)]`$. Thus, if $`\{𝐗(x)|x𝐑\}`$ is CP, then the total map $`𝐓=_{x𝐑}𝐗(x)`$ is CP, since the sum of CP superoperators is CP. Conversely, if $`𝐓`$ is an $`A`$-compatible CP superoperator, then the corresponding $`A`$-compatible superoperator distribution $`\{𝐗(x)|x𝐑\}`$ is CP, since the superoperator $`\rho E^A(x)\rho E^A(x)`$ is CP and the composition of any CP superoperators is CP. Thus we have the following:
###### Theorem 8
Let $`A`$ be a nondegenerate discrete self-adjoint operator. Then, an $`A`$-compatible superoprator distribution is CP if and only if its total superoperator is CP.
From the above theorem, we conclude the following : The statistical equivalence classes of apparatuses $`𝐀(𝐱)`$ measuring a discrete observable $`A`$ with indirect measurement models are in one-to-one correspondence with the $`A`$-compatible CP superoperators, where the statistical property is represented by (44) and (45).
Now, let $`A`$ be a nondegenerate discrete observable and let $`𝐓`$ be an $`A`$-compatible positive superoperator. Then, $`𝐓`$ is of the form (48). Let $`\sigma _1,\mathrm{},\sigma _n\tau c()`$ and $`\xi _1,\mathrm{},\xi _n`$. Then, we have
$`{\displaystyle \underset{ij}{}}\xi _i|𝐓(\sigma _i^{}\sigma _j)|\xi _j`$ $`=`$ $`{\displaystyle \underset{n}{}}{\displaystyle \underset{ij}{}}\text{Tr}[E^A(a_n)\sigma _i^{}\sigma _j]\xi _i|\mathit{\varrho }_n|\xi _j`$
$``$ $`0,`$
where the last inequality follows from the fact that the trace of the product of two positive definite matrices $`(\text{Tr}[E^A(a_n)\sigma _i^{}\sigma _j])_{ij}`$ and $`(\xi _i|\mathit{\varrho }_n|\xi _j)_{ij}`$ is nonnegative. It follows that $`𝐓`$ is a CP superoperator. Thus, every $`A`$-compatible superoperator is CP. Since every $`A`$-compatible superoperator distribution is obtained from an $`A`$-compatible superoperator by Theorem 4, it follows from Theorem 8 that every $`A`$-compatible superoperator distribution is CP. We have therefore obtained the following statements.
###### Theorem 9
Let $`A`$ be a nondegenerate discrete self-adjoint operator. Every $`A`$-compatible positive superoperator is completely positive. Every $`A`$-compatible superoperator distribution is completely positive.
From the above theorem and Theorem 6 we conclude: Every apparatus measuring $`A`$ is statistically equivalent to the one having an indirect measurement model.
Every sequence $`\{\mathit{\varrho }_n\}`$ of density operators defines an $`A`$-compatible positive superoperator by Theorem 5, and it is automatically completely positive so that it is realized by an idirect measurement model. Thus, we have reached the answer to the question what sequence of states can be obtained from an apparatus measuring $`A`$ that every sequence can. Thus, we conclude the following: The statistical equivalence classes of apparatuses $`𝐀(𝐱)`$ measuring a nondegenerate discrete observable $`A`$ are in one-to-one correspondence with the sequences $`\{\mathbf{\varrho }_n\}`$ of density operators, where the statistical property is represented by (50) and (51).
Given any sequence $`\{\mathit{\varrho }_n\}`$, an indirect measurement model with the quantum state reduction
$$\rho \rho _{\{𝐱=a_n\}}=\mathit{\varrho }_n$$
is constructed explicitly as follows. Let $`\{\varphi _n\}`$ be an orthonormal basis of $``$ consisting of the eigenvectors of $`A`$. Let $`𝒦=`$. Let
$$\mathit{\varrho }_n=\underset{j}{}\lambda _{nj}|\eta _{nj}\eta _{nj}|$$
be the spectral decomposition of $`\mathit{\varrho }_n`$. Then, there exists a unitary operator $`U`$ on $`𝒦`$ satisfying
$$U|\varphi _n\varphi _0\varphi _0=\underset{j}{}\sqrt{\lambda _{nj}}|\eta _{nj}\varphi _j\varphi _n.$$
Now, we define the density operator $`\sigma `$ on $`𝒦`$ by $`\sigma =|\varphi _0\varphi _0\varphi _0\varphi _0|`$ and define a self-adjoint operator $`M`$ on $`𝒦`$ by $`M=IA`$. Then, we have the indirect measurement model $`(𝒦,\sigma ,U,M)`$ such that the statistical property of its apparatus satisfies (50) and (51).
## X Conclusions
Let $`𝐀(𝐱)`$ be an apparatus with the discrete output variable $`𝐱`$. Then, depending on the input state $`\rho `$ and the outcome $`x`$, the apparatus $`𝐀(𝐱)`$ determines the output probability $`\mathrm{Pr}\{𝐱=x\rho \}`$ and the output state $`\rho _{\{𝐱=x\}}`$. The transformation from the input state $`\rho `$ to the output distribution $`\mathrm{Pr}\{𝐱=x\rho \}`$ is called the objective state reduction and the one from the input state $`\rho `$ to the output states $`\rho _{\{𝐱=x\}}`$ is called the quantum state reduction. The pair of the objective state reduction and the quantum state reduction is called the statistical property of the apparatus $`𝐀(𝐱)`$. Two apparatuses with the same statistical property is said to be statistically equivalent. In order to obtain a mathematical description of quantum state reductions for the most general class of measurements we have considered two requirements: one is necessary and the other is sufficient.
The necessary one is the mixing law of the joint probability. Suppose that a measurement carried out by an apparatus $`𝐀(𝐱)`$ in the input state $`\rho `$ is followed immediately by another measurement carried out by another apparatus $`𝐀(𝐲)`$. The joint probability distribution of the outcomes $`𝐱`$ and $`𝐲`$ is determined by their statistical properties as follows.
$$\mathrm{Pr}\{𝐱=x,𝐲=y\rho \}=\mathrm{Pr}\{𝐲=y\rho _{\{𝐱=x\}}\}\mathrm{Pr}\{𝐱=x\rho \}.$$
This joint probability distribution is considered to respect the mixture of input states and the mixing law of the joint probability requires that this is the case for any apparatuses $`𝐀(𝐱)`$ and $`𝐀(𝐲)`$. Under this hypothesis, any apparatus $`𝐀(𝐱)`$ has a superoperator distribution $`\{𝐗(x)|x𝐑\}`$, called the operational distribution of $`𝐀(𝐱)`$, satisfying
$$𝐗(x)\rho =\mathrm{Pr}\{𝐱=a\rho \}\rho _{\{𝐱=x\}}.$$
(65)
The sufficient condition is the unitary realizability condition. The apparatus $`𝐀(𝐱)`$ is said to have an indirect measurement model $`(𝒦,\sigma ,U,M)`$ iff the statistical property of $`𝐀(𝐱)`$ is given as follows.
output distribution:
$`\mathrm{Pr}\{𝐱=x\rho \}=\text{Tr}[(IE^M(x))U(\rho \sigma )U^{}]`$
output state:
$`\rho _{\{𝐱=x\}}={\displaystyle \frac{\text{Tr}_𝒦[(IE^M(x))U(\rho \sigma )U^{}]}{\text{Tr}[(IE^M(x))U(\rho \sigma )U^{}]}}`$
In general, an apparatus has an indirect measurement model if and only if its operational distribution is completely positive. The unitary realizability hypothesis states that every indirect measurement model defines an apparatus with the above statistical property. It follows that the statistical equivalence classes of apparatuses with indirect measurement models are in one-to-one correspondence with the CP superoperator distributions $`\{𝐱(x)|x𝐑\}`$, under the relation (65).
Let $`A`$ be a discrete observable. A trace preserving positive superoperator $`𝐋`$ is called $`A`$-compatible iff the range of its duel $`𝐋^{}`$ is included in the commutant $`\{A\}^{}`$ of $`A`$. The statistical property of an apparatus measuring an observable $`A`$ is represented by an $`A`$-compatible positive superoperator $`𝐓`$ as follows.
output distribution: $`\mathrm{Pr}\{𝐱=x\rho \}=\text{Tr}[E^A\rho ]`$ (66)
output state: $`\rho _{\{𝐱=x\}}={\displaystyle \frac{𝐓[E^A(x)\rho ]}{\text{Tr}[E^A(x)\rho ]}}`$ (67)
In particular, the statistical equivalence classes of apparatuses with indirect measurement models measuring $`A`$ are in one-to-one correspondence with the $`A`$-compatible completely positive superoperators $`𝐓`$, under the above description.
According to the above, the class of possible quantum state reductions is included in the class of conditional state transformations satisfying the mixing law, i.e., the general superoperator distributions, and includes the one satisfying the unitary realizability condition, i.e., the completely positive superoperator distributions. Since these two classes are generally different, there seems to be still a room for the debate in measurement theory on what class between them is the true class of all the possible quantum sate reductions.
Nevertheless, for the case where $`A`$ is nondegenerate, this paper shows, the above two conditions are actually equivalent. Thus, both of them are necessary and sufficient and we reach a clear-cut conclusion. In fact, if $`A`$ is nondegenerate, all the $`A`$-compatible positive superoperators $`𝐓`$ are completely positive and they are in one-to-one correspondence with the sequences $`\{\mathit{\varrho }_n\}`$ of density operators, under the relation $`𝐓(E^A(a_n))=\mathit{\varrho }_n`$ where $`\{a_n\}`$ is the sequence of the eigenvalues of $`A`$. In this case, every apparatus measuring $`A`$ is statistically equivalent with the one with an indirect measurement model. The statistical equivalence classes of the apparatuses measuring $`A`$ are, therefore, in one-to-one correspondence with the sequences $`\{\mathit{\varrho }_n\}`$ of density operators and their statistical properties are represented as follows.
output distribution: $`\mathrm{Pr}\{𝐱=a_n\rho \}=\text{Tr}[E^A(a_n)\rho ]`$ (68)
output states: $`\rho _{\{𝐱=a_n\}}=\mathit{\varrho }_n`$ (69)
The above measurement statistics has the following two remarkable features: (i) The output states are independent of the input state. (ii) The family of output states can be arbitrarily chosen by the choice of the apparatus. The possibility of this kind of generalized measurements was first pointed out in part by Gordon and Louisell relative to the measurement of an overcomplete family of states generalizing the conventional measurement of an orthonormal basis. Yuen generalized the Gordon-Louisell description to the following measurement described by the set of operators $`\{|\mathrm{\Psi }_x\mathrm{\Phi }_x|\}`$, where $`\{\mathrm{\Phi }_x\}`$ is an overcomplete family of vectors and $`\{\mathrm{\Psi }_x\}`$ is a Borel family of state vectors, as follows.
output distribution: $`\mathrm{Pr}\{𝐱dx\rho \}=\mathrm{\Phi }_x|\rho |\mathrm{\Phi }_xdx`$
output states: $`\rho _{\{𝐱=x\}}=|\mathrm{\Psi }_x\mathrm{\Psi }_x|`$
The unitary realizability of the above measurement statistics was assumed by Yuen to claim the realizability of the contractive state measurement and proved rigorously in ; see for survey. We can see that for the nondegenerate discrete observable $`A=_na_n|\mathrm{\Phi }_n\mathrm{\Phi }_n|`$ and the output states $`\mathit{\varrho }_n=|\mathrm{\Psi }_n\mathrm{\Psi }_n|`$ the measurement statistics given in (68) and (69) corresponds to the (discrete version of) measurement described by $`\{|\mathrm{\Psi }_n\mathrm{\Phi }_n|\}`$. The present paper has proved rigorously, even without assuming the unitary realizability, that every measurement of a nondegenerate discrete observable is always of this form.
Along with the analogous arguments, it can be shown that the statistical equivalence classes of the apparatuses measuring a nondegenerate (but not necessarily discrete) observable including the position or the momentum observable are in one-to-one correspondence with the Borel families of density operators (modulo the spectral measure). Since the precise mathematical formulation for that result is beyond the scope of this paper, we shall discuss the nondiscrete case in a separate article.
Therefore, we can conclude that as long as the statistical properties of measurements of nondegenerate observables are concerned, we can always assume that the measuring process are described by an indirect measurement model in which the interaction between the object and the apparatus is described by a unitary operator. For measurements of degenerate observables and even for measurements of general probability operator valued measures, it appears to be an important question whether every apparatus is statistically equivalent with the one having the indirect measurement model that has the unitary measuring interaction. Since in this case there are many superoperator distributions (or normalized PSV measures) that are not completely positive , we need further physical requirements to settle this problem.
Following von Neumann , some authors appear to support the hypothesis that every apparatus has an indirect measurement model, the converse of the unitary realizability hypothesis. If this is the case, the description of measuring processes will be simplified considerably as shown in Section VIII. In particular, we have an instant of time at which the measuring process is divided into the measuring interaction and the amplification process (including the so-called decoherence process) and the output state has been prepared for the next measurement before the amplification mode of the first measurement . It is also interesting whether non-conventional quantum mechanics such as nonlinear quantum mechanics will provide a different measurement statistics from the unitarily realizable ones.
## A Linear extension of the quantum state reduction
For any $`x𝐑`$ and any density operator $`\rho `$, the trace class operator $`𝐗(x,\rho )`$ is defined by (13). In this section, we shall prove that the mapping $`𝐗(x):\rho 𝐗(x,\rho )`$ defined on the space of density operators can be extended uniquely to a linear transformation on the space $`\tau c()`$ of trace class operators on $``$. By the linearity of the extension, for any trace class operator $`\sigma `$ with decomposition (17) it is necessary for $`𝐗(x)\sigma `$ to be defined by (19). Since the decomposition $`(\text{17})`$ is not unique, in order for the extension (19) to be well-defined we need to show that the right hand side of (19) is uniquely determined independent of the decomposition of $`\sigma `$. Namely, we need to prove that if $`\sigma `$ has another decomposition
$$\sigma =\lambda _1^{}\sigma _1^{}\lambda _2^{}\sigma _2^{}+i\lambda _3^{}\sigma _3^{}i\lambda _4^{}\sigma _4^{},$$
(A1)
then we have
$`\lambda _1𝐗(x)\sigma _1\lambda _2𝐗(x)\sigma _2+i\lambda _3𝐗(x)\sigma _3i\lambda _4𝐗(x)\sigma _4`$ (A2)
$`=`$ $`\lambda _1^{}𝐗(x)\sigma _1^{}\lambda _2^{}𝐗(x)\sigma _2^{}+i\lambda _3^{}𝐗(x)\sigma _3^{}i\lambda _4^{}𝐗(x)\sigma _4^{}.`$ (A3)
The proof runs as follows . By equating the right hand sides of (17) and (A1) and comparing the real and imaginary parts in both sides, we have
$`\lambda _1\sigma _1+\lambda _2^{}\sigma _2^{}`$ $`=`$ $`\lambda _1^{}\sigma _1^{}+\lambda _2\sigma _2`$ (A5)
$`\lambda _3\sigma _3+\lambda _4^{}\sigma _4^{}`$ $`=`$ $`\lambda _3^{}\sigma _3^{}+\lambda _4\sigma _4.`$ (A6)
Taking the trace of both sides of (A5), we have
$$\lambda _1+\lambda _2^{}=\lambda _1^{}+\lambda _2.$$
(A7)
By dividing both sides of (A5) by this value, we have
$$\alpha \sigma _1+(1\alpha )\sigma _2^{}=\beta \sigma _1^{}+(1\beta )\sigma _2,$$
where we define $`\alpha `$ and $`\beta `$ by
$`0<\alpha `$ $`=`$ $`{\displaystyle \frac{\lambda _1}{\lambda _1+\lambda _2^{}}}<1`$
$`0<\beta `$ $`=`$ $`{\displaystyle \frac{\lambda _1^{}}{\lambda _1^{}+\lambda _2}}<1.`$
Thus, from (16) we have
$$\alpha 𝐗(x)\sigma _1+(1\alpha )𝐗(x)\sigma _2^{}=\beta 𝐗(x)\sigma _1^{}+(1\beta )𝐗(x)\sigma _2.$$
Multiplying both sides by the value of (A7), we have
$$\lambda _1𝐗(x)\sigma _1\lambda _2𝐗(x)\sigma _2=\lambda _1^{}𝐗(x)\sigma _1^{}\lambda _2^{}𝐗(x)\sigma _2^{}.$$
By the similar manipulations for (A6), we have
$$i\lambda _3𝐗(x)\sigma _3i\lambda _4𝐗(x)\sigma _4=i\lambda _3^{}𝐗(x)\sigma _3^{}i\lambda _4^{}𝐗(x)\sigma _4^{}.$$
Thus, we have proved equation (A2). It is concluded, therefore, that $`𝐗(x)\sigma `$ is defined uniquely for every $`\sigma `$ by (19).
## B Proof of Theorem 3
Let $`\{𝐗(x)|x𝐑\}`$ be an $`A`$ compatible family of positive maps and $`𝐓`$ its total map. Let $`C`$ be a bounded operator such that $`0CI`$ and let $`x𝐑`$. We define
$`A_{11}`$ $`=`$ $`𝐗(x)^{}C,`$
$`A_{12}`$ $`=`$ $`𝐗(x)^{}(IC),`$
$`A_{21}`$ $`=`$ $`{\displaystyle \underset{yx}{}}𝐗(y)^{}C,`$
$`A_{22}`$ $`=`$ $`{\displaystyle \underset{yx}{}}𝐗(y)^{}(IC),`$
$`P_1`$ $`=`$ $`E^A(x),`$
$`P_2`$ $`=`$ $`IE^A(x),`$
$`Q_1`$ $`=`$ $`𝐓^{}(C),`$
$`Q_2`$ $`=`$ $`I𝐓^{}(C).`$
Then $`0A_{ij}P_i`$, so that $`[A_{ij},P_i]=[A_{ij},P_j]=0`$. It follows that $`Q_j=A_{1j}+A_{2j}`$ commutes with $`P_1`$ and $`P_2`$ as well. Thus,
$$A_{ij}=P_iA_{ij}P_iQ_j.$$
On the other hand, we have $`_{ij}A_{ij}=I`$ and $`_{ij}P_iQ_j=I`$, whence $`A_{ij}=P_iQ_j`$. It follows that $`𝐗(x)^{}C=E^A(x)𝐓^{}(C)`$. Since any bounded operator $`B`$ can be represented by $`B=_{n=0}^3i^n\lambda _nC_n`$ with positive operators $`0C_nI`$ and positive reals $`\lambda _n`$, we have $`𝐗(x)^{}B=E^A(x)𝐓^{}(B)`$ for any real number $`x`$ and bounded operator $`B`$. Since $`[E^A(x),𝐓^{}(B)]=0`$, other assertions follow immediately.
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# 1 Lines from the top to the bottom correspond to 𝑄²={10⁷,10⁶,10⁵} etc.
Quasinormal modes of Reissner-Nordstr$`\ddot{o}`$m Anti-de Sitter black holes
Bin Wang<sup>a,b,</sup><sup>1</sup><sup>1</sup>1e-mail:binwang@fma.if.usp.br, Chi-Yong Lin<sup>a,</sup><sup>2</sup><sup>2</sup>2e-mail:lcyong@fma.if.usp.br and Elcio Abdalla<sup>a,</sup><sup>3</sup><sup>3</sup>3e-mail:eabdalla@fma.if.usp.br
<sup>a</sup> Instituto De Fisica, Universidade De Sao Paulo, C.P.66.318, CEP 05315-970, Sao Paulo, Brazil
<sup>b</sup> Department of Physics, Shanghai Teachers’ University, P. R. China
## Abstract
Complex frequencies associated with quasinormal modes for large Reissner-Nordstr$`\ddot{o}`$m Anti-de Sitter black holes have been computed. These frequencies have close relation to the black hole charge and do not linearly scale with the black hole temperature as in Schwarzschild Anti-de Sitter case. In terms of AdS/CFT correspondence, we found that the bigger the black hole charge is, the quicker for the approach to thermal equilibrium in the CFT. The properties of quasinormal modes for $`l>0`$ have also been studied.
PACS number(s): 04.30.Nk, 04.70.Bw
Quasinomal modes of black holes have been an intriguing subject of discussions for the last few decades. It has become an evidence that the quasinormal ringing will dominate most processes involving perturbed black holes. This means that quasinormal modes will carry a unique fingerprint which would lead to the direct identification of the black hole existence. Detection of these quasinormal modes is expected to be realized through gravitational wave observation in the near future. In order to extract as much information as possible from gravitational wave signal, it is important that we understand exactly how the quasinormal modes behave for the parameters of black holes in different models. For black holes in asymptotically flat spacetime, especially spherical cases, they have been studied extensively, for a review see . The study for the nonspherical black holes is developing . Considering the case when the black hole is immersed in an expanding universe, the quasinormal modes of black holes in de Sitter space have also been investigated recently . It was found that there are qualitative differences from the asymptotically flat case, in particular the scalar field decay is always exponential, rather than a power-law tail in asymptotically flat spacetime. This result has also been uncovered by a very recent study of the Schwarzschild AdS black hole model . These observations support the earlier argument by Ching et al. that usual inverse power-law tails as seen in asymptotically flat black hole spacetime, are not a general feature of wave propagation in curved spacetime.
Motivated by the recent discovery of the AdS/CFT correspondence, the investigation of the quasinormal modes of AdS black holes becomes more appealing nowdays. The quasinormal frequencies of AdS black hole have direct interpretation in terms of the dual conformal field theory (CFT). In terms of the AdS/CFT correspondence \[8-10\], a large black hole corresponds to an approximately thermal state in the field theory, and the decay of the scalar field corresponds to the decay of perturbation of the state. After computing the scalar quasinormal modes of Schwarzschild AdS black holes in four, five and seven dimensions, Horowitz and Hubeny claimed that for large black hole both the real and the imaginary parts of the quasinormal frequences scale linearly with the black hole temperature. The timescale for approaching to the thermal equilibrium is determined by the imaginary part of the lowest quasinormal frequency and is proportional to the inverse of the black hole temperature. However, for a small black hole the results are no longer the same as that for a large black hole. The quasinormal frequencies do not continue to scale with temperature. Some comments on small black holes were presented therein.
The Schwarzschild AdS black hole studied in is the simplest model in Anti-de Sitter space, which is determined by only two dimensionful parameters, the black hole event horizon $`r_+`$ and the AdS radius $`R`$ relating to the cosmological constant by $`\mathrm{\Lambda }=3/R^2`$. It is of interest to generalize the study of ref. to a more general model, say Reissner-Nordstr$`\ddot{o}`$m (RN) AdS black holes. Besides $`r_+`$ and $`R`$, the RN AdS black hole has another parameter, the charge $`Q`$. Thus it possesses richer physics to be explored. In this paper we are going to study RN AdS black hole in the hope of getting more understanding of how the quasinormal modes depend on this additional parameter and whether there is more information on AdS/CFT correspondence.
The RN black hole solution of Einstein’s equations in free space with a negative cosmological constant $`\mathrm{\Lambda }=3/R^2`$ is given by
$$\mathrm{d}s^2=h\mathrm{d}t^2+h^1\mathrm{d}r^2+r^2\mathrm{d}\mathrm{\Omega }^2,A=Q/r\mathrm{d}t,$$
(1)
with
$$h=1\frac{r_+}{r}\frac{r_+^3}{R^2r}\frac{Q^2}{r_+r}+\frac{Q^2}{r^2}+\frac{r^2}{R^2}.$$
(2)
The asymptotic form of this spacetime is AdS. There is an outer horizon located at $`r=r_+`$. The mass of the black hole is
$$M=\frac{1}{2}(r_++\frac{r_+^3}{R^2}+\frac{Q^2}{r_+}).$$
(3)
The Hawking temperature is given by the expression
$$T_H=\frac{1{\displaystyle \frac{Q^2}{r_+^2}}+{\displaystyle \frac{3r_+^2}{R^2}}}{4\pi r_+}$$
(4)
and the potential by
$$\varphi =\frac{Q}{r_+}$$
(5)
In the extreme case $`r_+,Q`$ satisfy the relation
$$1\frac{Q^2}{r_+^2}+\frac{3r_+^2}{R^2}=0.$$
(6)
In the following discussions, we will concentrate our attention on the large black hole with $`r_+R`$. We will not consider the case of the small black hole here, partly because it is unstable, having negative specific heat and partly because it is not of direct interest for the AdS/CFT correspondence .
Let us consider a massless scalar field $`\mathrm{\Phi }`$ in the RN AdS spacetime, obeying the wave equation
$$\mathrm{}\mathrm{\Phi }=0$$
(7)
where $`\mathrm{}=g^{\alpha \beta }_\alpha _\beta `$ is the d’Alembertian operator. If we decompose the scalar field according to
$$\mathrm{\Phi }=\underset{lm}{}\frac{1}{r}\psi _l(t,r)Y_{lm}(\theta ,\varphi )$$
(8)
then each wave function $`\psi _l(r)`$ satisfies the equation
$$\frac{^2\psi _l}{t^2}+\frac{^2\psi _l}{r^2}=\mathrm{}_l\psi _l$$
(9)
where
$`\mathrm{}_l`$ $`=`$ $`h[{\displaystyle \frac{l(l+1)}{r^2}}+{\displaystyle \frac{1}{r}}{\displaystyle \frac{dh}{dr}}]`$ (10)
$`=`$ $`h[{\displaystyle \frac{l(l+1)}{r^2}}+{\displaystyle \frac{r_++r_+^3/R^2+Q^2/r_+}{r^3}}{\displaystyle \frac{2Q^2}{r^4}}+{\displaystyle \frac{2}{R^2}}]`$
and $`r`$ here is the tortoise coordinate defined by $`r={\displaystyle \frac{dr}{h}}`$.
The potential $`\mathrm{}`$ has the same characteristic as that in Schwarzschild AdS black hole. It is positive and vanishes at the horizon, however it diverges at $`r=\mathrm{}`$, which requires that $`\mathrm{\Phi }`$ vanishes at infinity. This is the boundary condition to be satisfied by the wave equation for the scalar field in AdS space.
Quasinormal modes of AdS space are difined to be modes with only ingoing wave near the horizon. There only exists a discrete set of complex quasinormal frequencies. The quasinormal modes behave like $`e^{i\omega (t+r)}`$ near the horizon in RN AdS background. We thus introduce the ingoing Eddington coordinates by setting $`v=t+r`$. The metric (1) can be rewritten as
$$\mathrm{d}s^2=h\mathrm{d}v^2+2\mathrm{d}v\mathrm{d}r+r^2(\mathrm{d}\theta ^2+\mathrm{sin}^2\theta \mathrm{d}\varphi ^2)$$
(11)
We separare the scalar field in a product form as
$$\mathrm{\Phi }=\frac{1}{r}\psi (r)Y(\theta ,\varphi )e^{i\omega v}$$
(12)
The minimally-coupled scalar wave equation (7) may thereby be reduced to an ordinary, second order, linear differential equation with the radial terms yielding
$$h(r)\frac{d^2\psi (r)}{dr^2}+(h^{}(r)2i\omega )\frac{d\psi (r)}{dr}V(r)\psi (r)=0,$$
(13)
where the potential function is given by
$`V(r)`$ $`=`$ $`{\displaystyle \frac{h^{}(r)}{r}}+{\displaystyle \frac{l(l+1)}{r^2}}`$ (14)
$`=`$ $`{\displaystyle \frac{1}{r}}({\displaystyle \frac{r_+}{r^2}}+{\displaystyle \frac{r_+^3}{R^2r^2}}+{\displaystyle \frac{Q^2}{r_+r^2}}{\displaystyle \frac{2Q^2}{r^3}}+{\displaystyle \frac{2r}{R^2}})+{\displaystyle \frac{l(l+1)}{r^2}}.`$
Note that by setting $`Q^2=0,R=1`$, Eqs.(13,14) go back to the 4D Schwarzschild AdS black hole case addressed in . In the following discussion we adopt $`R=1`$.
To find the complex values of $`\omega `$ such that (13) has a solution with $`\psi `$ finite at the horizon $`r=r_+`$, and vanishing at infinity, we have to count on the numerical calculations. By using the numerical method suggested in to compute the quasinormal modes, we will expand the solution in power series about the horizon and impose the boundary condition that the solution vanishes at infinity. Adopting the new variable $`x=1/r`$, (13) can be reexpressed as
$$s(x)\frac{d^2}{dx^2}\psi (x)+\frac{t(x)}{xx_+}\frac{d}{dx}\psi (x)+\frac{u(x)}{(xx_+)^2}\psi (x)=0$$
(15)
where the coefficient functions are
$`s(x)`$ $`=`$ $`{\displaystyle \frac{r_0x^5x^4x^2Q^2x^6}{xx_+}}`$ (16)
$`t(x)`$ $`=`$ $`3r_0x^42x^34Q^2x^52i\omega x^2`$ (17)
$`u(x)`$ $`=`$ $`(xx_+)V(x)`$ (18)
and the parameter $`r_0={\displaystyle \frac{1+x_+^2+Q^2x_+^4}{x_+^3}}`$. As done in , we can expand $`s,t`$ and $`u`$ about the horizon $`x=x_+`$ in the form $`s(x)=s_n(xx_+)^n`$ etc. The first terms are $`s_0=2x_+^2\kappa ,t_0=2x_+^2(\kappa i\omega )`$ and $`u_0=0`$, where $`\kappa `$ is the surface gravity and has the form $`\kappa =(x_++3/x_+Q^2x_+^3)/2`$. The solution of (15) can then be expressed as a power series
$$\psi (x)=\underset{n=0}{\overset{\mathrm{}}{}}a_n(xx_+)^n$$
(19)
Substituting (19) into (15) and equating coefficients of $`(xx_+)^n`$ for each $`n`$, we have the recursion relations for $`a_n`$
$$a_n=\frac{1}{P_n}\underset{k=0}{\overset{n1}{}}[k(k1)s_{nk}+kt_{nk}+u_{nk}]a_k$$
(20)
where
$$P_n=n(n1)s_0+nt_0=2x_+^2n(n\kappa i\omega ).$$
(21)
This relation has the same form as that of the Schwarzschild AdS case, but with an additional parameter $`Q`$ herein.
In order to find the quasinormal modes for the AdS spacetime, we must select solutions which satisfy the boundary condition $`\psi =0`$ as $`r\mathrm{}(x0)`$. Thus we need to look for the zeros of Eq(19) at $`x=0`$ in complex $`\omega `$ planes. For a given $`l,x_+,Q`$ the algorithm to find these frequencies follows the simple steps: (i) truncate (19) at a number $`N`$ of terms, construct the polynormial equation of $`\omega `$ using (20,21) so that (19) reduces to $`_{n=0}^Na_n(\omega )(x_+)^n=0`$, (ii) find roots of interest of this function, (iii) increase $`N`$ until these roots become constant within the desired precision. Since the problem is to find numerical solutions of a polynomial equation, it becomes easy to use a built-in Mathematica function to locate zeros of $`_{n=0}^Na_n(\omega )(x_+)^n`$ directly. This procedure may reduce tedious trials and make the numerical calculations neat.
We decompose the quasinormal frequencies into real and imaginary parts in the form
$$\omega =\omega _ri\omega _i$$
(22)
which makes $`\omega _i`$ positive for all quasinormal frequencies. For large black holes $`(r_+R)`$, the relation of the values of the lowest quasinomal mode frequencies for $`l=0`$ and selected values of $`r_+`$ for different charge $`Q`$ are exhibited in fig.1 and fig.2. The dots represent the lowest modes. In fig.1, lines from the top to the bottom correspond to $`Q^2=10^7,10^6,10^5`$ etc. However the lines shown in fig.2 are corresponding to $`Q^2=10,10^2,\mathrm{..10}^7`$ from the top to the bottom, respectively. It is easy to see that with an additional parameter, the charge $`Q`$, neither the real nor the imaginary part of the frequency is a linear function of $`r_+`$ as found in Schwarzschild AdS case . The bigger the charge $`Q`$ is, the larger is the deviation from the linear relation we observe.
The imaginary and real parts of the quasinormal frequencies relate to the damping time scale $`(\tau _1=1/\omega _i)`$ and oscillation time scale $`(\tau _2=1/\omega _r)`$, respectively. From fig.1 we learn that as $`Q`$ increase, $`\omega _i`$ increases as well, which corresponds to the decrease of the damping time scale. According to the AdS/CFT correspondence, this means that for big $`Q`$, it is quicker for the quasinormal ringing to settle down to the thermal equilibrium. Fig.2 tells us that the bigger charge $`Q`$ leads to the smaller $`\omega _r`$, which means that the frequency of the oscillation becomes small as $`Q`$ increases. Therefore from fig.1 and fig.2 we can have a picture that if we perturbe a RN AdS black hole with high charge, the surrounding geometry will not “ring” as much and long as that of the black hole with small $`Q`$. It is easy for the perturbation on the highly charged AdS black hole background to return to thermal equilibrium. This is the new physics brought by the additional parameter $`Q`$ in RN AdS black hole.
¿From (4), we learn that the temperature of the large charged AdS black hole does not scale linearly with the event horizon as that in Schwarzschild AdS black hole case. The behavior of the temperature is shown in fig.3. The linear relation between $`T`$ and $`r_+`$ is broken as $`Q`$ increases. The relations between the real and the imaginary parts of the frequency and the temperature for large RN AdS black hole are displayed in fig.4,5. Again we see that in contrast to the results in , the relations are no longer linear when the additional parameter $`Q`$ is taken into account.
We have so far discussed only the lowest quasinormal modes with $`l=0`$. Increasing $`l`$, we obtain the surprising effect of increasing the damping time scale ($`\omega _i`$ decreases), and decreasing the oscillation time scale ($`\omega _r`$ increases) as in the case of Schwarzschild AdS black holes. As pointed out in , here we may also meet the problem of possible negative $`\omega _i`$ as it continuously decreases with $`l`$. However from fig.6, it looks that for the large black hole, the problem is not as serious as that shown in for small Schwarzschild AdS case. Despite the similar behavior, once again the additional parameter $`Q`$ led us to further new properties. As fig.6 shows, different values of $`Q`$ do not change the qualitative characteristic of decreasing $`\omega _i`$ with $`l`$. Their decreasing rates are the same and do not depend on the value $`Q`$. However it is clear that the higher is the charge of AdS black hole, the later we will confront the tough question of $`\omega _i`$ as $`l\mathrm{}`$. The dependence of $`\omega _r`$ on $`l`$ for different values of $`Q`$ are exhibited in fig.7. Lines from upper to the bottom are $`Q^2=10^3,10^4\mathrm{}10^6`$, respectively. As shown for the case of the imaginary part, the lines of different charges are parallel.
In our numerical calculations, we found that we need for a large number $`N`$ of terms in the partial sum to reduce the relative error in the computation of quasinormal frequencies as $`r_+`$ decreases and $`Q`$ increases. However when the charge $`Q`$ increases to nearly the extreme value satisfying (6), we cannot accurately determine the mode, no matter how large is the value of $`N`$ that we adopt. The numerical convergence problem can be attributed to the method we adopted to compute the quasinormal modes. In the expansion of the solution of the differential equation into power series, we have the convergence radius
$$L=\underset{n\mathrm{}}{lim}|\frac{\psi _n}{\psi _{n+1}}||\frac{3+x_+^2Q^2x_+^4}{93x_++4x_+^2x_+^3Q^2(5x_+^4x_+^5)}|$$
(23)
To expand the solution in a power series about the horizon, we need $`L>x_+`$ to ensure that the expansion is valid. In the range $`0<Q^2<{\displaystyle \frac{93x_++4x_+^2x_+^3}{5x_+^4x_+^5}}`$, $`L`$ increases from $`L_0=(3+x_+^2)/(93x_++4x_+^2x_+^3)`$ to infinity, where $`L_0>x_+`$ for a big black hole, which means that $`L>x_+`$ always hold in this range. However for $`Q^2>{\displaystyle \frac{93x_++4x_+^2x_+^3}{5x_+^4x_+^5}}`$, $`L`$ continuously decreases with increasing values of $`Q`$. Finally when $`Q^2{\displaystyle \frac{3+x_+^2}{x_+^4}}=3r_+^4+r_+^2`$, which is the exact extreme value, $`L0<x_+`$. Therefore although the numerical approach is very efficient for determining the quasinormal modes for Schwarzschild AdS black hole and lowly charged AdS black hole, it breaks down for the nearly extreme RN AdS black hole case.
It is important to notice that Eq(15) has the same form as the generalized spheroidal wave equations which also arises both in the quantum scattering theory of nonrelativistic electrons from polar molecules and ions and in the theory of radiation process involving black hole in asymptotically flat spacetime. Similar convergence problem also appeared and hindered the earlier study of the quasinormal modes for extreme charged black hole and extreme rotating black hole .
Besides problems related to the method, is there some deep physics to account for the similar no convergence problem in determining quasinormal modes for extreme black holes both in asymptotically flat and AdS spaces? We know that the frequencies and damping times of the quasinomal modes are entirely fixed by the black hole, and are independent of the initial perturbation. It has been shown that there is a second order phase transition in the extreme limit of black holes \[14-16\]. This result has been further promoted in a recent study for the charged AdS black holes . Because of the phase transition, the fluctuation of thermodynamic quantities become tremendous. It is hard to expect that we can obtain the fixed quansinormal frequencies to characteristize the thermalization timescale in the strong coupled CFT on the extreme RN AdS background. We speculate that the black hole phase transition maybe the candidate physical reason behind this problem.
In summary, we have computed the scalar quasinormal modes of large RN AdS black hole. These modes govern the late time decay of a minimally coupled scalar field, such as the dilaton. Compared to the Schwarzschild AdS black hole, we found that the additional parameter $`Q`$ in RN AdS black hole has brought in some new properties. We observed that these modes no longer linearly scale with the black hole temperature. By AdS/CFT correspondence, we can interpretate the decay of the quasinormal modes to the time scale approaching thermal equilibrium in CFT. We learnt that perturbations are associated with the black hole charge. The larger the charge of the RN AdS black hole, the sooner it returns to the thermal equilibrium. The dependence of the quasinormal frequencies on nonzero angular momentum $`l`$ has also been discussed. Due to the convergence problem, the method we adopted cannot be extended directly to study of the extreme black holes cases. Further refinement of the numerical method of solving the differential equation (13) is called for. Based upon some critical phenomena uncovered, we have given some speculation of the possible physical reason behind the problem for the extreme black holes.
ACKNOWLEDGMENT: This work was partially supported by Fundac$`\stackrel{~}{a}`$o de Amparo $`\stackrel{`}{a}`$ Pesquisa do Estado de S$`\stackrel{~}{a}`$o Paulo (FAPESP) and Conselho Nacional de Desenvolvimento Cientifico e Tecnologico (CNPq). B. Wang would like to acknowledge the support given by Shanghai Science and Technology Commission.
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# Standard Conjectures for the Arithmetic Grassmannian 𝐺(2,𝑁) and Racah Polynomials
## 0. Introduction
Let $`X`$ be an arithmetic variety, by which we mean a regular, projective and flat scheme over $`\text{Spec}`$, of absolute dimension $`d+1`$. Assume that $`\overline{M}=(M,)`$ is a hermitian line bundle on $`X`$ which is arithmetically ample, in the sense of \[Z\] and \[So, §5.2\]. For each $`p`$ the line bundle $`\overline{M}`$ defines an arithmetic Lefschetz operator
$$\begin{array}{cccc}\widehat{L}:& \widehat{CH}^p(X)_{}& & \widehat{CH}^{p+1}(X)_{}\\ & \alpha & & \alpha \widehat{c}_1(\overline{M}).\end{array}$$
Here $`\widehat{CH}^{}(X)_{}`$ is the real arithmetic Chow ring of \[GS\] and $`\widehat{c}_1(\overline{M})`$ is the arithmetic first Chern class of $`\overline{M}`$.
In this setting, Gillet and Soulé \[GS\] proposed arithmetic analogues of Grothendieck’s standard conjectures \[Gr\] on algebraic cycles. A more precise version of the conjectures was formulated in \[So, §5.3\]; assuming $`2pd+1`$, the statement is
###### Conjecture 1.
(a) (Hard Lefschetz) The map
$$\widehat{L}^{d+12p}:\widehat{CH}^p(X)_{}\widehat{CH}^{d+1p}(X)_{}$$
is an isomorphism;
(b) (Hodge index) If the nonzero $`x\widehat{CH}^p(X)_{}`$ satisfies $`\widehat{L}^{d+22p}(x)=0`$, then
$$(1)^p\widehat{\mathrm{deg}}(x\widehat{L}^{d+12p}(x))>0.$$
We study these conjectures when $`X=G(r,N)`$ is the arithmetic Grassmannian, parametrizing $`r`$-dimensional subspaces of an $`(r+N)`$-dimensional vector space, over any field, and $`\overline{M}=\overline{𝒪}(1)`$ is the very ample line bundle giving the Plücker embedding, equipped with its natural hermitian metric. The latter is the metric induced from the standard metric on complex affine space, so that the first Chern form $`c_1(\overline{M})`$ on $`X()`$ is dual to the hyperplane class.
Our main result is that Conjecture 1 holds when $`r=2`$. For projective space ($`r=1`$) this was shown by Künnemann \[K\]. Moreover, it is proved in \[KM\] and \[Ta\] that Conjecture 1 holds for $`G(r,N)`$ after a suitable scaling of the metric on $`\overline{𝒪}(1)`$. To obtain the precise result for $`G(2,N)`$ we use the arithmetic Schubert calculus of \[T\] and linear algebra to reduce the problem to combinatorial estimates. In this case the inequality in part (b) asserts the positivity of a linear combination of harmonic numbers with coefficients certain Racah polynomials. The latter are a system of orthogonal polynomials in a discrete variable introduced by Wilson \[Wi\]\[AW\] which generalize the classical Racah coefficients or $`6`$-$`j`$ symbols \[Ra\] of quantum physics.
The results of Künnemann \[K\] show that each statement in Conjecture 1 (for given $`X`$, $`p`$ and $`\overline{M}`$) is true if and only if it holds when $`\widehat{CH}^p(X)_{}`$ is replaced by the Arakelov subgroup $`CH^p(\overline{X})_{}`$ associated to the Kähler form $`c_1(\overline{M})`$. We therefore restrict attention to this subgroup throughout the paper. In Section 1 we study arithmetic Lefschetz theory for varieties which admit a cellular decomposition and derive a cohomological criterion (Corollary 1) which we use to check Conjecture 1. This criterion does not suffice to check the Hodge index inequality on more general Grassmannians. In Section 2 we apply classical and arithmetic Schubert calculus to reduce the conjecture for $`G(2,N)`$ to estimates for a class of Racah polynomials. The required bounds for these polynomials are established in Section 3.
We thank Christophe Soulé for suggesting this problem to us and for general encouragement. Thanks are also due to Klaus Künnemann, Jennifer Morse and Herb Wilf for helpful discussions. It is a pleasure to acknowledge the support of National Science Foundation Postdoctoral Research Fellowships for both authors.
## 1. Arithmetic standard conjectures on cellular spaces
We study Conjecture 1 for arithmetic varieties $`X`$ which have a cellular decomposition over $`\text{Spec}`$, in the sense of \[F, Ex. 1.9.1\]; the Grassmannian $`G(r,N)`$ is a typical example. See \[KM\] for more information on these spaces and an approach to a weaker version of the conjecture. Recall that for each $`p`$ the class map
$$\mathrm{cl}:CH^p(X)_{}H^{p,p}(X_{})$$
is an isomorphism of the real Chow ring $`CH^p(X)_{}=CH^p(X)_{}`$ with the space $`H^{p,p}(X_{})`$ of real harmonic differential $`(p,p)`$-forms on $`X()`$. We denote by
$$\begin{array}{cccc}L:& CH^p(X)_{}& & CH^{p+1}(X)_{}\\ & \alpha & & \alpha c_1(M)\end{array}$$
the classical Lefschetz operator associated to an ample line bundle $`M`$ over $`X`$.
Let us equip the holomorphic line bundle $`M()`$ with a smooth hermitian metric, invariant under complex conjugation, to obtain a hermitian line bundle $`\overline{M}`$. As we have indicated, to check Conjecture 1 for the operator $`\widehat{L}(\alpha )=\alpha \widehat{c}_1(\overline{M})`$ it suffices to work with the Arakelov Chow group $`CH^p(\overline{X})_{}`$ defined using the Kähler form $`c_1(\overline{M})`$. Since $`X`$ has a cellular decomposition, we have an exact sequence
(1)
$$0CH^{p1}(X)_{}\stackrel{\stackrel{~}{a}}{}CH^p(\overline{X})_{}\stackrel{\zeta }{}CH^p(X)_{}0$$
(see \[KM, Prop. 6\]). Here $`\stackrel{~}{a}=a\mathrm{cl}`$ is the composite of the class map with the natural inclusion $`a:H^{p1,p1}(X_{})CH^p(\overline{X})_{}`$ and $`\zeta `$ is the projection defined in \[GS, §1\]. We choose a splitting
$$s_p:CH^p(X)_{}CH^p(\overline{X})_{}$$
for the sequence (1) and thus arrive at a direct sum decomposition
(2)
$$CH^p(\overline{X})_{}CH^p(X)_{}CH^{p1}(X)_{}.$$
for every $`p`$.
Summing (1) over all $`p`$ produces a sequence
(3)
$$0CH^1(X)_{}\stackrel{\stackrel{~}{a}}{}CH^{}(\overline{X})_{}\stackrel{\zeta }{}CH^{}(X)_{}0$$
which is compatible with the actions of $`L`$ and $`\widehat{L}`$. The splitting $`s:=_ps_p`$ of (3) does not commute with $`\widehat{L}`$ in general. Rather, the image of $`\widehat{L}ssL`$ is contained in $`\text{Ker}(\zeta )`$, hence
(4)
$$\widehat{L}ssL=\stackrel{~}{a}U$$
for a uniquely defined degree-preserving linear operator $`U`$ on $`CH^{}(X)_{}`$.
We now give some conditions equivalent to the arithmetic hard Lefschetz theorem (Theorem 1). When checking these for $`G(2,N)`$, we obtain something stronger, which establishes the arithmetic Hodge index theorem as well; this is quantified in Theorem 2. Recall the classical Lefschetz decomposition on $`CH^m(X)_{}H^{2m}(X(),)`$:
$$CH^m(X)_{}=\underset{p0}{}L^{mp}CH_{prim}^p(X)_{},$$
where the group of primitive codimension $`p`$ classes is
$$CH_{prim}^p(X)_{}=\text{Ker}(L^{d+12p}:CH^p(X)_{}CH^{d+1p}(X)_{}).$$
For $`m=dp`$ this decomposition induces a projection map
$$\pi _p:CH^{dp}(X)_{}L^{d2p}CH_{prim}^p(X)_{}.$$
###### Theorem 1.
Let $`X`$ be an arithmetic variety of dimension $`d+1`$ which admits a cellular decomposition. Let $`\widehat{L}`$ be the arithmetic Lefschetz operator associated to an ample hermitian line bundle on $`X`$. Then the following statements are equivalent:
* $`\widehat{L}^{d+12p}:CH^p(\overline{X})_{}CH^{d+1p}(\overline{X})_{}`$ is an isomorphism for all $`p`$.
* There exists a linear map $`\widehat{\mathrm{\Lambda }}:CH^{}(\overline{X})_{}CH^1(\overline{X})_{}`$ such that for every $`p`$ and $`\alpha CH^p(\overline{X})_{}`$ we have $`[\widehat{\mathrm{\Lambda }},\widehat{L}]\alpha =(d+12p)\alpha `$.
* For some (equivalently, any) choice of splitting $`s`$ of (3), with $`U`$ as in (4),
$$\pi _p\underset{i=0}{\overset{d2p}{}}L^{d2pi}UL^i:CH_{prim}^p(X)_{}L^{d2p}CH_{prim}^p(X)_{}$$
is an isomorphism for all $`p`$.
###### Proof.
We show that (iii) implies (ii). To do this, we first prove that if $`s^{}`$ is another splitting of (3) as above, with associated linear operator $`U^{}`$ on $`CH^{}(X)_{}`$, then for any $`p`$,
(5)
$$\pi _p\underset{i=0}{\overset{d2p}{}}L^{d2pi}U^{}L^i(\alpha )=\pi _p\underset{i=0}{\overset{d2p}{}}L^{d2pi}UL^i(\alpha )$$
for all $`\alpha CH_{prim}^p(X)_{}`$. Indeed, from (4) we have
(6)
$$\widehat{L}^kssL^k=\stackrel{~}{a}\underset{i=0}{\overset{k1}{}}L^{k1i}UL^i$$
for all $`k`$. Taking $`k=d2p+1`$ and using the fact that $`\alpha `$ is primitive, (6) gives
(7)
$$\widehat{L}^{d2p+1}(s^{}(\alpha )s(\alpha ))=\stackrel{~}{a}\underset{i=0}{\overset{d2p}{}}L^{d2pi}(U^{}U)L^i(\alpha ).$$
But $`s^{}(\alpha )s(\alpha )`$ is the class of a pure harmonic form, so the left-hand side of (7) is the harmonic form associated to an element which is in the image of $`L^{d2p+1}`$, and hence is killed by $`\pi _p`$.
We now claim there exists a splitting $`s^{}`$ such that for any $`p`$ and $`\alpha CH_{prim}^p(X)_{}`$ we have
(8)
$$U^{}L^i\alpha =0\text{ for all }i<d2p\text{and}U^{}L^{d2p}\alpha L^{d2p}CH_{prim}^p(X)_{}.$$
Indeed, if we let $`D`$ be the linear transformation such that $`s^{}s=\stackrel{~}{a}D`$, then
$$U^{}=U+[L,D],$$
and it is an exercise to check that the space of transformations $`[L,D]`$ is equal to the set of operators $`V`$ on $`CH^{}(X)_{}`$ satisfying $`\pi _p_{i=0}^{d2p}L^{d2pi}VL^i(\alpha )=0`$ for all $`\alpha CH_{prim}^p(X)_{}`$ and every $`p`$.
Suppose (iii) holds and choose $`s^{}`$ satisfying (8). Let us choose a primitive basis for $`CH^{}(X)_{}`$. Applying $`s^{}`$, we get half of a basis for $`CH^{}(\overline{X})_{}`$. By (iii), we may apply $`(L^{d2p})^1\pi _pU^{}L^{d2p}`$ to the basis elements in $`CH_{prim}^p(X)_{}`$ for each $`p`$ to obtain another basis for $`CH^{}(X)_{}`$, which we view (via $`\stackrel{~}{a}`$) as the other half of our basis for $`CH^{}(\overline{X})_{}`$.
Let $`vCH_{prim}^p(X)_{}`$ be one of the basis elements, and let $`r=d2p`$. By our conditions on $`s^{}`$, a subset of our basis for $`CH^{}(\overline{X})_{}`$ consists of $`\widehat{v}:=s^{}(v)`$, the iterates $`\widehat{L}^i(\widehat{v})=s^{}(L^iv)`$ of $`\widehat{L}`$ applied to $`\widehat{v}`$, the primitive element $`w`$ satisfying $`L^r(w)=\pi _p(U^{}L^r(v))`$, and the iterates of $`\widehat{L}`$ applied to $`w`$:
(9)
$$\widehat{v},\widehat{L}\widehat{v},\mathrm{},\widehat{L}^r\widehat{v},w,Lw,\mathrm{},L^rw.$$
The action of $`\widehat{L}`$ is to send each element in (9) to the element on its right, except that $`\widehat{L}^r\widehat{v}`$ is sent to $`L^rw`$, and $`L^rw`$ to $`0`$. We construct $`\widehat{\mathrm{\Lambda }}`$ explicitly: define
$`\widehat{\mathrm{\Lambda }}(\widehat{L}^i\widehat{v})`$ $`=i(r+2i)\widehat{L}^{i1}\widehat{v},`$
$`\widehat{\mathrm{\Lambda }}(L^iw)`$ $`=(r+1)\widehat{L}^i\widehat{v}+i(ri)L^{i1}w.`$
Then $`\widehat{\mathrm{\Lambda }}`$ (defined this way for every basis element $`v`$) satisfies the condition of (ii).
Statement (i) follows from (ii) by standard representation theory of $`sl(2)`$, as in \[We, §V.3\]. To show that (i) implies (iii), note that by (5) the condition in (iii) is independent of choice of splitting. If $`\alpha \text{Ker}\left(\pi _p_{i=0}^{d2p}L^{d2pi}UL^i\right)`$ is a nonzero primitive element and if we take $`s^{}`$ to be a splitting which satisfies (8), then
$$\widehat{L}^{d+12p}(s^{}(\alpha ))=\left(\widehat{L}^{d+12p}s^{}s^{}L^{d+12p}\right)\alpha =\underset{i=0}{\overset{d2p}{}}L^{d2pi}U^{}L^i(\alpha )=0,$$
and (i) fails. $`\mathrm{}`$
###### Theorem 2.
Suppose the arithmetic variety $`X`$ and $`p`$ are such that $`CH_{prim}^{p1}(X)_{}=0`$. If, for each nonzero $`\alpha CH_{prim}^p(X)_{}`$, we have
(10)
$$(1)^p\underset{i=0}{\overset{d2p}{}}_XL^{d2pi}\alpha UL^i\alpha >0$$
then the statements in the arithmetic hard Lefschetz and Hodge index conjectures are true for that $`X`$, $`p`$ and $`\overline{M}`$.
###### Proof.
Let $`(\alpha ,\beta )CH^p(\overline{X})_{}`$ be a nonzero element of the kernel of $`\widehat{L}^{d+22p}`$; the notation $`(\alpha ,\beta )`$ refers to the direct sum decomposition (2), with respect to some splitting. We claim that $`\alpha `$ must be in $`CH_{prim}^p(X)_{}`$. Indeed, $`\widehat{L}^{d+22p}(\alpha ,\beta )=(L^{d+22p}\alpha ,\gamma )`$ for some $`\gamma `$ and $`L^{d+22p}\alpha =0`$ implies $`L^{d+12p}\alpha =0`$ since $`CH_{prim}^{p1}(X)_{}`$ vanishes. Also, by the classical hard Lefschetz theorem, $`\alpha 0`$. Now, if
$$,:CH^{}(\overline{X})_{}CH^{}(\overline{X})_{}$$
denotes the arithmetic intersection pairing, then we have
$`(\alpha ,\beta ),\widehat{L}^{d+12p}(\alpha ,\beta )`$ $`=(\alpha ,\beta ),(0,{\displaystyle \underset{i}{}}L^{d2pi}UL^i\alpha +L^{d+12p}\beta )`$
$`={\displaystyle \frac{1}{2}}{\displaystyle \underset{i}{}}{\displaystyle _X}L^{d2pi}\alpha UL^i\alpha .`$
Hence, assuming $`CH_{prim}^{p1}(X)_{}=0`$, we have $`\widehat{L}^{d+12p}(\alpha ,\beta )0`$ for every nonzero $`(\alpha ,\beta )CH^p(\overline{X})_{}`$. Moreover, if $`(\alpha ,\beta )`$ is primitive, then the pairing of $`(\alpha ,\beta )`$ with $`\widehat{L}^{d+12p}(\alpha ,\beta )`$ has the required sign. $`\mathrm{}`$
###### Corollary 1.
Suppose $`X`$ is such that, for every $`p`$,
(11)
$$CH_{prim}^p(X)_{}0\text{implies}CH_{prim}^{p1}(X)_{}=0.$$
If condition (10) holds for every $`p`$ and each nonzero $`\alpha CH_{prim}^p(X)_{}`$, then both the arithmetic hard Lefschetz and Hodge index conjectures are true for $`X`$, $`\overline{M}`$.
Example. We illustrate both theorems for projective space $`^n`$ over $`\text{Spec}`$ (compare \[K, §4\]). In this case we choose the splitting
$$CH^{}(\overline{^n})_{}=\underset{i=0}{\overset{n}{}}\widehat{\omega }^i\underset{i=0}{\overset{n}{}}\omega ^i$$
where $`\widehat{\omega }^i=\widehat{c}_1(\overline{𝒪}(1))^i`$ and $`\omega ^i=\stackrel{~}{a}(c_1(𝒪(1))^i)`$. Then the sequence (9) is given by
$$\widehat{1},\widehat{\omega },\mathrm{},\widehat{\omega }^n,\tau _n,\tau _n\omega ,\mathrm{},\tau _n\omega ^n.$$
Here $`\tau _n=_{k=1}^n_k`$, where each $`_k=1+\frac{1}{2}+\mathrm{}+\frac{1}{k}`$ is a harmonic number. In the proof of Theorem 1 we constructed an explicit adjoint map $`\widehat{\mathrm{\Lambda }}`$ for the arithmetic Lefschetz operator $`\widehat{L}(x)=\widehat{\omega }x`$; in our example it is given by
$`\widehat{\mathrm{\Lambda }}(\widehat{\omega }^i)`$ $`=i(n+2i)\widehat{\omega }^{i1},`$
$`\widehat{\mathrm{\Lambda }}(\omega ^i)`$ $`={\displaystyle \frac{n+1}{\tau _n}}\widehat{\omega }^i+i(ni)\omega ^{i1}.`$
Observe that the nonzero primitive elements of $`CH^{}(^n)_{}`$ are multiples of $`1CH^0(^n)_{}`$; hence $`^n`$ satisfies (11). The operator $`U`$ is given by $`U(\omega ^i)=\delta _{i,n}\tau _n\omega ^n`$, and condition (10) for $`p=0`$, $`\alpha =1`$ becomes
$$\underset{i=0}{\overset{n}{}}_^n\omega ^{ni}U(\omega ^i)=\tau _n_^n\omega ^n=\tau _n>0.$$
The arithmetic Hodge index conjecture for $`^n`$ follows by applying Corollary 1.
## 2. The arithmetic Grassmannian $`G(2,N)`$
In this section we study Conjecture 1 for the Grassmannian of lines in projective space. For computational purposes we will work with the isomorphic Grassmannian $`G=G(N,2)`$ parametrizing $`N`$-planes in $`(N+2)`$-space throughout. Note that $`d=dim_{}G()=2N`$. There is a universal exact sequence of vector bundles
$$0SEQ0$$
over $`G`$; the complex points of $`E`$ and $`Q`$ are metrized by giving the trivial bundle $`E()`$ the trivial hermitian metric and the quotient bundle $`Q()`$ the induced metric. The hermitian vector bundles that result are denoted $`\overline{E}`$
The real vector space $`CH^{}(G)_{}H^2(G(),)`$ decomposes as
$$CH^{}(G)_{}=\underset{a,b}{}s_{a,b}(Q),$$
summed over all partitions $`\lambda =(a,b)`$ with $`aN`$, i.e., whose Young diagrams are contained in the $`2\times N`$ rectangle $`(N,N)`$. Moreover $`s_\lambda (Q)=s_{a,b}(Q)`$ is the characteristic class coming from the Schur polynomial $`s_{a,b}`$ in the Chern roots of $`Q`$; this is dual to the class of a codimension $`|\lambda |=a+b`$ Schubert variety in $`G`$.
The line bundle $`M=det(Q)`$ giving the Plücker embedding has $`c_1(M)=s_1(Q)`$; let $`L:CH^p(G)_{}CH^{p+1}(G)_{}`$ be the associated classical Lefschetz operator. Further for all $`p`$ let $`:CH^p(G)_{}CH^{2Np}(G)_{}`$ denote the Hodge star operator induced by the Kähler form $`s_1(\overline{Q})`$. We then have
###### Proposition 1.
The space $`CH_{prim}^p(G)_{}`$ is nonzero if and only if $`p=2kN`$. In the latter case it is one dimensional and spanned by the class
$$\alpha _k=\underset{j=0}{\overset{k}{}}(1)^j\left(\genfrac{}{}{0pt}{}{N+1j}{N2k}\right)\left(\genfrac{}{}{0pt}{}{N2k+j}{N2k}\right)s_{2kj,j}(Q).$$
Proof. By computing the Betti numbers for $`G`$ one sees that
$$dimCH^p(G)_{}dimCH^{p1}(G)_{}>0$$
if and only if $`p=2kN`$, and in this case the above difference equals 1. For such $`p`$ we have
(12)
$$\text{Ker}(L:CH^{2N2k}(G)_{}CH^{2N2k+1}(G)_{})=\mathrm{Span}\{\underset{j=0}{\overset{k}{}}(1)^js_{Nj,N2k+j}(Q)\}.$$
One checks (12) easily using the Pieri rule:
$$L(s_{a,b}(Q))=s_{a+1,b}(Q)+s_{a,b+1}(Q)$$
where it is understood that $`s_{c,c^{}}(Q)=0`$ if $`c<c^{}`$ or $`c>N`$.
From \[KT\] we know the action of the Hodge star operator on $`CH^{}(G)_{}`$ is given by
(13)
$$s_{a,b}(Q)=\frac{(a+1)!b!}{(Na)!(Nb+1)!}s_{Nb,Na}(Q).$$
Since
$$CH_{prim}^{2k}(G)_{}=\text{Ker}(L:CH^{2N2k}(G)_{}CH^{2N2k+1}(G)_{}),$$
the proof is completed by applying (13) to (12) and noting that the result is proportional to $`\alpha _k`$. $`\mathrm{}`$
We now pass to the arithmetic setting, where we use the arithmetic Schubert calculus of \[T, §3,4\]. The real Arakelov Chow group $`CH^p(\overline{G})_{}`$ decomposes as
(14)
$$CH^p(\overline{G})_{}=\underset{a+b=p}{}\widehat{s}_{a,b}(\overline{Q})\underset{a^{}+b^{}=p1}{}s_{a^{},b^{}}(\overline{Q}).$$
Here the indexing sets satisfy $`Nab0`$, $`\widehat{s}_{a,b}(\overline{Q})`$ is an arithmetic characteristic class and we identify the harmonic differential form $`s_{a^{},b^{}}(\overline{Q})`$ with its image in $`CH^p(\overline{G})_{}`$. The decomposition (14) is induced by the splitting map $`s_{a,b}(Q)\widehat{s}_{a,b}(\overline{Q})`$ which agrees with the one used in \[T\].
The hermitian line bundle $`\overline{M}`$ has $`\widehat{c}_1(\overline{M})=\widehat{s}_1(\overline{Q})`$ and is arithmetically ample; this follows from \[BGS, Prop. 3.2.4\]. We now apply the arithmetic Pieri rule of \[T, §4\] to compute the action of the arithmetic Lefschetz operator $`\widehat{L}(x)=\widehat{s}_1(\overline{Q})x`$ on the above basis elements. The induced operator $`U:CH^{}(G)_{}CH^{}(G)_{}`$ of (4) satisfies $`U(s_{a,b})=0`$ for $`a<N`$ and
(15)
$$U(s_{N,b})=\left(\underset{i=0}{\overset{N+1}{}}_i\right)s_{N,b}\underset{i=0}{\overset{(Nb)/2}{}}(_{Nb+1i}_i)s_{Ni,b+i}.$$
Here and in the rest of this section $`s_{a,b}`$ will denote the Schubert class $`s_{a,b}(Q)CH^{a+b}(G)_{}`$ and $`_i`$ is a harmonic number; by convention $`_0=0`$. Recall that the classical intersection pairing on $`CH(G)_{}`$ satisfies
$$s_{a,b},s_{a^{},b^{}}=_Gs_{a,b}s_{a^{},b^{}}=\delta _{(a,b),(Nb^{},Na^{})}.$$
The sequence of Betti numbers for $`G`$ shows that $`G`$ satisfies condition (11) of Corollary 1. We proceed to check the inequality (10) for all even $`p=2k`$; this will establish Conjecture 1 for $`G(N,2)`$. In our case (10) may be written as
$$\mathrm{\Sigma }(N,k):=_G\underset{b=0}{\overset{N2k}{}}L^{N2kb}\alpha _kUL^{N2k+b}\alpha _k>0.$$
To compute iterates of the classical Lefschetz operator $`L`$ on the Schubert basis, note that
(16)
$$L^rs_\mu =\underset{\genfrac{}{}{0pt}{}{\lambda \mu }{|\lambda |=|\mu |+r}}{}f^{\lambda /\mu }s_\lambda .$$
When $`\lambda =(\lambda _1,\lambda _2)`$ and $`\mu =(\mu _1,\mu _2)`$ are partitions with at most two parts the skew $`f`$-number in (16) satisfies
(17)
$$f^{\lambda /\mu }=\left(\genfrac{}{}{0pt}{}{|\lambda ||\mu |}{\lambda _1\mu _1}\right)\left(\genfrac{}{}{0pt}{}{|\lambda ||\mu |}{\lambda _1\mu _2+1}\right).$$
This follows from the determinantal formula for $`f^{\lambda /\mu }`$, given for example in \[St, Corollary 7.16.3\].
We now apply Proposition 1 and (16), (17) to calculate
$$L^{c2k}\alpha _k=\underset{j=0}{\overset{k}{}}(1)^j\left(\genfrac{}{}{0pt}{}{N+1j}{N2k}\right)\left(\genfrac{}{}{0pt}{}{N2k+j}{N2k}\right)L^{c2k}s_{2kj,j}$$
$$=\underset{j=0}{\overset{k}{}}\underset{ij}{}(1)^j\left(\genfrac{}{}{0pt}{}{N+1j}{2k+1j}\right)\left(\genfrac{}{}{0pt}{}{N2k+j}{j}\right)\left[\left(\genfrac{}{}{0pt}{}{c2k}{ij}\right)\left(\genfrac{}{}{0pt}{}{c2k}{i2k1+j}\right)\right]s_{ci,i}.$$
Therefore,
(18)
$$L^{c2k}\alpha _k=\underset{i,j}{}(1)^j\left(\genfrac{}{}{0pt}{}{N+1j}{N2k}\right)\left(\genfrac{}{}{0pt}{}{N2k+j}{N2k}\right)\left(\genfrac{}{}{0pt}{}{c2k}{ij}\right)s_{ci,i}.$$
We use (18) with $`c=N+b`$ to identify the coefficient of $`s_{N,b}`$ in the expansion of $`L^{N+b2k}\alpha _k`$ as
$`{\displaystyle _G}L^{N+b2k}\alpha _ks_{Nb}`$ $`={\displaystyle \underset{j}{}}(1)^j\left({\displaystyle \genfrac{}{}{0pt}{}{N+1j}{N2k}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{N2k+j}{j}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{N+b2k}{N2k+j}}\right)`$
$`=\left({\displaystyle \genfrac{}{}{0pt}{}{N2k+b}{N2k}}\right){\displaystyle \underset{j}{}}(1)^j\left({\displaystyle \genfrac{}{}{0pt}{}{b}{j}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{N+1j}{N2k}}\right)`$
$`=\left({\displaystyle \genfrac{}{}{0pt}{}{N+1b}{2k+1}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{N2k+b}{N2k}}\right)=:C_b.`$
It now follows from (15) that
(19)
$$UL^{N2k+b}\alpha _k=C_b\left(\underset{i=0}{\overset{N+1}{}}_i\right)s_{N,b}\underset{i=0}{\overset{(Nb)/2}{}}C_b(_{Nb+1i}_i)s_{Ni,b+i}.$$
Note also that we have the identity
(20)
$$\underset{b=0}{\overset{N2k}{}}C_b=\left(\genfrac{}{}{0pt}{}{2N2k+2}{N+2}\right).$$
Now we substitute $`c=Nb`$ in (18), pair with (19) and use (20) to sum over $`b`$ and obtain
(21)
$$\mathrm{\Sigma }(N,k)=A_{N,k}\underset{i=1}{\overset{N+1}{}}_i+\underset{i=0}{\overset{(Nb)/2}{}}C_{N,k}^i,$$
where
$$A_{N,k}=\left(\genfrac{}{}{0pt}{}{N+1}{N2k}\right)\left(\genfrac{}{}{0pt}{}{2N2k+2}{N+2}\right)$$
and
$$C_{N,k}^i=\underset{j,b}{}(1)^j\left(\genfrac{}{}{0pt}{}{N+1j}{N2k}\right)\left(\genfrac{}{}{0pt}{}{N2k+j}{N2k}\right)\left(\genfrac{}{}{0pt}{}{N2kb}{ij}\right)C_b(_i_{Nb+1i}).$$
By dividing the expression for $`C_{N,k}^i`$ into two sums and substituting in (21) one gets
(22)
$$\mathrm{\Sigma }(N,k)=A_{N,k}\underset{i=1}{\overset{N+1}{}}_i+\underset{i=1}{\overset{N+1}{}}B_{N,k}^i_i,$$
where
$$B_{N,k}^i=\underset{j,b}{}(1)^j\left(\genfrac{}{}{0pt}{}{N+1j}{N2k}\right)\left(\genfrac{}{}{0pt}{}{N2k+j}{N2k}\right)\left(\genfrac{}{}{0pt}{}{N2kb}{ij}\right)C_b.$$
(Notice that when $`Nb=2r1`$ is odd, there is a missing summand (for $`i=r`$)
$$\underset{j}{}(1)^j\left(\genfrac{}{}{0pt}{}{N+1j}{N2k}\right)\left(\genfrac{}{}{0pt}{}{N2k+j}{N2k}\right)\left(\genfrac{}{}{0pt}{}{2r2k1}{rj}\right)C_b_r$$
which vanishes, as can be seen by the change of variable $`j2k+1j`$.)
At this point it is convenient to introduce the variable change
$$n=N2k\text{and}T=N+2$$
and write equation (22) in the new coordinates as
(23)
$$\mathrm{\Sigma }(n,T)=A_{n,T}\underset{i=1}{\overset{T1}{}}_i+\underset{i=1}{\overset{T1}{}}B_{n,T}^i_i.$$
Observe that
$`B_{n,T}^i`$ $`={\displaystyle \underset{j}{}}(1)^j\left({\displaystyle \genfrac{}{}{0pt}{}{n+j}{n}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{T1j}{n}}\right){\displaystyle \underset{b}{}}\left({\displaystyle \genfrac{}{}{0pt}{}{nb}{ij}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{T1b}{nb}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{n+b}{n}}\right)`$
$`={\displaystyle \underset{j}{}}(1)^j\left({\displaystyle \genfrac{}{}{0pt}{}{n+j}{n}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{T1j}{n}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{T1n+ij}{ij}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{T+n}{ni+j}}\right).`$
We now substitute $`r=ij`$ and write the resulting sum in hypergeometric notation \[Ro\] \[VK, Chap. 3\]:
$`(1)^iB_{n,T}^i`$ $`={\displaystyle \underset{r}{}}(1)^r\left({\displaystyle \genfrac{}{}{0pt}{}{n+ir}{n}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{T1+ri}{n}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{T1n+r}{r}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{T+n}{nr}}\right)`$
$`=\left({\displaystyle \genfrac{}{}{0pt}{}{n+i}{n}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{T+n}{n}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{T1i}{n}}\right){}_{4}{}^{}F_{3}^{}\left(\begin{array}{c}n,i,Tn,Ti\\ ni,T+1,Tni\end{array}\right|1).`$
The Whipple transformation \[Wh, §10\] applied to the above $`{}_{4}{}^{}F_{3}^{}`$ gives
(24)
$$(1)^i\frac{B_{n,T}^i}{A_{n,T}}={}_{4}{}^{}F_{3}^{}\left(\begin{array}{c}n,n+1,i,i+1\\ 1,\mathrm{\hspace{0.17em}1}+T,\mathrm{\hspace{0.17em}1}T\end{array}\right|1).$$
The hypergeometric term (24) belongs to a class of orthogonal polynomials called Racah polynomials, which are studied in the next section.
## 3. Bounds for Racah polynomials
The Racah coefficients \[Ra\] or $`6`$-$`j`$ symbols have long been used by physicists as the transformation coefficients between two different coupling schemes of three angular momenta; see \[BL\] for an exposition. In mathematical language they are the entries of a change of basis matrix for the tensor product of three irreducible representations of $`SU(2)`$; the two bases involved come from the associativity relation for this product (see \[VK, §8.4\]). It was recognized later by Wilson \[Wi\] that these coefficients are special cases of a class of orthogonal polynomials $`R_n(x;\alpha ,\beta ,\gamma ,\delta )`$, called Racah polynomials \[VK, §8.5\]:
$$R_n(s(s+\gamma +\delta +1);\alpha ,\beta ,\gamma ,\delta )={}_{4}{}^{}F_{3}^{}\left(\begin{array}{c}n,n+\alpha +\beta +1,s,s+\gamma +\delta +1\\ \alpha +1,\beta +\delta +1,\gamma +1\end{array}\right|1).$$
The Racah polynomials in (24) have $`\alpha =\beta =\gamma +\delta =0`$, with $`\gamma =T`$, a positive integer. We let
$$R_n(s,T)=R_n(s(s+1);0,0,T,T).$$
Observe that $`R_n(s,T)`$ is symmetric in $`n`$ and $`s`$. The orthogonality condition (loc. cit. or \[AW\]) reads:
(25)
$$\underset{s=0}{\overset{T1}{}}(2s+1)R_n(s,T)R_m(s,T)=\frac{T^2}{(2n+1)}\delta _{nm}.$$
The arithmetic Hodge index inequality $`\mathrm{\Sigma }(n,T)>0`$ can be rephrased using (23) and (24) as
(26)
$$\underset{s=1}{\overset{T1}{}}(1)^{s+1}R_n(s,T)_s<\underset{s=1}{\overset{T1}{}}_s.$$
We give a proof of (26) which does not depend on the precise values of the harmonic numbers. Let us say that a sequence $`\{_k\}_{k1}`$ of positive real numbers (with $`_0=0`$) is concave increasing if $`_k=_{i=1}^kh_i`$ for some monotone decreasing sequence $`\{h_i\}`$ of positive reals.
###### Theorem 3.
Let $`\{_k\}`$ be any concave increasing sequence of real numbers and $`n`$, $`T`$ integers with $`0nT1`$ and $`T3`$. Then inequality (26) holds.
We believe that, in fact, (26) holds for an arbitrary sequence of positive real numbers $`_k`$, that is, the arithmetic standard conjectures for $`G(2,N)`$ do not depend on the relative sizes of the harmonic numbers involved:
###### Conjecture 2.
For any integers $`n,s`$ with $`0n,sT1`$ we have $`|R_n(s,T)|1`$.
In Proposition 2 we check this conjecture for some values of $`n`$ near the endpoints $`0`$ and $`T1`$. Computer calculations support the validity of Conjecture 2 for general $`n`$.
###### Proof of Theorem 3.
We shall see that (25) implies (26) except when $`T`$ is exponentially large compared to $`n`$. For large $`T`$, the Racah polynomials are close approximations of classical orthogonal polynomials, in this case the Legendre polynomials, and we know how to bound these.
By Cauchy’s inequality, (25) gives
$$\left(\underset{s=0}{\overset{T1}{}}|R_n(s,T)|_s\right)^2\frac{T^2}{2n+1}\underset{s=0}{\overset{T1}{}}\frac{_s^2}{2s+1}.$$
So, (26) holds whenever
(27)
$$\underset{s=0}{\overset{T1}{}}\frac{_s^2}{2s+1}<(2n+1)\left(\frac{1}{T}\underset{s=0}{\overset{T1}{}}_s\right)^2.$$
Since $`\{_k\}`$ is concave increasing, the average value of $`_0`$, $`\mathrm{}`$, $`_{T1}`$ is at least $`_{T1}/2`$. As $`_{s=1}^{T1}2/(2s+1)\mathrm{log}T`$, the inequality (27) holds whenever
(28)
$$\mathrm{log}T<n+\frac{1}{2}.$$
To analyze the case where $`T`$ is exponentially large compared to $`n`$, it is convenient to introduce the change of variable
(29)
$$x=s(s+1)=1/4+T^2(1+t)/2$$
and the rescaling
$$p_n(t)=(1)^n\underset{i=1}{\overset{n}{}}\frac{T^2i^2}{T^2}R_n(x;0,0,T,T).$$
Let $`P_n(t)=P_n^{(0,0)}(t)`$ denote the $`n^{\mathrm{th}}`$ Legendre polynomial. It is known \[NSU, §3.8\] that
$$p_n(t)=P_n(t)+O(1/T^2)$$
where the constant in the error term depends on both $`n`$ and $`t`$. For our purposes, we demonstrate
###### Lemma 1.
a) Let $`n`$ and $`T`$ be positive integers such that $`1+2n+2n^2<T^2/10`$. Then
(30)
$$|p_n(t)P_n(t)|(3/2)4^n/T^2$$
for all $`t`$ with $`1t1`$.
b) We have $`|p_n(t)P_n(t)|1/10`$ whenever $`T90`$ and $`n<\mathrm{log}T`$.
###### Proof.
We have the following recurrences (loc. cit.; for $`P_n(t)`$ this is classical)
(31) $`tP_n(t)`$ $`={\displaystyle \frac{n+1}{2n+1}}P_{n+1}(t)+{\displaystyle \frac{n}{2n+1}}P_{n1}(t)`$
(32) $`tp_n(t)`$ $`={\displaystyle \frac{n+1}{2n+1}}p_{n+1}(t){\displaystyle \frac{2n^2+2n+1}{2T^2}}p_n(t)+\left(1{\displaystyle \frac{n^2}{T^2}}\right)^2{\displaystyle \frac{n}{2n+1}}p_{n1}(t)`$
with initial data
(33)
$$\begin{array}{ccc}P_0(t)=1\hfill & ;\hfill & P_1(t)=t\hfill \\ p_0(t)=1\hfill & ;\hfill & p_1(t)=t+1/(2T^2).\hfill \end{array}$$
Subtracting (31) from (32) leads to a recurrence in $`p_n(t)P_n(t)`$. Then (30) follows by induction on $`n`$, using the known bound $`|P_n(t)|1`$ for all $`n`$ and all $`t`$ with $`1t1`$. The statement (b) is a corollary of (a). $`\mathrm{}`$
###### Proposition 2.
We have $`|R_n(s,T)|1`$ when $`n3`$ or $`n=T1`$.
###### Proof.
For $`n=0`$ we have $`R_0(s,T)=1`$. When $`n=1`$, we see from (33) that $`p_1(t)`$ is a increasing linear function in $`t`$, attaining minimum when $`s=0`$, giving $`R_1(0,T)=1`$, and maximum when $`s=T1`$, giving $`R_1(T1,T)=(1T)/(1+T)<1`$. For $`n=T1`$ the Pfaff-Saalschütz identity \[Ro\] \[VK, 8.3.3\] gives
$`R_{T1}(s,T)`$ $`={\displaystyle \underset{j}{}}(1)^j{\displaystyle \frac{T}{T+j}}\left({\displaystyle \genfrac{}{}{0pt}{}{s}{j}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{s+j}{j}}\right)`$
$`={}_{3}{}^{}F_{2}^{}\left(\begin{array}{ccc}s,s+1,T& & \\ 1,T+1& & \end{array}\right|1)`$
$`={\displaystyle \frac{(1T)(2T)\mathrm{}(sT)}{(1+T)(2+T)\mathrm{}(s+T)}}`$
so the inequality is clear.
When $`n=2`$ or $`n=3`$, $`p_n(t)`$ is a quadratic or cubic polynomial, and it is a calculus exercise to check that $`|R_n(s,T)|1`$ for every integer $`s`$ with $`0sT1`$. In fact, the integrality condition on $`s`$ is required only when $`n=3`$, $`T=4`$. $`\mathrm{}`$
###### Lemma 2.
a) We have $`|P_n(t)|3/4`$ for $`t`$, $`|t|0.9`$ and $`n2`$.
b) For $`T10`$, we have $`|t|0.9`$ in (29) whenever $`\sqrt{5}/10s/T4/5`$.
c) Assume $`T90`$ and $`n<\mathrm{log}T`$. Then $`{\displaystyle \frac{1}{T^{2n}}}{\displaystyle \underset{i=1}{\overset{n}{}}}(T^2i^2)>40/41`$.
###### Proof.
The indicated bound on Legendre polynomials is evident for $`n=2`$, and for larger $`n`$ it follows from the inequality
(34)
$$\sqrt{\mathrm{sin}\theta }|P_n(\mathrm{cos}\theta )|<\sqrt[]{\frac{2}{\pi n}},0\theta \pi .$$
One obtains (34) by using the transformed differential equation for $`\sqrt{\mathrm{sin}\theta }P_n(\mathrm{cos}\theta )`$ \[Sz, IV.4.24.2\]; this is indicated in \[H, Chap. 5, Exer. 15–16\]. The proofs of (b) and (c) are routine; the inequality $`\mathrm{log}(T^{2n}_{i=1}^n(T^2i^2))(2/T^2)_{i=1}^ni^2`$ may be used for the latter. $`\mathrm{}`$
To complete the proof of Theorem 3, assume that (28) fails, so that $`n<\mathrm{log}T1/2`$. If $`T90`$ then $`n3`$ and (26) follows from Proposition 2 (note that the inequality in the proposition is strict unless $`n=0`$ or $`s=0`$). When $`T90`$ we combine Lemma 1(b) with Lemma 2 to deduce the inequality (26). Indeed, $`(40/41)|R_n(s,T)|`$ is bounded by $`1+1/10`$ for every $`s`$, and by $`3/4+1/10`$ over the middle half of the summation range. By pairing terms $`_s`$ with $`_{T1s}`$ and using the fact that $`_s+_{T1s}`$ is monotone increasing for $`0s(T1)/2`$, we obtain (26). $`\mathrm{}`$
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# Hydrogen-Accreting Carbon-Oxygen White Dwarfs of Low Mass: Thermal and Chemical Behavior of Burning Shells
## 1 Introduction
The knowledge of the physical consequences of the accretion of hydrogen-rich matter onto a white dwarf has played an important role in understanding the main properties of several types of eruptive stars such as slow and fast novae and symbiotic stars (e.g., Starrfield 1971, Starrfield, Truran, & Sparks 1978, Sparks, Starrfield, & Truran 1978). In addition, the scenario in which a red giant star transfers mass to its carbon-oxygen (CO) white dwarf companion (Whelan & Iben 1973) is regarded by some as one of the most plausible precursor candidates for Type Ia supernovae (e.g., Hachisu, Kato, & Nomoto 1996). Fairly extensive surveys of the dependence of behavior on white dwarf mass and accretion rate have been conducted (see, e.g., Iben & Tutukov 1996; Cassisi, Iben, & Tornambé 1998 \[hereinafter CIT\] and references therein). For a low mass CO white dwarf of typical mass in the range $`0.5÷0.8`$ $`M_{\mathrm{}}`$, the consequences of the accretion of hydrogen-rich matter of solar metallicity can be summarized as follows:
* for sufficiently large mass-accretion rates (say, $`10^7`$ $`M_{\mathrm{}}`$ yr<sup>-1</sup> or larger, depending on the white dwarf mass), the accreted layer adopts an expanded configuration similar to that of the envelope of a red giant star;
* for intermediate mass-accretion rates (say, in the range $`4÷10\times 10^8`$ $`M_{\mathrm{}}`$ yr<sup>-1</sup>), the accretor burns hydrogen in a steady state regime at the same rate as it accretes hydrogen;
* for small mass-accretion rates (say, in the range $`1÷4\times 10^8`$ $`M_{\mathrm{}}`$ yr<sup>-1</sup>), the accretor experiences recurrent mild flashes;
* for even smaller mass-accretion rates (say, smaller than $`10^9`$ $`M_{\mathrm{}}`$ yr<sup>-1</sup>), the accretor experiences very strong nova-like hydrogen shell flashes.
Although considerable attention has been paid to models that experience recurrent mild hydrogen-burning flashes (e.g., Paczyński & $`\dot{Z}`$ytkow 1978, Iben 1982; José, Hernanz, & Isern 1993, CIT), a deep insight into the evolution of the main physical properties in the accreting models over an outburst-cooling cycle is still missing. Analytical studies of hydrogen-burning shells by Sugimoto & Fujimoto (1978) and by Fujimoto (1982a,b) have established the general properties of the hydrogen-burning shell as a function of the fundamental parameters of the accreting star, but have not provided profiles of structural and chemical variables in the shell itself. Finally, the extant numerical experiments do not explore systematically how the outcome of the accretion process depends on the abundances of heavy elements in the accreted matter.
To investigate in detail the evolution of the main physical characteristics of the hydrogen-burning shell during a flash episode, we have adopted as an initial model a white dwarf of mass 0.516 $`M_{\mathrm{}}`$ which has accreted matter at the rate $`\dot{M}=10^8`$ $`M_{\mathrm{}}`$ yr<sup>-1</sup> for $`7.6\times 10^5`$ yr, at which point the interior has been cooled to a temperature of $`8.6\times 10^6`$ K and the density at the center is $`2.56\times 10^6`$ g cm<sup>-3</sup>. For the composition of accreted matter, we have adopted a helium abundance by mass of $`Y=0.28`$ and three different metallicities ($`Z=0.02`$, 0.001, and 0.0001).
In §2 we discuss the input physics and the assumptions. In §3, the thermal properties of the hydrogen-burning shell are presented and discussed in detail for the Z=0.02 case. In §4, analytical relations in the ($`M_{\mathrm{WD}}`$-$`\dot{M}`$) plane are obtained for the case Z=0.02 and, in §5, the dependence on metallicity of the evolutionary behavior of the accreting models is analyzed. Conclusions and a brief discussion follow in §6.
## 2 Input Physics and Assumptions
The initial cold white dwarf model and the accretion experiments have been computed with an updated version of the FRANEC code (Chieffi & Straniero 1989). A detailed discussion of the main differences with respect to the original version of the code is given in CIT. The initial model is the same as that used in CIT and a description of the pre-accretion properties and the first pulse episode can be found in CIT. The main properties of this model are listed in Table 1.
The accretion process is computed on the assumption that the accreted matter and the white dwarf surface have the same specific entropy; that is, it is assumed that all of the energy liberated by matter as it falls onto the surface of the white dwarf is radiated away.
The input physics differs from that used in CIT only in the low temperature opacities: for $`Z=0.02`$ and $`Z=0.0001`$, we adopt the opacity tables provided by Cox & Stewart (1970a,b) and by Cox & Tabor (1976); for $`Z=0.001`$, the opacity values provided by Alexander & Fergusson (1994) have been used. The initial distribution of heavy elements adopted for models of metallicities Z=0.02 and Z=0.0001 is the solar one as given by Ross & Aller (1976); for Z=0.001, the initial distribution is taken from Grevesse (1991).
## 3 A Typical Hydrogen-Flash-Driven Pulse Cycle
In this section, we analyze the thermal and nuclear evolution of the hydrogen-burning shell of a model which has experienced enough (about 60) pulses that pulse properties have reached a local asymptotic limit. The model accretes hydrogen-rich matter of solar metallicity (Z=0.02) and its total mass at this point has grown to $`M_{\mathrm{WD}}0.5236M_{}`$. We assume that no mass is lost by the accreting star, despite the fact that a relatively nearby companion is required to supply the hydrogen-rich matter to the white dwarf and that, during evolution at high luminosity and low surface temperature, the surface of the star in outburst, in the real analogue, extends in most instances beyond the Roche lobe of the accretor and even beyond the donor star.
The most relevant structural properties of this model are reported in Table 1. The well known evolution in the HR diagram during one complete pulse cycle is reported in Figure 1. Several points of special interest are noted along the track: the positions where flash-driven convection begins (IC) and where it disappears (EC); the positions where maxima and minima of $`\mathrm{\Phi }_\mathrm{H}=L_\mathrm{H}/L_\mathrm{s}`$ and $`\mathrm{\Phi }_{\mathrm{He}}=L_{\mathrm{He}}/L_\mathrm{s}`$ occur. Here, $`L_\mathrm{H}`$, $`L_{\mathrm{He}}`$, and $`L_\mathrm{s}`$ are, respectively, the hydrogen-burning luminosity, the helium-burning luminosity, and the surface luminosity. Several interior characteristics during two passages of the cycle are shown in Figure 2 as a function of the total mass of the accretor.
We begin our description at the upper right hand portion of the evolutionary track as the model evolves from red to blue along the high luminosity plateau and $`\mathrm{\Phi }_\mathrm{H}0.994`$. In this phase, the release of gravothermal energy $`\mathrm{\Phi }_{\mathrm{gr}}=L_{\mathrm{gr}}/L_\mathrm{s}`$ is equal to $`0.006`$, whereas the Helium burning contributes negligibly to the surface luminosity ($`\mathrm{\Phi }_{\mathrm{He}}<<1`$). Figure 2a discloses that, along the plateau portion of the evolutionary track, the location in mass of the hydrogen-burning shell (which we define as the point where the maximum rate of energy production via hydrogen-burning is located) moves outwards much more quickly than the total mass grows due to the accretion of fresh matter. In the next section, it will be shown how, during the plateau phase, the relationship between the surface luminosity and the mass of the hydrogen-exhausted core depends on the metallicity of the accreted matter.
Once the mass of the layer between the hydrogen-burning shell and the surface decreases below a critical value, energy production by hydrogen burning declines rapidly (the reasons for this are discussed by Iben and the observational consequences are discussed by Iben & Tutukov ). The critical point in the HR diagram is the point of maximum effective temperature labeled BP (for “blue point”) in Figure 1.
After reaching the blue point, the hydrogen-burning efficiency plumets and gravothermal energy takes over as the main source of energy (Figure 2d). In the absence of mass accretion, further evolution would be similar to that of a single star after leaving the AGB to become the central star of a planetary nebula; i.e., apart from the onset of crystallization, evolution would not have presented any additional curiosities. However, continuous mass accretion leads to a quite different behavior for the outer layers of the white dwarf relative to that of a non-accreting, cooling white dwarf. In particular, the combined action of the accretion process (growth in mass of the hydrogen-rich envelope) and the contraction of the layers above the hydrogen-burning shell (see Figs. 2e and 2c) relatively quickly brakes the rate of decline of the temperature of the hydrogen-burning shell (Fig. 2b), and leads to a slowly, but continuously, growing hydrogen-burning luminosity (Fig. 2d). In addition, it is worth noticing that the accretion process induces an increase of the evolutionary lifetime in the blue side of the HR loop from the bluest point to the faintest one. In particular along this portion of the cycle the accreting model evolves in a 50% longer time-scale.
The energy delivered by hydrogen burning adds to the rate at which the local temperature in the shell increases until eventually a new hydrogen shell flash occurs. The increase in pressure related to the increase in temperature in the hydrogen-burning shell leads to a rapid expansion of the layers above the shell (see the surface radius increase in Fig. 2f and the decrease in density at the center of the shell in Fig. 2c). In the HR diagram, the model evolves upward as energy diffuses from the burning shell to the surface, and the increase in the surface radius eventually causes the model to evolve to the blue until the plateau portion of the track is reached once again.
Thus, the accretion process modifies the thermal content of the hydrogen-rich layer, producing the physical conditions suitable for a new hydrogen shell flash. An examination of the evolution of the temperature profile in the outer layers of the model and of the profile of the rate of nuclear energy generation in the hydrogen-rich layer shows how this takes place. In Figure 1, circles, triangles, and squares indicate the locations of all models for which profiles are given explicitly in Figures 3 and 4. All models marked by a given symbol in Figure 1 are represented by profiles plotted in a specific panel in Figures 3 and 4. Models marked by solid disks in Figure 1 are represented by the temperature profiles in Figure 3a and the energy-generation profiles in Figure 4a. Models marked by open triangles in Figure 1 are represented by profiles in Figures 3b and 4b. This continues in a counter clockwise fashion with, eventually, the models designated by open squares in Figure 1 being represented by the profiles in Figures 3f and 4f.
The temperature profile in Figure 3a which has the narrowest “peak” (and acts as a lower envelope of the ensemble of profiles) describes the model designated by the solid disk at smallest surface temperature in Figure 1. Thus, as the model evolves in the HR diagram along the plateau branch from red to blue, thermal energy in the model interior flows inward in the form of a “thermal wave,” heating up hydrogen-free matter. At the same time, thermal energy diffuses outward from the hydrogen-burning shell, which itself propagates outward in mass toward the model surface; the temperature profile between the shell and the surface steepens.
During the next portion of the evolution (the open triangles in Fig. 1 and the profiles in Fig. 3b), thermal energy stored over a fairly large fraction of the outer layers of the model (clearly beyond $`M_\mathrm{r}0.523`$ $`M_{\mathrm{}}`$ in Fig. 3b) leaks outward and the drop in temperatures in the hydrogen-burning shell is reflected in a decline in the hydrogen-burning luminosity. Matter interior to $`M_\mathrm{r}0.522`$ $`M_{\mathrm{}}`$ is still being heated from above. This behavior continues for the next designated set of models (solid boxes in Fig. 1 and profiles in Fig. 3c), with the “watershed” for thermal energy flow (the place from which energy flows both inward and outward) moving inward. For the following set of models (open circles in Fig. 1 and profiles in Fig. 3d), the energy-flow watershed has moved into interior regions not shown and, until the last few models, when hydrogen burning is beginning to make itself evident again (the bumps near the surface in the temperature profiles in Fig. 3d), cooling prevails over the entire region above $`M_\mathrm{r}0.521`$ $`M_{\mathrm{}}`$.
As the new hydrogen-burning shell flash gets underway (the last three open circles and the closed triangles in Fig. 1), a large fraction of the nuclear energy which is released is stored locally, since the time scale on which thermal energy can be transferred is greater than the time scale on which nuclear energy is released. This is true even though matter in the burning region is not degenerate. A convective shell is formed early on in the development of the flash (the point labeled IC in Fig. 1). The mass of the convective shell grows until the shell extends from the base of the burning shell up to photospheric layers. Up to this moment, the evolution of the model has occurred on the nuclear burning time scale. Thereafter, the convective shell (now, really, the convective envelope) cannot accommodate a further increase in its thermal energy content, but must expand to giant dimensions, using up local thermal energy to do the work of expansion against gravity. The readjustment to a new, expanded configuration occurs on the thermal time scale of the envelope. When the energy surplus has been dissipated, (excursion to higher luminosity and lower effective temperature) and the envelope has readjusted to an expanded configuration, the thermal wave begins to propagate inward again (see Fig. 3f). The model is now back to where it began (the profile in Fig 3a with the narrowest peak) .
The profiles of the nuclear energy-generation rate $`ϵ_\mathrm{n}`$ in Figure 4 are also instructive. As it is evident from Figure 4a, along the high luminosity plateau branch, the width in mass of the nuclear energy-generating region increases slightly and the hydrogen-burning luminosity drops as the burning shell works its way toward the surface. During most of the plateau phase, the full CNO cycles are active. By the time the model has attained its maximum effective temperature (Fig. 4b), hydrogen burning rapidly declines in importance as temperatures in the shell (Fig. 3b) decrease. The decline continues along the next portion of the cooling phase (Fig. 4c) as temperatures in the shell continue to decrease (Fig. 3c). An interesting aspect of nuclear burning during the long cooling phase (Figs. 4b, 4c, and 4d) is that it takes place in two distinct regions: a left hand region where the CNO cycles operate and a right hand spike where <sup>12</sup>C in freshly accreted fuel is burned into <sup>14</sup>N. Along this phase the pp contribution to energy production plays a not negligible role (see the broad secondary peak among the spikes in Fig. 4d).
Ultimately, thanks to the increase in density at the base of the hydrogen-rich layers, heating replaces cooling in the hydrogen-burning layers (Figs. 3d and 4d) and both the CNO cycles and pp-driven $`ϵ_\mathrm{n}`$ profiles begin to increase in height and in breadth. When convection appears, the second $`ϵ_\mathrm{n}`$ peak is “swallowed” by the first (Fig 4d). After the hydrogen-burning luminosity has attained its maximum value and convection begins to recede (Fig. 4e), the mass-width of the region undergoing the strongest hydrogen burning decreases, attaining its minimum width (Fig. 4f) at the same time the model reaches in the HR diagram the minimum effective temperature.
## 4 Steady-State Regime for Small White Dwarf Masses
The evolution described in the previous section illustrates the well known result that, for accretion rates smaller than a critical value which depends on white dwarf mass, accreting white dwarfs can be viewed as evolving alternatively in two stable states (a low state and a high state, in the nomenclature of Fujimoto \[1982b\]), separated by short lived transitional phases. The low state is the cooling phase, when the gravothermal energy source supplies essentially the entire surface luminosity, with hydrogen burning being almost extinguished. The high (or “excited”) state is the high luminosity plateau phase during which the hydrogen-burning shell is the main energy source, and the contribution of gravothermal energy is minor.
The hydrogen shell flash acts as the excitation mechanism which induces the transition from the low to the high state. The trigger for the transition is that, when the mass of the accreted layer exceeds a critical value which depends on the accretion rate, the rate of local heating by the hydrogen-burning shell exceeds the rate at which heat can diffuse out, initiating a thermonuclear runaway
The “strength” of a flash (the maximum hydrogen-burning luminosity) and the maximum extension to the red of the evolutionary track during the excited state depend on the accretion rate, in the sense of being greater, the smaller the accretion rate, but the final plateau luminosity during the evolution from red to blue depends only on the mass of the white dwarf. Since the amount of mass accreted between flashes is a quite small fraction of the total mass, the white dwarf mass may be thought of as constant over a large number of cycles; to a very good approximation, the structure of the envelope during the plateau phase does not “remember” the accretion rate preceding the flash which produced it, being sensitive only to the mass of the underlying white dwarf.
The transition from the excited to the low state sets in at the blue point along the evolutionary track in the HR diagram (point BP in Fig. 1). The sudden drop in $`\mathrm{\Phi }_\mathrm{H}`$ that is shown in Figure 2d is initiated at the blue point. The reason for this second transition can be understood in terms of the dependence on accretion rate of the mass of the hydrogen-rich layer $`\mathrm{\Delta }M_\mathrm{H}`$ in static models which are forced to burn hydrogen at the same rate as they accrete it (Iben 1982). The blue point in the HR diagram of the locus formed by a sequence of static models of fixed mass but different accretion rate is a bifurcation point (see Fig. 2 in Iben 1982), such that models of successively higher luminosity and lower surface temperature than at the blue point have larger $`\mathrm{\Delta }M_\mathrm{H}`$, whereas models with successively lower luminosity and surface temperature also have larger $`\mathrm{\Delta }M_\mathrm{H}`$. If one imagines turning off the accretion rate in any particular model on the upper branch and letting this model evolve statically, $`\mathrm{\Delta }M_\mathrm{H}`$ in that model would decrease because of nuclear burning and the model would evolve stably to the blue along the sequence, arriving at positions occupied by static steady state models of the same (successively smaller) $`\mathrm{\Delta }M_\mathrm{H}`$.
If, however, mass accretion were switched off in a model along the lower branch and this model were allowed to evolve statically, the model would evolve upward until it reached the bifurcation point, at which position it would be faced with a quandary. It could not evolve statically in either direction from the bifurcation point.
The implication of these thought experiments is that a real star evolving from red to blue along the plateau branch can do so in a roughly static fashion, with $`\mathrm{\Phi }_\mathrm{H}1`$ until, on reaching the bifurcation point, hydrogen burning can no longer control the course of evolution. The static approximation is no longer valid, $`\mathrm{\Phi }_\mathrm{H}`$ plumets until hydrogen burning is no longer of significance and gravothermal energy has taken over as the prime source of surface luminosity.
These considerations allow us to make use of the properties of accreting models in the quasistatic approximation to estimate, for small white dwarf masses, the lower boundary of the region in the $`M_{\mathrm{WD}}`$-$`\dot{M}`$ plane where steady state accretion can occur. In Figure 5, the rate $`ϵ_\mathrm{H}`$ of energy generation at the center of the hydrogen-burning shell (panel a) and the rate $`\dot{M}_\mathrm{H}`$ at which the hydrogen-burning shell processes matter (panel b) in a model of mass 0.5236 $`M_{\mathrm{}}`$ accreting $`Z=0.02`$ matter at the rate $`10^8`$ $`M_{\mathrm{}}`$ yr<sup>-1</sup> are shown. The position at the maximum surface temperature occurs along the curves in both panels corresponds to the bluest point in the HR diagram of Figure 1, and, there, $`\mathrm{log}\dot{M}_\mathrm{H}=7.554`$
In the $`M_{\mathrm{WD}}`$-$`\dot{M}_{\mathrm{H}\mathrm{sh}}`$ plane of Figure 6 are shown curves formed by four additional sets of models for larger white dwarf masses but for the same accretion rate of $`10^8`$ $`M_{\mathrm{}}`$ yr<sup>-1</sup>. The values of $`\dot{M}_\mathrm{H}`$ at the blue points for these model tracks and for several others are given in Table 2, where $`\mathrm{\Phi }_\mathrm{H}`$, $`\mathrm{log}T_\mathrm{e}`$ and $`\mathrm{log}(L/L_{\mathrm{}})`$ at the blue points are also given.
A linear fit to the properties of models at the bluest points in the HR diagram gives:
$$\mathrm{log}\dot{M}_{\mathrm{low}}(M_{\mathrm{}}\mathrm{yr}^1)=2.073\frac{M_{\mathrm{WD}}}{M_{\mathrm{}}}8.639,$$
(1)
as the lower boundary of the region in which steady state accretion solutions exist over the mass range $`0.52M_{\mathrm{WD}}/M_{\mathrm{}}0.68`$.
An approximation to the upper boundary of the region where steady state solutions exist can also be derived from the information in Figure 6. The “kink” in each curve near the red end of each curve for the three largest white dwarf masses in Figure 6 actually defines the point beyond which static solutions are of the red giant variety, with the accretion rate being larger than the rate at which nuclear burning consumes fuel (Fujimoto 1982 a,b). In Table 3, the values of $`\dot{M}_{\mathrm{H}\mathrm{sh}}`$ for several models is shown, along with the values of $`\mathrm{\Phi }_\mathrm{H}`$, $`\mathrm{log}T_\mathrm{e}`$ and $`\mathrm{log}(L/L_{\mathrm{}})`$ at the kink. A linear fit between the maximum allowed accretion rate and the WD mass at the kinks provides
$$\mathrm{log}\dot{M}_{\mathrm{high}}(M_{\mathrm{}}\mathrm{yr}^1)=1.512\frac{M_{\mathrm{WD}}}{M_{\mathrm{}}}7.800.$$
(2)
This line defines the upper boundary of the region in which steady state accretion can occur over the white dwarf mass range $`0.52M_{\mathrm{WD}}/M_{\mathrm{}}0.68`$.
The method adopted here to estimate the boundaries of the steady burning zone was introduced by Fujimoto (1982 a,b), who studied the properties of the hydrogen-burning shell in accreting models using an analytical solution for the envelope. He found (see Fig. 4 in Fujimoto 1982b) the border lines to be parallel for white dwarf masses over the range 0.5-1.5 $`M_{\mathrm{}}`$, while our border lines have different slopes. It is probable that our results differ because of the approximations used in the analytical solution. The fact that our lower border in the $`M_{\mathrm{WD}}`$-$`\dot{M}_\mathrm{H}`$ plane is steeper than the upper border is consistent with other estimates in the literature (see, e.g. Fig. 7 in Iben & Tutukov 1996 and Fig. 10 in CIT).
## 5 The Thermal Behavior of the Hydrogen-Burning Shell as a Function of Metallicity
To investigate the dependence on metallicity of the behavior of accreting white dwarfs, we have computed two additional sets of models in which hydrogen-rich matter characterized, respectively, by $`Z=0.001`$ and $`Z=0.0001`$ is accreted onto the same initial model of mass $`0.516M_{\mathrm{}}`$ at the same rate $`\dot{M}=10^8`$ $`M_{\mathrm{}}`$ yr<sup>-1</sup>. To make more meaningful the comparison between models accreting mass with different metallicities, we have adopted for all sets of models the same helium abundance: $`Y=0.28`$.
In the HR diagram of Figure 7 are shown the paths during one pulse cycle of models of the three different metallicities. In all three cases, the total mass of the model is $`M0.5236M_{}`$. Several characteristics of the models are given in Table 4. The dependence on the metallicity of various path characteristics can be understood relatively simply. At the very lowest luminosities, all paths converge because the models adopt the essentially metal-independent radius of a cold white dwarf. Because of smaller CNO abundances, the temperatures and densities at the base of the accreted layer (see Table 4) must be larger in models of lower metallicity in order for a CNO cycle thermonuclear runaway to be initiated: in order to achieve larger densities and temperatures, more mass must be accumulated by the lower metallicity models. This is why the time between pulses is larger, the lower the metallicity. During the transition between the low and high states, the radius of the expanding envelope is larger, the larger the mass of the envelope, and this explains why, at any given luminosity, the lower the metallicity, the redder the model. The fact that the reddest point along a path is bluer, the lower the metallicity, can be accounted for as an envelope-opacity effect. Finally, during the plateau phase and during the cooling phase, the fact that, at any luminosity, the model of lower metallicity is redder, is again ascribable to the larger mass of the hydrogen-rich envelope and the consequent larger radius of the envelope.
In Figure 8 we have reported for comparison the evolution of the main physical quantities for the hydrogen-burning shell for the cases with $`Z=0.001`$ and $`Z=0.0001`$. The evolution of $`\mathrm{\Phi }_\mathrm{H}`$ in panel (d) of this figure demonstrates graphically how the durations of both the high state (plateau phase) and the low state (cooling phase) increase with decreasing metallicity. As we have argued, the plateau phase lasts longer, the lower the metallicity, because the duration of the cooling phase and therefore the mass of hydrogen-rich material accreted between thermonuclear runaways increases with decreasing metallicity. That the amount of mass accreted during the low phase increases with decreasing metallicity is also evident by analyzing panel (a) which show that, the smaller the metallicity, the greater is the amount of fuel burned during the plateau phase. Figure 9 emphasizes this point once again.
It is evident from Figure 7, and the discussion in §4, that the band in the $`M_{\mathrm{WD}}\dot{M}`$ plane where steady state burning solutions exist drops to lower $`\dot{M}`$ as metallicity decreases. The semianalytical relations obtained in §4 suggest that a $`0.52`$ $`M_{\mathrm{}}`$ white dwarf accreting hydrogen-rich matter with $`Z=0.02`$ at $`8\times 10^8`$ $`M_{\mathrm{}}`$ yr<sup>-1</sup> settles into a steady state configuration, while, for an accretion rate of $`2\times 10^8`$ $`M_{\mathrm{}}`$ yr<sup>-1</sup>, it experiences recurrent mild flashes. To explore quickly the effect of the choice of metallicity on the location of the steady state band, we have calculated models of initial mass $`M_{\mathrm{WD}}=0.516M_{\mathrm{}}`$ and accretion rates $`2\times 10^8`$ and $`8\times 10^8`$ $`M_{\mathrm{}}`$ yr<sup>-1</sup>, for metallicities of $`Z=0.001`$ and $`Z=0.0001`$. The model accreting hydrogen at $`8\times 10^8`$ $`M_{\mathrm{}}`$ yr<sup>-1</sup> settles into a red giant configuration after only one pulse, while the model accreting at $`\dot{M}=2\times 10^8`$ $`M_{\mathrm{}}`$ yr<sup>-1</sup> settles into a steady state accretion configuration after one pulse.
Adopting the method outlined in §4, we have estimated the limits of the steady burning band for the three metallicities when $`M_{\mathrm{WD}}0.5236`$ $`M_{\mathrm{}}`$, obtaining the values listed in Table 5.
We have followed the long term evolution of the low $`Z`$ models, and from the results (which will be described in detail elsewhere), we have estimated the upper and lower bounds of the steady state band (see Tables 6 and 7, respectively).
The upper boundary may be approximated by
$$\mathrm{log}(\dot{M}_{\mathrm{high}}(M_{\mathrm{}}\mathrm{yr}^1))=2.235M_{\mathrm{WD}}/M_{\mathrm{}}8.350Z=0.0001$$
(3)
$$\mathrm{log}(\dot{M}_{\mathrm{high}}(M_{\mathrm{}}\mathrm{yr}^1))=1.832M_{\mathrm{WD}}/M_{\mathrm{}}8.043Z=0.001$$
(4)
and the lower boundary can be approximated by
$$\mathrm{log}(\dot{M}_{\mathrm{low}}(M_{\mathrm{}}\mathrm{yr}^1))=3.847M_{\mathrm{WD}}/M_{\mathrm{}}9.874Z=0.0001$$
(5)
$$\mathrm{log}(\dot{M}_{\mathrm{low}}(M_{\mathrm{}}\mathrm{yr}^1))=2.969M_{\mathrm{WD}}/M_{\mathrm{}}9.236Z=0.001$$
(6)
## 6 Summary and Conclusions
We have investigated and discussed in detail the evolutionary behaviour of a white dwarf accreting hydrogen-rich matter of three different metallicities: $`Z=0.02,0.001,\text{and }0.0001`$.
An analysis of the evolutionary behavior of several physical characteristics of the models has shown that, for fixed values of $`M_{WD}`$ and $`\dot{M}`$, lowering the metallicity causes the recurrence period to become longer because, in order to achieve the larger temperatures and densities necessary to offset the reduction of CNO catalysts in the accreted matter, the thickness of the hydrogen-rich accreted layer must increase.
For the steady-state burning regime, we have been able to derive borders in the $`M_{WD}`$-$`\dot{M}`$ plane as they depend on the metal content of the accreted matter. In agreement with earlier estimates, we find that the area of the region in this plane in which steady-state burning takes place becomes narrower as the white dwarf mass is increased. In addition, the location and the extension of the steady-state burning regime have been found to depend critically on the metallicity of the accreted matter, as shown clearly in Figure 10, where the topology of the steady-state region in the $`M_{WD}\dot{M}`$ plane is provided for the three metallicities considered. Reducing the metallicity, the steady-state burning region drops to smaller accretion rates and its extension is drastically decreased.
The consequences of our results for the final behavior of real low metallicity accretors are not easy to predict. On the one hand, as metallicity is decreased, the hydrogen-burning shell becomes hotter. This means that the underlying helium-burning layer is hotter and less degenerate when a helium-burning thermonuclear runaway is initiated. On the other hand, the fact that, for fixed core mass and accretion rate, the power of hydrogen-burning flash decreases as the metallicity is reduced suggests that low metallicity accretors may experience relatively mild helium shell flashes for a range of helium layer masses more extended than in the solar metallicity case.
However, as extensively discussed in Piersanti et al.(1999), in the mild pulse regime, there is a parameter region in which the effects of the hydrogen-burning shell on the helium layer are negligible, and a model with characteristics in theis region behaves as if pure helium is accreted. These models lead to a sub-Chandrasekhar explosion if the initial mass of the white dwarf and the accretion rate are within a given range (see Fig. 1 in Tornambé et al. 1998 for the solar metallicity case). Therefore, over the long term evolution, once the helium-burning layer becomes thermally decoupled from the hydrogen-burning shell, the accreted layer behaves in a way that depends only on the accretion rate and not on the metallicity.
Due to the prohibitively long computing time required, we have only partially studied the long term evolution of models accreting hydrogen-rich matter of metallicities $`Z=0.001`$ and $`Z=0.0001`$ at the mass-accretion rate of $`\dot{M}=10^8`$ $`M_{\mathrm{}}`$ yr<sup>-1</sup>. Both models show that the helium layer and the hydrogen-rich layer become decoupled as in the case of accretion of hydrogen-rich matter of solar metallicity (Piersanti et al. 1999). Such models will likely experience similar outcomes independent of the metallicity of the hydrogen-rich accreted matter.
On the basis of the results obtained so far, we suggest that, on lowering the metallicity, the area in the $`M_{WD}\dot{M}`$ plane suitable for sub-Chandrasekhar dynamical outcomes is shifted toward slightly lower values of $`\dot{M}`$, remaining almost unchanged in extention, as indicated in Tornambé et al. (1998) for the solar metallicity. It has to be finally considered that metallicity could even affect other parameters of the binary system (as, for instance, initial white dwarf masses, accretion rates, etc) with the consequence that in the real world, this scenario could be also significantly changed.
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# Integration of the SL(2,ℝ)/U(1) Gauged WZNW Theory by Reduction and Quantum Parafermions
## 1 Introduction
Gauged Wess-Zumino-Novikov-Witten (WZNW) models are an important subclass of two-dimensional integrable conformal field theories. But there is an essential structural difference between nilpotent and non-nilpotent gauging. Although the equations of motion for both cases have a linear Lax pair representation , only the Lax pairs for nilpotently gauged (Toda) theories have been integrated directly .
In this paper we will focus on non-nilpotently gauged WZNW theory. In contrast to the Toda case a systematic integration method is lacking. Instead of integrating a Lax pair, we extend the methods of Lagrangian and Hamiltonian reduction in a gauge invariant manner and prove that a non-nilpotently gauged WZNW theory can be integrated directly. We rederive here the general solution of the SL(2,$``$)/U(1) theory, which was found in refs in a non-systematic manner guided by the solution of the $`B_2`$ non-abelian Toda theory . While we will restrict ourselves to the simplest SL(2,$``$)/U(1) case, we presume that this approach can be generalised to any gauged WZNW theory.
Gauge invariant reduction also proves that the parafermionic conserved quantities are the coset currents, and the reduced energy-momentum tensor retains the simple Sugawara form in terms of these coset currents.
The SL(2,$``$)/U(1) theory is not only of mathematical interest. This model attracted much attention when it was realised that such theories have a black hole interpretation . In the early 1990’s studies of the quantum SL(2,$``$)/U(1) WZNW theory relied heavily on rather formal path integral manipulations or related operator identities . Since the attempted path integration over the $`U(1)`$ gauge field provided an incomplete effective action , we define the SL(2,$``$)/U(1) theory by the classical Lagrangian given in ref. . Our goal is to perform for this theory an exact canonical quantisation, much in the same way as it has been done for Liouville theory . Here one recasts the general solution as a canonical transformation exchanging interacting and free fields and lifts the classical conformal transformation to an operator transformation.
We take the parafermion algebra as a starting point for quantisation. In doing this we are forced to deform the classical free field representation of the parafermions. Similar deformations have been found before in OPE based Feigin-Fuks constructions of WZNW Kac-Moody currents . In addition we derive the quantum analogue of the parafermion algebra. The Sugawara form of the energy-momentum tensor motivates us to build the quantum energy-momentum tensor using only the parafermions, and we find an improvement term which in the $`\sigma `$-model picture could correspond to a non-perturbative dilaton. Finally, we point out that the general solution of the model has a parafermion interpretation. More precisely, the general solution contains fields of conformal weight zero related to an alternative set of ‘auxiliary’ parafermions which also undergo quantum deformation.
## 2 A Lagrangian Reformulation of the SL(2,$``$) <br>WZNW Theory
WZNW models are defined by the action
$$S_{\mathrm{WZNW}}[g]=\frac{k}{8\pi }\underset{M}{}h^{\mu \nu }\mathrm{tr}\left(g^1_\mu gg^1_\nu g\right)\sqrt{h}𝑑\sigma 𝑑\tau +kI_{\mathrm{WZ}}[g],$$
(1)
which includes the topological Wess-Zumino term
$$I_{\mathrm{WZ}}=\frac{1}{12\pi }\underset{B}{}\mathrm{tr}\left(g^1\text{d}gg^1\text{d}gg^1\text{d}g\right).$$
(2)
Here $`h_{\mu \nu }=diag(+,)`$ is the Minkowskian metric of the world surface $`M`$, $`h`$ its determinant, $`B`$ a volume with the boundary $`B=M`$, $`k`$ the coupling parameter, and $`g(\tau ,\sigma )`$ is a field which takes values in a semi-simple Lie group $`G`$. We shall restrict ourselves in this paper to the case $`G=`$ SL(2,$``$).
To simplify the construction of the SL(2,$``$)/U(1) theory it will prove useful to rewrite the topological WZNW term (2) as an integral of a local Lagrangian with global $`U(1)`$ symmetry. We are going to verify that the differentiation of the 2-form (for a more general treatment see )
$$F=\frac{2}{1+agag^1}L_aR_a$$
(3)
provides the integrand of the topological WZ term (2)
$$dF=\frac{2}{3}g^1dgg^1dgg^1dg,$$
(4)
and then Stokes’ theorem reduces the Wess-Zumino term to a two dimensional integral of $`F`$ over $`M=B`$. Here $`a`$ is a fixed normalised time-like element of the sl(2,$``$) algebra with $`aa=1`$, where $`=\frac{1}{2}\text{tr}()`$ denotes a normalised trace, and the left and right 1-forms of (3) are given by
$$L_a=adgg^1,R_a=ag^1dg.$$
(5)
Let us introduce the basis of the sl(2,$``$) algebra
$$t_0=\left(\begin{array}{cc}0& \hfill 1\\ 1& \hfill 0\end{array}\right),t_1=\left(\begin{array}{cc}0& \hfill 1\\ 1& \hfill 0\end{array}\right),t_2=\left(\begin{array}{cc}1& \hfill 0\\ 0& \hfill 1\end{array}\right).$$
(6)
It satisfies the relations
$$t_mt_n=\eta _{mn}I+ϵ_{mn}{}_{}{}^{l}t_{l}^{},$$
(7)
where $`I`$ is the unit matrix, $`\eta _{mn}=\text{diag}(+,,)`$ the $`3d`$ Minkowskian metric, and $`ϵ_{012}=1`$. The normalised traces of the matrices $`t_n`$ are then given by
$$t_mt_n=\eta _{mn},t_lt_mt_n=ϵ_{lmn}.$$
(8)
This defines an isometry between the sl(2,$``$) algebra and $`3d`$ Minkowski space.
The left and right 1-forms
$$L_n=t_ndgg^1,R_n=t_ng^1dg$$
(9)
are related by $`L_m=\mathrm{\Lambda }_m^n(g)R_n`$, where
$$\mathrm{\Lambda }_m^n(g)=t_mgt^ng^1$$
(10)
is a Lorentz transformation matrix. Since from (9) we have $`g^1dg=t^nR_n`$ and $`dgg^1=t^nL_n`$, using (8) the right hand side of (4) can be written in terms of right (or left) 1-forms
$$\frac{2}{3}g^1dgg^1dgg^1dg=4R_0R_1R_2.$$
(11)
Moreover, the differentials of (9) and (10)
$`dL_n`$ $`=`$ $`ϵ_n^{lm}L_lL_m,d\mathrm{\Lambda }_{mn}=2ϵ_n^{kl}\mathrm{\Lambda }_{mk}R_l`$
$`dR_n`$ $`=`$ $`ϵ_n^{lm}R_lR_m,`$ (12)
give for the left hand side of (4) the same result $`4R_0R_1R_2`$. This proves our statement (4). With Stokes’ theorem, we finally obtain the alternative Lagrangian formulation of the WZNW theory
$$S=_M𝑑z𝑑\overline{z},=_0+_{WZ},$$
(13)
where the kinetic term $`_0`$ remains unchanged
$$_0=\frac{1}{\gamma ^2}g^1_zgg^1_{\overline{z}}g,$$
(14)
but the WZ term becomes
$$_{WZ}=\frac{1}{\gamma ^2}\frac{a_zgg^1ag^1_{\overline{z}}ga_{\overline{z}}gg^1ag^1_zg}{1+agag^1}.$$
(15)
Here we used light-cone coordinates $`z=\sigma +\tau ,\overline{z}=\tau \sigma `$, and a new coupling constant $`\gamma ^2=2\pi /k`$. The Euler-Lagrange equations obtained from (13-15) reproduce consistently the dynamical equations of the WZNW theory (1)
$$_{\overline{z}}(_zgg^1)=0,_z(g^1_{\overline{z}}g)=0.$$
(16)
Note that the timelike property of $`a`$ ($`aa=1`$) guarantees regularity of (15).
## 3 The Gauged SL(2,$``$)/U(1) WZNW Theory
The Lagrangian (13-15) is invariant under the global transformations
$$gh_a(\epsilon )gh_a(\epsilon ),\text{with}h_a(\epsilon )=e^{\epsilon a},$$
(17)
which for timelike $`a`$ form the $`U(1)`$ subgroup of SL(2,$``$).
By a standard gauging procedure we introduce the $`U(1)`$ gauge fields $`A_z`$, $`A_{\overline{z}}`$ and get the new Lagrangian
$$_G(g,A_z,A_{\overline{z}},_zg,_{\overline{z}}g)=(g,_zgA_z(ag+ga),_{\overline{z}}gA_{\overline{z}}(ag+ga)),$$
(18)
which is invariant under the local transformations
$$A_zA_z+_z\epsilon ,A_{\overline{z}}A_{\overline{z}}+_{\overline{z}}\epsilon ,gh_a(\epsilon )gh_a(\epsilon ),\epsilon =\epsilon (z,\overline{z}).$$
(19)
The non-dynamical gauge fields can easily be eliminated from (18) through their algebraic equations of motion
$$A_z=\frac{a_zgg^1}{1+agag^1},A_{\overline{z}}=\frac{ag^1_{\overline{z}}g}{1+agag^1}.$$
(20)
So we obtain a gauge invariant Lagrangian only in terms of the field $`g`$
$`_G|`$ $`=`$ $`{\displaystyle \frac{1}{\gamma ^2}}(g^1_zgg^1_{\overline{z}}g`$ (21)
$`{\displaystyle \frac{a_zgg^1ag^1_{\overline{z}}g+a_{\overline{z}}gg^1ag^1_zg}{1+agag^1}}).`$
Without any loss of generality we assume $`a=t_0`$. Since $`t_0gt_0g^1=\mathrm{\Lambda }_{00}`$ is strictly positive the denominator in (21) is never zero. It is important for the following that this Lagrangian can be rewritten entirely in terms of the gauge invariant variables
$$v_1=t_1g,v_2=t_2g.$$
(22)
The gauge invariance follows from $`e^{ϵt_0}t_ne^{ϵt_0}=t_n\text{for}n=1,2`$. Introducing the additional variables $`v_0=t_0g`$ and $`c=g`$, we parameterise $`g(z,\overline{z})`$ SL(2,$``$) as
$$g=cI+v^nt_n=\left(\begin{array}{cc}cv_2& \hfill v_1v_0\\ v_1+v_0& \hfill c+v_2\end{array}\right),\text{with}c^2+v^nv_n=1.$$
(23)
Inserting this in (21), we find that the dependence on the gauge variant fields $`v_0`$ and $`c`$ is eliminated, and the reduced Lagrangian becomes
$$_G|=\frac{_zv_1_{\overline{z}}v_1+_zv_2_{\overline{z}}v_2}{\gamma ^2(1+v_1^2+v_2^2)}.$$
(24)
This Lagrangian has a natural complex structure in terms of the Kruskal coordinates $`u=v_1+iv_2`$, $`\overline{u}=v_1iv_2`$ of . The resulting equation of motion
$$_z_{\overline{z}}u=\overline{u}\frac{_zu_{\overline{z}}u}{1+u\overline{u}}$$
(25)
is just that obtained in .
## 4 Integration of the Theory by Lagrangian Reduction
The WZNW equations of motion (16) have the well-known general solution
$$g(z,\overline{z})=g_L(z)g_R(\overline{z}),$$
(26)
where $`g_L(z)`$, $`g_R(\overline{z})`$ are arbitrary SL(2,$``$) group valued (anti-)chiral functions.
Now let $`g(z,\overline{z})`$ be a solution (26) which satisfies the conditions
$$t_0_zg_L(z)g_L^1(z)=0\text{and}t_0g_R^1(\overline{z})_{\overline{z}}g_R(\overline{z})=0.$$
(27)
Then, due to (20), the set of functions $`g(z,\overline{z}),A_z(z,\overline{z})=A_{\overline{z}}(z,\overline{z})=0`$ form a solution of the dynamical equations derived from (18). Since the Lagrangians (18) and (24) have in terms of the gauge invariant fields $`v_1`$ and $`v_2`$ the same dynamical equations (25), the solutions of (25) can be written as
$$u(z,\overline{z})=(t_1+it_2)g_L(z)g_R(\overline{z}),$$
(28)
where the fields $`g_L`$ and $`g_R`$ satisfy (27). Equations (20) and (25) imply vanishing field strength $`F_{z\overline{z}}=_zA_{\overline{z}}_{\overline{z}}A_z=0`$. Then due to gauge invariance, (28) describes the general solution of (25).
We seek these solutions in terms of unconstrained (anti-) chiral fields. Therefore, we parameterise $`g_L`$ and $`g_R`$ as in (23)
$$g_L(z)=c(z)I+v^n(z)t_n,g_R(\overline{z})=\overline{c}(\overline{z})I+\overline{v}^n(\overline{z})t_n,$$
(29)
and introduce for convenience polar coordinates for the chiral fields
$`c`$ $`=`$ $`R\mathrm{cos}\beta ,v_0=R\mathrm{sin}\beta ,`$
$`v_1`$ $`=`$ $`r\mathrm{cos}\alpha ,v_2=r\mathrm{sin}\alpha ,`$ (30)
and similarly for the anti-chiral ones. The conditions (27) lead to $`R^2\beta ^{}r^2\alpha ^{}=0`$, and with $`R^2r^2=1`$ which follows from (23), we deduce the relations
$$R=\sqrt{\frac{\alpha ^{}}{\alpha ^{}\beta ^{}}},r=\sqrt{\frac{\beta ^{}}{\alpha ^{}\beta ^{}}}.$$
(31)
Here denotes differentiation. The insertion of (29) and (4) in (28) yields finally the general solution of (25)
$$u(z,\overline{z})=R(z)\overline{r}(\overline{z})e^{i\overline{\alpha }(\overline{z})i\beta (z)}+r(z)\overline{R}(\overline{z})e^{i\alpha (z)+i\overline{\beta }(\overline{z})},$$
(32)
which is correctly parameterised by the two chiral and two anti-chiral functions $`\alpha (z),\beta (z)`$, $`\overline{\alpha }(\overline{z}),\overline{\beta }(\overline{z})`$.
As a conformal field theory (24) has a traceless energy-momentum tensor $`T_{z\overline{z}}=0`$, with the chiral component
$$T=T_{zz}=\frac{1}{\gamma ^2}\frac{_z\overline{u}_zu}{1+u\overline{u}}=\frac{1}{\gamma ^2}\left(\alpha ^{}\beta ^{}+\frac{(\alpha ^{\prime \prime }\beta ^{}\beta ^{\prime \prime }\alpha ^{})^2}{4\alpha ^{}\beta ^{}(\alpha ^{}\beta ^{})^2}\right),$$
(33)
and similarly for the anti-chiral part $`\overline{T}=T_{\overline{z}\overline{z}}`$. A free-field form of this energy-momentum tensor, $`T(z)=\varphi _{1}^{}{}_{}{}^{2}(z)+\varphi _{2}^{}{}_{}{}^{2}(z)`$, can be obtained by passing to canonical free fields ($`k=1,2`$)
$$\psi _k(\sigma ,\tau )=\varphi _k(z)+\overline{\varphi }_k(\overline{z}),$$
(34)
where $`\varphi _k(z)`$ and $`\overline{\varphi }_k(\overline{z})`$ are chiral and anti-chiral components, respectively. The free-field transformation which solves this problem is given by
$`e^{i\alpha }`$ $`=`$ $`e^{i\gamma \varphi _2}{\displaystyle \frac{2e^{\gamma \varphi _1}e^{\gamma \varphi _1}+2ie^{\gamma \varphi _1}\mathrm{\Phi }}{\sqrt{\left(2e^{\gamma \varphi _1}e^{\gamma \varphi _1}\right)^2+4e^{2\gamma \varphi _1}\mathrm{\Phi }^2}}},`$ (35)
$`e^{i\beta }`$ $`=`$ $`e^{i\gamma \varphi _2}{\displaystyle \frac{2e^{\gamma \varphi _1}+e^{\gamma \varphi _1}2ie^{\gamma \varphi _1}\mathrm{\Phi }}{\sqrt{\left(2e^{\gamma \varphi _1}+e^{\gamma \varphi _1}\right)^2+4e^{2\gamma \varphi _1}\mathrm{\Phi }^2}}},`$
where $`\mathrm{\Phi }(z)`$ is defined as the integral of
$$_z\mathrm{\Phi }(z)=\gamma e^{2\gamma \varphi _1(z)}_z\varphi _2(z).$$
(36)
The general solution (32) then takes the form of a canonical transformation mapping SL(2,$``$)/U(1) fields onto free fields
$`u=e^{i\gamma (\varphi _2+\overline{\varphi }_2)}(e^{\gamma (\varphi _1+\overline{\varphi }_1)}(1+\mathrm{\Phi }\overline{\mathrm{\Phi }})`$ $``$ $`{\displaystyle \frac{1}{4}}e^{\gamma (\varphi _1+\overline{\varphi }_1)}`$ (37)
$`+`$ $`{\displaystyle \frac{i}{2}}(\mathrm{\Phi }e^{\gamma (\varphi _1\overline{\varphi }_1)}+\overline{\mathrm{\Phi }}e^{\gamma (\varphi _1\overline{\varphi }_1)})).`$
In the following we will impose periodicity in the spatial direction
$$u(\sigma +2\pi ,\tau )=u(\sigma ,\tau ),\psi _k(\sigma +2\pi ,\tau )=\psi _k(\sigma ,\tau ).$$
(38)
While the $`\psi _k`$ are strictly periodic their chiral and anti-chiral pieces have the following monodromy
$`\varphi _k(z+2\pi )=\varphi _k(z)+{\displaystyle \frac{p_k}{2}},\overline{\varphi }_k(\overline{z}2\pi )=\overline{\varphi }_k(\overline{z}){\displaystyle \frac{p_k}{2}},`$ (39)
where the $`p_k`$ are momentum zero modes. With these boundary conditions the integration of (36) gives
$$\mathrm{\Phi }(z)=\frac{\gamma }{2\mathrm{sinh}(\frac{1}{2}\gamma p_1)}_0^{2\pi }𝑑z^{}_z^{}\varphi _2(z^{})e^{\frac{1}{2}\gamma p_1ϵ_{2\pi }(zz^{})2\gamma \varphi _1(z^{})}.$$
(40)
Inserting (40) into (37) exactly reproduces the results obtained in .
Thus, we have demonstrated that the gauge invariant Lagrangian reduction indeed yields an integration method which allows one to solve the equations of motion of a non-nilpotently gauged WZNW theory in a straightforward manner.
## 5 Integration of the Theory by Hamiltonian Reduction
The Hamiltonian reduction of the WZNW theory is an alternative method for the construction and integration of coset models. The phase space of the system (13) is given by a set of functions $`R(\sigma ),g(\sigma )`$, where $`R(\sigma )`$ and $`g(\sigma )`$ take values in the sl(2,$``$) algebra and the SL(2,$``$) group, respectively. The reformulated WZNW action (13-15) can be written as $`S=\left(\theta H\text{d}\tau \right)`$ where the 1-form $`\theta `$ and the Hamiltonian $`H`$ are
$`\theta `$ $`=`$ $`{\displaystyle _0^{2\pi }}𝑑\sigma \left(Rg^1\text{d}g+{\displaystyle \frac{t_0g^1g^{}t_0\text{d}gg^1t_0g^{}g^1t_0g^1\text{d}g}{\gamma ^2(1+t_0gt_0g^1)}}\right)`$
$`H`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _0^{2\pi }}𝑑\sigma \left(\gamma ^2RR+{\displaystyle \frac{1}{\gamma ^2}}g^1g^{}g^1g^{}\right).`$ (41)
Here $`g^{}=_\sigma g`$, and d denotes the exterior derivative. Variation of $`R(\sigma )`$ yields the Hamiltonian equation
$$\gamma ^2R(\sigma )=g^1_\tau g.$$
(42)
Accordingly, we parameterise the functions $`R(\sigma ),g(\sigma )`$ by the SL(2,$``$) group valued fields $`g_L`$ and $`g_R`$
$`g(\sigma )`$ $`=`$ $`g_L(\sigma )g_R(\sigma ),`$
$`R(\sigma )`$ $`=`$ $`g_R^1(\sigma )g_L^1(\sigma )g_L^{}(\sigma )g_R(\sigma )+g_R^1(\sigma )g_R^{}(\sigma ).`$ (43)
Then the Hamiltonian in (5) splits into chiral and anti-chiral parts $`H=H_L+H_R`$, where
$$H_L=\frac{1}{\gamma ^2}_0^{2\pi }𝑑\sigma g_L^1g_L^{}g_L^1g_L^{},$$
(44)
and similarly for $`H_R`$. The corresponding splitting holds (up to boundary terms) also for the symplectic form $`\omega =\text{d}\theta `$. Using (5) and (5) we obtain
$`\omega `$ $`=`$ $`{\displaystyle \frac{1}{\gamma ^2}}{\displaystyle _0^{2\pi }}𝑑\sigma \left((g_L^1\text{d}g_L)^{}(g_L^1\text{d}g_L)(\text{d}g_Rg_R^1)^{}(\text{d}g_Rg_R^1)\right)`$ (45)
$`{\displaystyle \frac{1}{\gamma ^2}}\left((g_L^1\text{d}g_L)(\text{d}g_Rg_R^1){\displaystyle \frac{t_0g^1\text{d}gt_0\text{d}gg^1}{1+t_0gt_0g^1}}\right)|_0^{2\pi }.`$
The last term of this equation vanishes for periodic $`g`$ and in this case (45) reduces to the WZNW symplectic form of . Then the Hamiltonian equations split into $`\dot{g}_L=g_L^{},`$ and $`\dot{g}_R=g_R^{}`$, providing the general solution (26). That is why we used the same notation for the $`g_L`$, $`g_R`$ fields in the Hamiltonian and Lagrangian approaches.
The gauging procedure which led to the coset model (24) is equivalent to a Hamiltonian reduction with the same constraints (27). For the parameterisation (29), (4) the reduced chiral Hamiltonian and 2-form become
$$H_L|=\frac{1}{\gamma ^2}_0^{2\pi }d\sigma (f^2+\alpha ^{}\beta ^{}),\omega _L|=\frac{1}{\gamma ^2}_0^{2\pi }d\sigma (\text{d}f^{}\text{d}f+\text{d}\beta ^{}\text{d}\alpha ),$$
(46)
where $`\mathrm{tanh}^2f=\beta ^{}/\alpha ^{}`$, and we have neglected boundary terms in $`\omega _L`$. Note that the integrand of $`H_L`$ is just the energy-momentum tensor (33).
A canonical free-field form of (46) can be obtained by passing to the canonical free fields $`\varphi _1`$ and $`\varphi _2`$ using again the free-field transformation (35), and we get
$`T={\displaystyle \frac{1}{\gamma ^2}}(f^2+\alpha ^{}\beta ^{})`$ $`=`$ $`\varphi _1^2+\varphi _2^2,`$
$`{\displaystyle \frac{1}{\gamma ^2}}(\text{d}f^{}\text{d}f+\text{d}\beta ^{}\text{d}\alpha )`$ $`=`$ $`\text{d}\varphi _1^{}\text{d}\varphi _1+\text{d}\varphi _2^{}\text{d}\varphi _2`$ (47)
$`+\text{boundary terms}.`$
In fact the free-field transformation (35) was obtained as a solution of these equations.
This shows that the Hamiltonian reduction like the Lagrangian reduction provides a convenient approach for the integration of our non-nilpotently gauged SL(2,$``$)/U(1) WZNW theory. We envisage that these methods should be generalisable to other gauged WZNW theories.
## 6 The Parafermionic SL(2,$``$)/U(1) Coset Currents
It is well known that the chiral WZNW currents
$$\gamma ^2J_k(z)=t_k_zg(z,\overline{z})g^1(z,\overline{z})=t_k_zg_L(z)g_L^1(z)$$
(48)
satisfy the linear Kac-Moody algebra
$$\{J_k(z),J_l(z^{})\}=ϵ_{kl}^{}{}_{}{}^{m}J_m(z)\delta (zz^{})+\frac{1}{2\gamma ^2}\eta _{kl}\delta ^{}(zz^{}).$$
(49)
Here we would like to understand how these properties are impacted by the reduction. The chiral currents
$$\gamma ^2V_\pm (z)=(t_1\pm it_2)_zg_L(z)g_L^1(z)$$
(50)
are of particular interest. Taking into account the parameterisations (29,4), the constraints (31) and the free-field transformations (35), a straightforward calculation yields
$$V_\pm (z)=\frac{1}{\gamma }\left(_z\varphi _1(z)\pm i_z\varphi _2(z)\right)e^{\pm 2i\gamma \varphi _2(z)}.$$
(51)
Interestingly, these fields are the free-field transformed conserved parafermions of , which now obtain a coset current interpretation through the reduction.
The Sugawara energy-momentum tensor of the WZNW theory
$$T(z)=\gamma ^2J^k(z)J_k(z)$$
(52)
retains this simple form in terms of the parafermionic coset currents even after reduction
$$T(z)=\gamma ^2V_+(z)V_{}(z).$$
(53)
The free-field parameterisation of the parafermions (51) clearly leads to the free-field energy-momentum tensor (5). But in contrast to the WZNW Kac-Moody currents (48), the coset currents satisfy a non-linear and non-local Poisson bracket algebra, which can be calculated either through the Dirac bracket method or directly from the free field representations. The results for the Poisson brackets depend on the chosen boundary conditions. Here we work with the periodic ones prescribed in section 3. For the chiral and anti-chiral components $`\varphi _k(z)`$, $`\overline{\varphi }_k(\overline{z})`$ of the canonical free fields $`\psi _k(\sigma ,\tau )`$ (34) we choose the usual mode expansion
$`\varphi _k(z)`$ $`=`$ $`{\displaystyle \frac{1}{2}}q_k+{\displaystyle \frac{1}{4\pi }}p_kz+{\displaystyle \frac{i}{\sqrt{4\pi }}}{\displaystyle \underset{n0}{}}{\displaystyle \frac{a_n^{(k)}}{n}}\text{e}^{inz},`$
$`\overline{\varphi }_k(\overline{z})`$ $`=`$ $`{\displaystyle \frac{1}{2}}q_k+{\displaystyle \frac{1}{4\pi }}p_k\overline{z}+{\displaystyle \frac{i}{\sqrt{4\pi }}}{\displaystyle \underset{n0}{}}{\displaystyle \frac{\overline{a}_n^{(k)}}{n}}\text{e}^{in\overline{z}},`$ (54)
with the canonical mode algebra
$`\{q_k,p_l\}`$ $`=`$ $`\delta _{k,l},i\{a_n^{(k)},a_n^{(l)}\}=m\delta _{k,l}\delta _{m+n,0},`$
$`\{q_k,a_m^{(l)}\}`$ $`=`$ $`0,\{p_k,a_m^{(l)}\}=0.`$ (55)
Instead of the $`V_\pm (z)`$’s we will consider the periodic parafermions
$$W_\pm (z)=\frac{1}{\gamma }\left(_z\varphi _1(z)\pm i_z\varphi _2(z)\right)e^{\pm 2i\gamma \phi _2(z)},$$
(56)
where $`\phi _2`$ is $`\varphi _2`$ with the momentum zero mode $`p_2`$ removed and the whole $`q_2`$ zero mode of the free field (34) included, i.e. $`\phi _2(z)=\frac{1}{2}q_2+\varphi _2(z)|_{p_2=0}`$. These periodic coset currents have the algebra
$`\{W_\pm (z),W_\pm (z^{})\}`$ $`=`$ $`\gamma ^2W_\pm (z)W_\pm (z^{})h(zz^{}),`$
$`\{W_\pm (z),W_{}(z^{})\}`$ $`=`$ $`\gamma ^2W_\pm (z)W_{}(z^{})h(zz^{})`$
$`+{\displaystyle \frac{1}{\gamma ^2}}\left(_z+{\displaystyle \frac{i\gamma p_2}{2\pi }}\right)\delta _{2\pi }(zz^{}),`$
$`\{p_2,W_\pm (z^{})\}`$ $`=`$ $`2i\gamma W_\pm (z^{}),`$ (57)
where
$$h(z)=\left(ϵ_{2\pi }(z)\frac{z}{\pi }\right)$$
(58)
is the periodic sawtooth function and $`ϵ_{2\pi }(z)`$ the stairstep function<sup>1</sup><sup>1</sup>1$`ϵ_{2\pi }(\sigma )=2n+1`$ for $`2n\pi <\sigma <(2n+2)\pi `$ which coincides with $`\text{sign}(\sigma )`$ for $`2\pi <\sigma <2\pi `$.. Note that the momentum zero mode $`p_2`$ enters into the periodic parafermion algebra.
The energy-momentum tensor (53), now expressed in terms of the $`W_\pm `$’s, provides the Virasoro algebra
$$\{T(z),T(z^{})\}=_z^{}T(z^{})\delta _{2\pi }(zz^{})+2T(z^{})_z\delta _{2\pi }(zz^{}),$$
(59)
indicating for it conformal weight $`\mathrm{𝑡𝑤𝑜}`$, and the parafermions $`W_\pm (z)`$
$`\{T(z),W_\pm (z^{})\}`$ $`=`$ $`_z^{}W_\pm (z^{})\delta _{2\pi }(zz^{})+W_\pm (z^{})_z\delta _{2\pi }(zz^{})`$
$`{\displaystyle \frac{i\gamma p_2}{2\pi }}W_\pm (z^{})\delta _{2\pi }(zz^{})`$
have conformal weight $`\mathrm{𝑜𝑛𝑒}`$. Finally, we add a useful formula which generates the energy-momentum tensor $`T(z)`$ through a Poisson bracket
$`\{D_zW_+(z),W_{}(z^{})\}`$ $`=`$ $`\gamma ^2D_zW_+(z)W_{}(z^{})h(zz^{})`$
$`2T(z)\delta _{2\pi }(zz^{})+{\displaystyle \frac{1}{\gamma ^2}}D_z^2\delta _{2\pi }(zz^{}),`$
where
$$D_z=_z+\frac{i\gamma p_2}{2\pi }.$$
(62)
Equation (6) becomes important quantum mechanically because the operator product $`W_+(z)W_{}(z)`$ is ill defined and cannot be used to define a $`T(z)`$ operator.
## 7 Canonical Quantisation of the Parafermions
In this section we shall determine explicitly the quantum analogue of the parafermion algebra (6) instead of analysing operator product expansions. The quantisation of the theory will be defined by replacing Poisson brackets of the canonical free fields by commutators $`i\mathrm{}\{,\}[,]`$, and non-linear expressions in the free fields will be normal ordered. But calculations with normal ordered operators usually yield anomalous contributions. Such anomalies can be avoided by quantum mechanically deforming the composite operators of the theory. Therefore, let us define the normal ordered parafermion operators corresponding to (56) as
$$W_\pm (z)=\frac{1}{\gamma }:\left(\eta _z\varphi _1(z)\pm i_z\varphi _2(z)\right)e^{\pm 2i\gamma \phi _2(z)}:,$$
(63)
where $`\eta `$ is a deformation parameter with the classical limit $`\eta =1`$.
First we look for the quantum analogue of the Poisson brackets (6), starting with the simplest example $`\{W_+(z),W_+(z^{})\}=\gamma ^2W_+(z)W_+(z^{})h(zz^{})`$. In the appendix we have explicitly derived the relation
$`{\displaystyle \frac{W_+(z)W_+(z^{})}{e^{i\mathrm{}\gamma ^2h^+(zz^{})}}}`$ $``$ $`{\displaystyle \frac{W_+(z^{})W_+(z)}{e^{i\mathrm{}\gamma ^2h^{}(zz^{})}}}=`$
$`{\displaystyle \frac{i\mathrm{}}{2\gamma ^2}}\left(\eta ^21+{\displaystyle \frac{\gamma ^2\mathrm{}}{\pi }}\right):e^{2i\gamma \phi _2(z)}e^{2i\gamma \phi _2(z^{})}:_z\delta _{2\pi }(zz^{}),`$
where
$$h^\pm (z)=\frac{1}{2}h(z)\frac{i}{2\pi }\mathrm{log}\left(4\mathrm{sin}^2\frac{z}{2}\right)$$
(65)
are the positive and negative frequency parts of $`h(z)`$, respectively. The right hand side of (7) is proportional to the operator $`:e^{2i\gamma \phi _2(z)}e^{2i\gamma \phi _2(z^{})}:`$ which evidently cannot be rewritten bilocally in terms of the parafermions, as is necessary to have a closed operator algebra. However we can remove the offending term altogether by imposing the restriction
$$\eta ^21+\frac{\gamma ^2\mathrm{}}{\pi }=0,\text{or}\eta =\pm \sqrt{1\frac{\mathrm{}\gamma ^2}{\pi }}.$$
(66)
The classical limit $`\eta =1`$ corresponds to the positive square root. With this choice we have
$$\frac{W_+(z)W_+(z^{})}{e^{i\mathrm{}\gamma ^2h^+(zz^{})}}\frac{W_+(z^{})W_+(z)}{e^{i\mathrm{}\gamma ^2h^{}(zz^{})}}=0.$$
(67)
The quantum relations corresponding to the other Poisson brackets of (6) are (see the appendix)
$`{\displaystyle \frac{W_{}(z)W_{}(z^{})}{e^{i\mathrm{}\gamma ^2h^+(zz^{})}}}{\displaystyle \frac{W_{}(z^{})W_{}(z)}{e^{i\mathrm{}\gamma ^2h^{}(zz^{})}}}`$ $`=`$ $`0,`$ (68)
$`{\displaystyle \frac{W_+(z)W_{}(z^{})}{e^{i\mathrm{}\gamma ^2h^+(zz^{})}}}{\displaystyle \frac{W_{}(z^{})W_+(z)}{e^{i\mathrm{}\gamma ^2h^{}(zz^{})}}}`$ $`=`$ $`{\displaystyle \frac{i\mathrm{}}{\gamma ^2}}\left(_z+{\displaystyle \frac{i\gamma p_2}{2\pi }}\right)\delta _{2\pi }(zz^{}),`$ (69)
$`[p_2,W_\pm (z)]`$ $`=`$ $`\pm 2\mathrm{}\gamma W_\pm (z).`$ (70)
As in the derivation of (67) it is necessary to impose (66) to eliminate anomalous contributions.
To check that these operator relations correspond to the classical Poisson brackets we expand the exact formulae in powers of $`\mathrm{}`$, e.g. equation (67) gives
$$[W_+(z),W_+(z^{})]i\mathrm{}\gamma ^2W_+(z)W_+(z^{})h(zz^{})+O(\mathrm{}^2)=0.$$
(71)
Here we have used the splitting relation $`h(z)=h^+(z)+h^{}(z)`$, which follows immediately from (65).
## 8 The Energy-Momentum Tensor Operator
As was mentioned before, to generate the quantum energy-momentum tensor we should consider the quantum analogue of (6). From the calculations given at the end of the appendix it follows that
$`{\displaystyle \frac{D_zW_+(z)W_{}(z^{})}{e^{i\mathrm{}\gamma ^2h^+(zz^{})}}}`$ $``$ $`{\displaystyle \frac{W_{}(z^{})D_zW_+(z)}{e^{i\mathrm{}\gamma ^2h^{}(zz^{})}}}={\displaystyle \frac{i\mathrm{}}{\gamma ^2}}\left(1+{\displaystyle \frac{\mathrm{}\gamma ^2}{2\pi }}\right)D_{z}^{}{}_{}{}^{2}\delta _{2\pi }(zz^{})`$
$`2i\mathrm{}\eta ^2(:(_z\varphi _1)^2(z^{}):+:(_z\varphi _2)^2(z^{}):`$
$`+{\displaystyle \frac{\mathrm{}\gamma }{2\pi \eta }}_z^2\varphi _1(z^{})+{\displaystyle \frac{\mathrm{}\gamma ^2}{(4\pi \eta )^2}})\delta _{2\pi }(zz^{}).`$
The second entry on the right hand side just corresponds to the term $`2T(z^{})\delta _{2\pi }(zz^{})`$ of the classical Poisson bracket (6) which suggests that
$$T(z)=:(_z\varphi _1)^2(z):+:(_z\varphi _2)^2(z):+\frac{\mathrm{}\gamma }{2\pi \eta }_z^2\varphi _1(z)+\frac{\mathrm{}\gamma ^2}{(4\pi \eta )^2}$$
(73)
is the free-field energy-momentum tensor of our model with an additional improvement term. In the $`\sigma `$-model interpretation this improvement term could correspond to a non-perturbative dilaton .
Note that the energy-momentum tensor obeys the Virasoro algebra
$`[T(z),T(z^{})]`$ $`=`$ $`i\mathrm{}_z^{}T(z^{})\delta _{2\pi }(zz^{})+2i\mathrm{}T(z^{})_z\delta _{2\pi }(zz^{})`$
$`{\displaystyle \frac{i\mathrm{}c}{24\pi }}(_z^3+2_z)\delta _{2\pi }(zz^{}),`$
with central charge
$$c=\mathrm{}\left(2+\frac{3\gamma ^2}{\pi \eta ^2}\right),$$
(75)
in agreement with the results of .
Classically the parafermions are primary fields of weight one. Quantum mechanically the commutator
$`[T(z),W_+(z^{})]`$ $`=`$ $`i\mathrm{}\left(1+{\displaystyle \frac{\mathrm{}\gamma ^2}{2\pi }}\right)W_+(z^{})_z\delta _{2\pi }(zz^{})`$
$`i\mathrm{}_z^{}W_+(z^{})\delta _{2\pi }(zz^{})+{\displaystyle \frac{\mathrm{}\gamma }{2\pi }}:p_2W_+(z^{}):\delta _{2\pi }(zz^{})`$
shows that the quantum parafermions have the shifted conformal weight $`1+\mathrm{}\gamma ^2/(2\pi )`$.
## 9 Auxiliary Parafermions
In order to quantise the canonical transformation (37) we need the operators corresponding to the simple exponentials in (37) as well as the non-local fields $`\mathrm{\Phi }(z)`$ and $`\overline{\mathrm{\Phi }}(\overline{z})`$. However, it turns out that
$$A(z)=\frac{1}{2}e^{2\gamma \varphi _1(z)}i\mathrm{\Phi }(z),\overline{A}(\overline{z})=\frac{1}{2}e^{2\gamma \varphi _1(\overline{z})}i\mathrm{\Phi }(\overline{z})$$
(77)
are more amenable to a quantum treatment. Actually $`A(z)`$ is one of the chiral fields which parametrises the general solution given in . The derivative of $`A(z)`$ has a similar structure to that of the parafermions, and so we will refer to $`A^{}(z)`$ as an auxiliary parafermion
$$\stackrel{~}{V}_{}(z)=\gamma \left(_z\varphi _1(z)i_z\varphi _2(z)\right)e^{2\gamma \varphi _1(z)}.$$
(78)
It satisfies a closed chiral algebra if we replace $`\varphi _1(z)`$ by $`\frac{1}{2}q_1+\varphi _1(z)`$. The somewhat artificial doubling of the $`q_1`$ zero-mode is an artifact of our strict separation of chiral and anti-chiral objects, whereas each term in the general solution is a product of chiral and anti-chiral pieces. Following the recipe of section 7 the quantum auxiliary parafermion turns out to be
$$\stackrel{~}{V}_{}(z)=\gamma :\left(\eta _z\varphi _1(z)i_z\varphi _2(z)\right)e^{2\gamma \eta ^1(\frac{1}{2}q_1+\varphi _1(z))}:.$$
(79)
It commutes with $`V_{}(z^{})`$, but unlike $`V_{}(z)`$ retains its classical conformal weight one. Thus, $`A(z)`$ has conformal weight zero, suggesting that it plays the role of a screening charge.
Summarising, only one fixed deformation parameter is sufficient for consistent quantisation. Of course, our quantisation of the SL(2,$``$)/U(1) theory is still incomplete. In particular, it remains to quantise the black hole metric. It could answer the question whether a non-perturbative dilaton renders this metric dynamical.
Acknowledgments
G.W. thanks J. Schnittger for discussions on the quantum aspects of the problem. We would like to thank C. J. Biebl for reading the manuscript. G.J. is grateful to DESY Zeuthen for hospitality. His research was supported by grants from the DFG, GSRT and RFBR (99-01-00151).
## Appendix A Normal Ordered Operator Identities
In this appendix we elaborate on the normal ordered operator identities quoted in the text. To effect the normal ordering we decompose the chiral free fields $`\varphi _i(z)`$ defined in (6) as follows
$$\varphi _i(z)=\frac{1}{2}q_i+\frac{1}{4\pi }p_iz+\varphi _i^+(z)+\varphi _i^{}(z),$$
(80)
where
$$\varphi _i^\pm (z)=\pm \frac{i}{\sqrt{4\pi }}\underset{n>0}{}\frac{a_{\pm n}^{(i)}}{n}e^{inz}.$$
(81)
$`\varphi _i^{}(z)`$ and $`\varphi _i^+(z)`$ will be interpreted as creation and annihilation operators, respectively. The equivalent anti-chiral constructions will not be considered here.
Using the commutator algebra
$$[q_i,p_j]=i\mathrm{}\delta _{i,j},[a_m^{(i)},a_n^{(j)}]=m\mathrm{}\delta _{i,j}\delta _{m+n,0},i,j=1,2$$
(82)
it follows that
$$[\varphi _i^\pm (z),\varphi _j^\pm (z^{})]=0,[\varphi _i^\pm (z),\varphi _j^{}(z^{})]=\frac{i\mathrm{}}{4}\delta _{ij}h^\pm (zz^{}),$$
(83)
where
$$h^\pm (z)=ϵ^\pm (z)\frac{z}{2\pi }.$$
(84)
Here the $`ϵ^\pm (z)`$ denote the positive and negative frequency parts of the stairstep function $`ϵ_{2\pi }(z)`$, and have the Fourier series representation
$$ϵ^+(z)=\frac{z}{2\pi }+\frac{i}{\pi }\underset{n>0}{}\frac{e^{in(zi\epsilon )}}{n},ϵ^{}(z)=\frac{z}{2\pi }+\frac{i}{\pi }\underset{n<0}{}\frac{e^{in(z+i\epsilon )}}{n}.$$
(85)
Note that we have included a convergence factor, $`\epsilon >0`$. The $`ϵ^\pm (z)`$ functions have the properties $`ϵ^+{}_{}{}^{}(z)=ϵ^{}(z),ϵ^+(z)=ϵ^{}(z)`$, and in the limit $`\epsilon 0`$ they exhibit free-field short distance singularities
$$ϵ^\pm (z)=\frac{1}{2}ϵ_{2\pi }(z)\frac{i}{2\pi }\mathrm{log}\left(4\mathrm{sin}^2\frac{z}{2}\right).$$
(86)
We also introduce ‘split’ delta functions $`\delta ^+(z)=1/2_zϵ^+(z)`$
$$\delta ^+(z)=\frac{1}{4\pi }+\frac{1}{2\pi }\underset{n>0}{}e^{in(zi\epsilon )}=\frac{1}{4\pi }+\frac{1}{2\pi }\frac{1}{1e^{i(zi\epsilon )}},$$
(87)
and similarly $`\delta ^{}(z)=1/2_zϵ^{}(z)`$, which have the properties $`\delta ^+{}_{}{}^{}(z)=\delta ^{}(z)`$, $`\delta ^+(z)=\delta ^{}(z)`$, and as $`\epsilon 0`$ the $`\delta ^\pm (z)`$ sum up to the periodic delta function $`\delta _{2\pi }(z)`$.
Now we are ready to establish the normal ordered operator identities quoted in the text. As usual normal ordering moves creation and annihilation operators respectively to the left and right, and the Hermitian normal ordering of zero modes $`:e^{2q}f(p):=e^qf(p)e^q`$ will be understood. With the definitions
$$e_\pm (z):=e^{i\gamma q_2}e^{2i\gamma \varphi _2^\pm (z)},\nu (z):=\frac{1}{\gamma }\left(\eta _z\varphi _1(z)+i_z\varphi _2(z)\right),$$
(88)
our periodic parafermion operator $`W_+(z)`$ (63) can be written
$$W_+(z)=e_{}(z)\nu (z)e_+(z).$$
(89)
Let us start with a quantum analogue of the Poisson bracket $`\{W_+(z),W_+(z^{})\}`$. Naively one would consider the commutator, $`[W_+(z),W_+(z^{})]`$, suggesting that we should compute the operator product $`W_+(z)W_+(z^{})`$. Using the identity $`e^Ae^B=e^Be^Ae^{[A,B]}`$, which holds if $`[A,B]`$ commutes with $`A`$ and $`B`$, we have $`e_+(z)e_{}(z^{})=e_{}(z^{})e_+(z)e^{i\mathrm{}\gamma ^2h^+(zz^{})}`$, so that
$$e^{i\mathrm{}\gamma ^2h^+(zz^{})}W_+(z)W_+(z^{})=e_{}(z)\nu (z)e_{}(z^{})e_+(z)\nu (z^{})e_+(z^{}).$$
(90)
But this operator is still not correctly normal ordered. We decompose $`\nu (z)`$ as in (80)
$`\nu (z)`$ $`=`$ $`\nu ^+(z)+\nu ^{}(z)+{\displaystyle \frac{\eta p_1+ip_2}{4\pi \gamma }},`$
$`\nu ^\pm (z)`$ $`=`$ $`{\displaystyle \frac{1}{\gamma }}\left(\eta _z\varphi _1^\pm (z)+i_z\varphi _2^\pm (z)\right).`$ (91)
With a little algebra the right hand side of (90) can be rewritten as follows
$`e^{i\mathrm{}\gamma ^2h^+(zz^{})}W_+(z)W_+(z^{})`$ $`=`$ $`:W_+(z)W_+(z^{}):`$
$`+e_{}(z)e_{}(z^{})[\nu ^+(z),\nu ^{}(z^{})]e_+(z)e_+(z^{})`$
$`+e_{}(z)[\nu (z),e_{}(z^{})]\nu (z^{})e_+(z)e_+(z^{})`$
$`+e_{}(z)e_{}(z^{})\nu (z)[e_+(z),\nu (z^{})]e_+(z^{})`$
$`+e_{}(z)[\nu (z),e_{}(z^{})][e_+(z),\nu (z^{})]e_+(z^{}).`$
Using the results
$`[\nu (z),e_{}(z^{})]`$ $`=`$ $`i\mathrm{}e_{}(z^{})\delta ^\pm (zz^{}),`$
$`[\nu ^+(z),\nu ^{}(z^{})]`$ $`=`$ $`{\displaystyle \frac{i\mathrm{}}{2\gamma ^2}}(\eta ^21)_z\delta ^+(zz^{}),`$ (93)
we have
$`{\displaystyle \frac{W_+(z)W_+(z^{})}{e^{i\mathrm{}\gamma ^2h^+(zz^{})}}}`$ $`=`$ $`:W_+(z)W_+(z^{}):`$
$`+i\mathrm{}:\left(e^{2i\gamma \phi _2(z)}W_+(z^{})W_+(z)e^{2i\gamma \phi _2(z^{})}\right):\delta ^+(zz^{})`$
$`+:e^{2i\gamma \phi _2(z)}e^{2i\gamma \phi _2(z^{})}:({\displaystyle \frac{i\mathrm{}}{2\gamma ^2}}(\eta ^21)_z\delta ^+(zz^{})`$
$`+\mathrm{}^2\left(\delta ^+(zz^{})\right)^2).`$
At this point the following identity is useful
$$\left[\delta ^+(z)\right]^2=\frac{1}{(4\pi )^2}+\frac{i}{2\pi }_z\delta ^+(z),$$
(95)
which is valid even for finite $`\epsilon `$. Using this formula the right hand side of (90) can be written linearly in $`\delta ^+(zz^{})`$ and its derivative. Recall that this distribution becomes $`\delta ^{}(zz^{})`$ on exchanging $`z`$ and $`z^{}`$. Thus, if we take (A) and subtract the equation obtained by exchanging $`z`$ and $`z^{}`$, we get
$`{\displaystyle \frac{W_+(z)W_+(z^{})}{e^{i\mathrm{}\gamma ^2h^+(zz^{})}}}`$ $``$ $`{\displaystyle \frac{W_+(z^{})W_+(z)}{e^{i\mathrm{}\gamma ^2h^{}(zz^{})}}}=`$
$`i\mathrm{}:\left(e^{2i\gamma \phi _2(z)}W_+(z^{})W_+(z)e^{2i\gamma \phi _2(z^{})}\right):\delta _{2\pi }(zz^{})`$
$`+{\displaystyle \frac{i\mathrm{}}{2\gamma ^2}}\left(\eta ^21+{\displaystyle \frac{\mathrm{}\gamma ^2}{\pi }}\right):e^{2i\gamma \phi _2(z)}e^{2i\gamma \phi _2(z^{})}:_z\delta _{2\pi }(zz^{}).`$
The first term on the right hand side is zero since the prefactor of $`\delta _{2\pi }(zz^{})`$ tends to zero as $`zz^{}`$, and so (7) follows immediately.
Since all the other results quoted in the text can be derived by the same technique, we will be rather sketchy from now on. Evaluating the operator products $`W_+(z)W_{}(z^{})`$ and $`W_{}(z^{})W_+(z)`$ one is led to
$`{\displaystyle \frac{W_+(z)W_{}(z^{})}{e^{i\mathrm{}\gamma ^2h^+(zz^{})}}}`$ $``$ $`{\displaystyle \frac{W_{}(z^{})W_+(z)}{e^{i\mathrm{}\gamma ^2h^{}(zz^{})}}}=`$
$`i\mathrm{}:\left(\nu (z)\nu ^{}(z^{})\right)e^{2i\gamma \phi _2(z)}e^{2i\gamma \phi _2(z^{})}:\delta _{2\pi }(zz^{})`$
$`+{\displaystyle \frac{i\mathrm{}}{2\gamma ^2}}\left(\eta ^2+1+{\displaystyle \frac{\mathrm{}\gamma ^2}{\pi }}\right):e^{2i\gamma \phi _2(z)}e^{2i\gamma \phi _2(z^{})}:_z\delta _{2\pi }(zz^{}).`$
Using the identity $`f(z)\delta _{2\pi }(zz^{})=f(z^{})\delta _{2\pi }(zz^{})`$ (valid for periodic $`f(z)`$) and derivatives thereof
$`{\displaystyle \frac{W_+(z)W_{}(z^{})}{e^{i\mathrm{}\gamma ^2h^+(zz^{})}}}{\displaystyle \frac{W_{}(z^{})W_+(z)}{e^{i\mathrm{}\gamma ^2h^{}(zz^{})}}}`$ $`=`$ $`{\displaystyle \frac{\mathrm{}p_2}{2\pi \gamma }}\delta _{2\pi }(zz^{})`$
$`+{\displaystyle \frac{\mathrm{}}{\gamma }}\left(\eta ^21+{\displaystyle \frac{\mathrm{}\gamma ^2}{\pi }}\right)_z\phi _2(z^{})\delta _{2\pi }(zz^{})`$
$`+{\displaystyle \frac{i\mathrm{}}{2\gamma ^2}}\left(\eta ^2+1+{\displaystyle \frac{\mathrm{}\gamma ^2}{\pi }}\right)_z\delta _{2\pi }(zz^{}).`$
(69) follows on imposing (66).
We conclude by laying the ground for the derivation of the energy-momentum operator. To prove (8) one first establishes that
$`{\displaystyle \frac{D_zW_+(z)W_{}(z^{})}{e^{i\mathrm{}\gamma ^2h^+(zz^{})}}}{\displaystyle \frac{W_{}(z^{})D_zW_+(z)}{e^{i\mathrm{}\gamma ^2h^{}(zz^{})}}}`$ $`=`$ (99)
$`2i\mathrm{}\gamma ^2A(z,z^{})`$ $`+`$ $`{\displaystyle \frac{i\mathrm{}}{\gamma ^2}}D_z^2\delta _{2\pi }(zz^{}),`$
where
$$A(z,z^{})=\frac{W_+(z)W_{}(z^{})}{e^{i\mathrm{}\gamma ^2h^+(zz^{})}}\delta ^+(zz^{})+\frac{W_{}(z^{})W_+(z)}{e^{i\mathrm{}\gamma ^2h^{}(zz^{})}}\delta ^{}(zz^{}).$$
(100)
In computing $`A(z,z^{})`$ we encounter the same operator products as in the calculation of (A). The result can be written as
$`A(z,z^{})`$ $`=`$ $`(:W_+(z)W_{}(z^{}):+{\displaystyle \frac{\mathrm{}}{16\pi ^2}})\delta _{2\pi }(zz^{})`$
$`{\displaystyle \frac{\mathrm{}}{2\pi }}:e^{2i\gamma \phi _2(z)}e^{2i\gamma \phi _2(z^{})}\left(\nu (z)\nu ^{}(z^{})\right):_z\delta _{2\pi }(zz^{})`$
$`{\displaystyle \frac{\mathrm{}}{8\pi \gamma ^2}}\left(\eta ^2+1+{\displaystyle \frac{\mathrm{}\gamma ^2}{\pi }}\right):e^{2i\gamma \phi _2(z)}e^{2i\gamma \phi _2(z^{})}:_z^2\delta _{2\pi }(zz^{}).`$
Again we used (95) and its derivative $`\delta ^+(z)_z\delta ^+(z)=i_z^2\delta ^+(z)/(4\pi )`$. With the help of (66) as well as standard properties of the delta function (A) reduces to
$`A(z,z^{})=`$ $`[{\displaystyle \frac{\eta ^2}{\gamma ^2}}(:(_z\varphi _1)^2(z^{}):+:(_z\varphi _2)^2(z^{}):)+{\displaystyle \frac{\mathrm{}\eta }{2\pi \gamma }}_z^2\varphi _1(z^{})`$ (102)
$`+{\displaystyle \frac{\mathrm{}}{16\pi ^2}}{\displaystyle \frac{\mathrm{}}{4\pi \gamma ^2}}(_z+{\displaystyle \frac{i\gamma p_2}{2\pi }})^2]\delta _{2\pi }(zz^{}).`$
Inserting this into (99) gives (8).
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# Testing the Symmetrization Postulate on Molecules with Three Identical Nuclei
## I Introduction
In the latest years, the possibility of small violations of the symmetrization postulate (SP) and the spin-statistics connection has been addressed from both the experimental and the theoretical point of view. The SP is at the basis of the quantum-mechanical description of systems composed by identical particles, asserting that only wave functions completely symmetric or antisymmetric in the permutation of the particles labels can describe physical states. In principle there would be no arguments against the existence of states with different symmetries, although they lack some of the properties which are peculiar of the completely symmetric and antisymmetric states. The only strict requirement that can be derived in a formal way in Quantum Mechanics is the so called superselection rule, which forbids transitions between different symmetry classes.
In addition, the experimental observation seems to indicate so far a well defined connection between the spin of the particles and the symmetry properties: half-integer spin particles can be described only by antisymmetric wave functions, while integer-spin particles by symmetric wave functions only. According to their behavior in a statistical ensemble, the two kinds of particles are indicated respectively as fermions and bosons.
Although both the SP and the spin-statistics connection seem to hold in the physical world, there is great interest in investigating possible very small violations, not detected so far. In case of two-particles systems, only the symmetrical and antisymmetrical states can be defined, and therefore the SP is redundant. As a matter of fact, the experimental tests of symmetry reported so far have been limited to the investigation of states containing only two identical particles, looking for violations of the usual connection between spin and statistics. In particular the latest high-sensitivity experiments of this kind have been performed on both fermionic systems, namely electrons, and bosonic systems, such as photons and oxygen nuclei. In the case of electrons, the existence of totally symmetric states in the spectrum of atomic helium has been searched , and the negative result of such investigation has led to the establishment of a bound to the degree of violation of the spin-statistics, i.e. to the presence of exchange-symmetric electrons. A different experiments was performed on a metallic Cu sample , looking for the presence of any of the electrons in the conduction band with no symmetry requirements with respect to the electrons in the already filled shells. This can be considered a one-particle test of the identity of the electrons, which is even a more general property of the particles, but not a test of their symmetry. Results which are conceptually identical to the first case have been obtained for bosonic particles, by searching for an exchange-antisymmetric two-photons transition in atomic barium , and for rotational states in molecular oxygen and carbon dioxide which are antisymmetrical in the exchange of O nuclei. The principle of all these experiments was to look for particular atomic or molecular transitions, whose probability would be zero according to the usual spin-statistics, by using high sensitivity spectroscopic detection techniques. The extrapolated bounds on the relative weight of the wrong statistics are almost at the level of 10<sup>-9</sup> . Reviews of different experiments which have been proposed, and interpretations of the results of various measurements as tests of the statistics, can be found in .
Theoretical studies of possible violations of the usual symmetry properties have focused on systems composed by a non-fixed number of particles, which can be treated within the Quantum-Field Theory formalism. In such context the spin-statistics connection is no more a postulate dictated by the experimental observation, but it can be rigorously proved through the spin-statistic theorem, the main assumption being the validity of the SP. Several theories allowing for small violations of such postulate have been proposed , and one of the most successful is based on the $`q`$-mutator algebra , in which deformed bilinear commutation relations are used in place of the ordinary ones. In this way it is possible to define exotic symmetries as a smooth interpolation between the symmetric and antisymmetric ones, or statistics interpolating between the Bose and Fermi cases. By assuming the validity of this algebra it has also been possible to translate the bound on the spin-statistics for oxygen nuclei , into a bound on the statistics of the nucleons composing the nuclei or, more generally, to establish a relationship between constraints on different kinds of particles.
To summarize, so far experiments have been devoted to search possible violations of the spin-statistics in simple systems, while theories allowing for violations of the more fundamental SP were focused on much more complex systems, which cannot be easily investigated experimentally. We note that no rigorous theoretical study of a simple quantum-mechanical system allowing for violations of the SP has been reported so far, although possible experiments on multi-particle molecular systems have been indicated .
In this paper we propose an experiment to be performed on molecular systems containing three identical nuclei, to search for possible violations of the SP in a system described within the quantum-mechanical theory. For this purpose, in Section II we briefly recall the concept of identity of particles, and the permutation properties of identical particles; we also recall that at least three particles are needed to define symmetries different from the usual ones. Then, in Section III, we discuss the properties of particular rotational states of symmetrical plane molecules with three spin-0 nuclei, which are forbidden by the SP, and with three spin-1/2 nuclei, which are forbidden by both the spin-statistics connection and the SP. We show that, nevertheless, such states have the proper symmetry to be defined consistently in a more general quantum-mechanical theory which does not include the SP. In Section IV we briefly discuss how to test the possible small symmetry violations. Finally, in Section V, we propose a high-sensitivity spectroscopic investigation to be performed on selected transitions of SO<sub>3</sub> molecule, and of BH<sub>3</sub> and NH<sub>3</sub> molecules, which best represent the prototypal molecule respectively for the spin-0 and spin-1/2 cases treated above. In the latter case, such states could be investigated in search of small violation of the statistics or of the symmetry of protons.
## II The indistinguishability of identical particles and the symmetrization postulate
### A General discussion
The indistinguishability of identical particles is one of the principles lying at the basis of Quantum Mechanics. It states that given a system of particles belonging to the same species (i.e. which have the same physical properties, like mass, charge, spin, etc.), a permutation of such particles cannot lead to any observable effect . To discuss the implications of such principle, here we briefly recall some properties of a system containing $`N`$ identical particles (for a general discussion see for example Ref. ).
To each observable of the system is associated an hermitian operator $`\widehat{A}`$, which can be written as
$$\widehat{A}=A(\{\widehat{O}_1\},\mathrm{}\{\widehat{O}_N\})\widehat{A}(1,\mathrm{},N)$$
(1)
where $`\{\widehat{O}_j\}`$ is a set of single-particle observables. In the same way, the vector
$$|1,2,\mathrm{},N$$
(2)
represents a state where the first particle is characterized by a set of quantum numbers $`\{1\}`$, the second by $`\{2\}`$, and so on. Moreover we can define the permutation operators as
$$\widehat{P}_{j_1,j_2,\mathrm{},j_N}f(1,2,\mathrm{},N)=f(j_1,j_2,\mathrm{},j_N)$$
(3)
and these operators form the symmetrization group $`S_n`$. Therefore, since a permutation of identical particles cannot be observed in any experiment, all quantum observables must be permutation-invariant
$$[\widehat{P},\widehat{A}]=0$$
(4)
and since the evolution operator $`\widehat{U}(t)`$ is related to the Hamiltonian of the system $`\widehat{H}`$, which is a physical observable, the above condition is fulfilled at any instant of time
$$[\widehat{P},\widehat{H}]=0[\widehat{P},\widehat{U}(t)]=0.$$
(5)
An important consequence of these relations is that the Hilbert space $``$ can be written as a direct sum of hortogonal subspaces, each one invariant under the permutation group, the hortogonality being preserved along the evolution of the system, and following a physical measurement. This result can be stated as a super-selection rule which forbids transitions between states transforming under inequivalent representations of the permutation group.
To investigate in detail the symmetry properties of a generic state of the system under permutation, it is convenient to introduce the operator of a cyclic permutation
$$\widehat{P}_{(k_1,k_2,\mathrm{},k_K)}f(k_1,k_2,\mathrm{},k_K)=f(k_2,\mathrm{},k_K,k_1)$$
(6)
and the exchange operators of two particle, $`\widehat{P}_{(j,k)}`$ (they are a particular case of the former), which are hermitian and unitary, with eigenvalues $`ϵ_{j,k}=\pm 1`$.
Each generic permutation can be written either as the product of exchange between two particles, or of cyclic permutation of distinct elements, in the form
$$\widehat{P}_{cyclic}\widehat{P}_{(k_1,\mathrm{},k_p)(k_{p+1},\mathrm{},k_{p+q})(k_{p+q+1},\mathrm{},k_{p+q+r})\mathrm{}}$$
(7)
with the only condition that the sum of each cycle length must be equal to $`N`$
$$p+q+r+\mathrm{}=N$$
(8)
(here we adopt the convention $`pqr\mathrm{}`$). The properties of transformation of such operators under permutation do not change the length of the cycles, and therefore all the cyclic permutation with the same values of $`p,q,r,\mathrm{}`$ are said to form a class. Obviously there are as many classes as many are the possible decomposition of $`N`$ as sum of integer numbers.
From the properties of the exchange operators $`\widehat{P}_{(j,k)}`$ it follows that the permutation operators admit only two common eigenstates, with the eigenvalues $`ϵ_{j,k}`$ all equal to $`1`$ or $`1`$. In fact, the eigenvalue equation for the exchange operators $`\widehat{P}_{(j,k)}`$ is ($`ϵ_{j,k}=\pm 1`$)
$$\widehat{P}_{(j,k)}f(1,2,\mathrm{},N)=ϵ_{j,k}f(1,2,\mathrm{},N)$$
(9)
then, from the relation
$$\widehat{P}_{(j,k)}=\widehat{P}_{(2,k)}\widehat{P}_{(1,j)}\widehat{P}_{(1,2)}\widehat{P}_{(2,k)}\widehat{P}_{(1,j)}$$
(10)
it follows
$$ϵ_{j,k}=ϵ_{1,2}ϵ_{2,k}^2ϵ_{1,j}^2=ϵ_{1,2}.$$
(11)
These eigenvalues $`ϵ_{j,k}=\pm 1`$ correspond respectively to vectors of state completely symmetric or antisymmetric under the exchange of two particles, and we indicate the subspaces spanned by these vectors by $`_+`$ and $`_{}`$ respectively. For $`N=2`$ the Hilbert space can be decomposed exactly as $`=_+_{}`$, i.e. only symmetric or antisymmetric states are possible. For $`N>2`$ this is no more true, and we have
$$=_+_{}_1^{}_2^{}\mathrm{},$$
(12)
where $`_j^{}`$ are permutation-invariant subspaces. Contrarily to $`_+`$ and $`_{}`$, which have a definite symmetry and are one-dimensional representation of the permutation group (for a given set of quantum numbers $`\{1,2,\mathrm{},N\}`$ there exist only one possible combination which gives symmetric or antisymmetric states), the subspaces $`_j^{}`$ do not posses a definite symmetry and have dimension greater that one. This is the reason why one usually assumes the symmetrization postulate (SP), requiring physical states to be either symmetric or antisymmetric. Nevertheless, there is no stringent reason to exclude a priori the possibility of physical states not obeying SP, since their presence do not violate any basic principle of Quantum Mechanics. The only difference that one has to take into account is the lack of a correspondence between physical states and vectors (modulo a phase factor), since the subspaces $`_j^{}`$ have dimension greater that one and do not admit a complete set of mutually commutating physical observables. In fact, let us imagine to have a common eigenvector $`|u`$ of a complete set of operators $`\{O_i`$} in $`_j^{}`$. Since $`|u`$ can not be eigenvector of all the permutation operator, there must exist some permutation such that $`\widehat{P}_{j_1,j_2,\mathrm{},j_N}|u`$ is linearly independent of $`|u`$, and therefore $`_j^{}`$ must have dimension greater that one, and the set of operators $`\{O_i\}`$ can not be complete (i.e. the state of the system is not completely determined by the set of eigenvalues $`\{o_i\}`$).
In the following we work out in more detail the permutation properties for the special case n=3 .
### B A system of 3 identical particles
The symmetrization group $`S_3`$ is a non-abelian group of order 6, formed by the $`3!`$ cyclic permutations of the labels $`(1,2,3)`$, which belong to three distinct classes
$$P_{(1)(2)(3)}=I$$
(13)
$$P_{(2,3)};P_{(1,3)};P_{(1,2)}$$
(14)
$$P_{(1,2,3)}=P_{231};P_{(3,2,1)}=P_{312}.$$
(15)
This group is isomorphic to the dihedral group $`D_3`$, generated by the symmetry transformation shown in Fig. 1. In fact, reflection about each axis ($`a,b,c`$) is equivalent to a two label permutation $`P_{(i,j)}`$, while rotations around the center by angles $`2\pi /3`$ and $`2\pi /3`$, lead to cyclic permutation of all the three labels. The visualization of the group $`S_3`$ by its association to the geometrical symmetries of $`D_3`$ is very helpful since it is indeed this realization on a physical system (a molecule in our case) that we are looking for.
The permutation operators can be written explicitly in a matrix form, on the basis formed by the $`3!`$ permutations of the vector $`|1,2,3`$. For example, by defining
$`|1,2,3|1`$ (16)
$`|1,3,2|2`$ (17)
$`|2,1,3|3`$ (18)
$`|2,3,1|4`$ (19)
$`|3,1,2|5`$ (20)
$`|3,2,1|6`$ (21)
we obtain
$$P_{(1,2,3)}=\left(\begin{array}{cccccc}0& 0& 0& 0& 1& 0\\ 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 0& 0& 1\\ 1& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 0& 0\\ 0& 1& 0& 0& 0& 0\end{array}\right),$$
(22)
and
$`P_{(3,2,1)}=\left(\begin{array}{cccccc}0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 0& 1\\ 0& 1& 0& 0& 0& 0\\ 0& 0& 0& 0& 1& 0\\ 1& 0& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 0\end{array}\right).`$ (23)
(24)
This matrix representation of the group $`S_3`$ can be reduced, since it must contain a number of irreducible representation equal to the number $`n_c`$ of distinct classes ($`n_c=3`$). Moreover, since this is the so called regular representation, each irreducible representation must appear a number of times equal to its dimensionality. From general results we know that these representations are the symmetric and antisymmetric ones (which are 1-dimensional), and that there must be also a 2-dimensional representation occurring two times. A formal way to systematically reduce the regular representation of $`S_n`$ into its irreducible components would be by using the Young diagrams . Here we work out the S<sub>3</sub> case with more heuristic arguments, as follows.
First of all, we can define the operators $`S`$ and $`A`$, which are respectively the projectors on the subspaces $`_+`$ and $`_{}`$
$`S`$ $`=`$ $`\left[I+P_{213}+P_{312}+P_{132}+P_{312}+P_{231}\right]/6`$ (25)
$`A`$ $`=`$ $`\left[IP_{213}P_{312}P_{132}+P_{312}+P_{231}\right]/6.`$ (26)
Each of them has therefore only one eigenvector associated to a non-zero eigenvalue, which are respectively
$`|s`$ $``$ $`\left[|1+|2+|3+|4+|5+|6\right]/\sqrt{6}`$ (27)
$`|a`$ $``$ $`\left[|1+|2+|3|4|5+|6\right]/\sqrt{6}.`$ (28)
and are the standard symmetric and antisymmetric states.
Then, in order to generate the two 2-dimensional irreducible representations of $`^{}`$, we can try to diagonalize the permutation operators $`P_{(1,2,3)}`$ and $`P_{(3,2,1)}`$. The eigenvectors of the former, and their corresponding eigenvalues $`\lambda `$, are
$`\{\begin{array}{cc}|v_1\lambda _{}|2+\lambda _+|3+|6\hfill & \\ & \\ |v_2\lambda _+|1+\lambda _{}|4+|5\hfill & \end{array};\lambda _{}=\mathrm{e}^{i{\displaystyle \frac{2\pi }{3}}}`$ (29)
(30)
$`\{\begin{array}{cc}|v_3\lambda _+|2+\lambda _{}|3+|6\hfill & \\ & \\ |v_4\lambda _{}|1+\lambda _+|4+|5\hfill & \end{array};\lambda _+=\mathrm{e}^{+i{\displaystyle \frac{2\pi }{3}}}`$ (31)
The other cyclic permutation, $`P_{(3,2,1)}`$, has the same eigenvectors, with the eigenvalues exchanged.
By looking at the behavior of such states under the exchange of two particles, it is easy to find out that the two invariant subspaces $`_1^{}`$ and $`_2^{}`$ are generated respectively by $`\{v_1,v_4\}`$ and $`\{v_2,v_3\}`$. Therefore, the effect of a generic exchange is a rotation of the vectors in $`_j^{}`$, with a change of the eigenvalues of $`P_{(1,2,3)}`$ and $`P_{(3,2,1)}`$. For example, the action of $`P_{(2,3)}`$ on $`v_1`$ gives
$$\left(\begin{array}{cccccc}0& 1& 0& 0& 0& 0\\ 1& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 0& 0\\ 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 0& 0& 1\\ 0& 0& 0& 0& 1& 0\end{array}\right)\left(\begin{array}{c}0\\ \mathrm{e}^{i\frac{2}{3}\pi }\\ \mathrm{e}^{+i\frac{2}{3}\pi }\\ 0\\ 0\\ 1\end{array}\right)=\left(\begin{array}{c}\mathrm{e}^{i\frac{2}{3}\pi }\\ 0\\ 0\\ \mathrm{e}^{+i\frac{2}{3}\pi }\\ 1\\ 0\end{array}\right).$$
(32)
The resulting vector is $`v_4`$, which corresponds to a different eigenvalue of both cyclic permutations. The same result applies for the other two exchanges.
In the next Section we investigate the consequences of these properties of the cyclic permutations for the specific case of three nuclei bound in a rigid molecule, and the possible physical implication of the existence of these states not obeying SP.
## III Symmetry-forbidden states in symmetrical 3-nuclei molecules
### A Three spin-0 identical nuclei
The prototypal molecule that we want to consider is composed by three identical spin-zero nuclei disposed at the vertices of an equilater triangle, as sketched in Fig. 1, whose symmetries belong to the dihedral group $`D_3`$. Anyway, the discussion below can be extended to any symmetrical molecule with additional nuclei in the plane (point group $`D_{3h}`$).
Within the Born-Oppenheimer approximation, the total wave function can be decomposed in the usual way as
$$\mathrm{\Psi }=\mathrm{\Psi }_e\mathrm{\Psi }_n\mathrm{\Psi }_v\mathrm{\Psi }_r,$$
(33)
where the partial wave functions are respectively the electronic, nuclear spin, vibrational and rotational components. To simplify our description, we can assume the molecule to be in the ground electronic and vibrational state, with species $`{}_{}{}^{1}\mathrm{\Sigma }`$ (which is the most common case), or, in other words, we assume the total electronic spin, electronic angular momentum and vibrational angular momentum to be zero. Therefore, the electronic and vibrational wave functions $`\mathrm{\Psi }_e`$ and $`\mathrm{\Psi }_v`$ are completely symmetric under any permutations of the labels of the nuclei. In the case of spin-zero nuclei also the nuclear spin wave-function $`\mathrm{\Psi }_n`$ is symmetric under permutations, and therefore we must consider the symmetry of the rotational wave-function alone.
For this purpose, we briefly recall the main properties of the rotational states of such a molecule, which are eigenstates of the rigid rotational Hamiltonian
$$H=B(𝐉/\mathrm{})^2(BC)(𝐊/\mathrm{})^2,$$
(34)
where $`B`$ and $`C`$ are the two rotational constants, J is the total angular momentum and K is its projection on the three-fold symmetry axis perpendicular to the molecule’s plane. The eigenvalues of the Hamiltonian are of the form
$$E(J,K)=BJ(J+1)(BC)K^2,$$
(35)
with $`K=J,J+1,\mathrm{},J`$, and are strictly degenerate in the sign of $`K`$ as far as the molecule is symmetric, even if non-rigidity and perturbations are considered (the degeneracy accounts for the undistiguishability of the two orientations of K with respect to the molecular axis). The eigenfunctions are somewhat complicated functions of $`J`$, $`K`$, and of the Eulerian angles ; they also depend parametrically on the equilibrium positions of the three nuclei in the molecular potential. Here we are interested in their transformation properties under a generic rotation of the molecule. The general rule for such transformations is the usual
$$\mathrm{\Psi }_r\mathrm{\Psi }_re^{i\alpha 𝐮_\alpha 𝐉},$$
(36)
where $`\alpha `$ is the angle of rotation and $`𝐮_\alpha `$ is the versor of the plane containing $`\alpha `$. To find the symmetry character of the rotational states, we want to compare this transformation rule to those associated to the permutations of particles. For such purpose we write the generic rotational states in a formal way as a function of the coordinates of the three minima in the molecular potential
$$|𝐱_\mathrm{𝟏},𝐱_\mathrm{𝟐},𝐱_\mathrm{𝟑}.$$
(37)
These states are defined in the 6-dimensional Hilbert space defined in the previous Section. We note that a permutation of the quantum labels corresponds to a classical permutation of the nuclei’s mean positions, and therefore all the permutations are equivalent to rotations of the molecule.
We start by considering rotations in the plane of the molecule by angles $`\theta `$=$`ϵ2\pi /3`$, with $`ϵ`$=$`\pm `$1, which are equivalent to the cyclic permutations $`P_{(1,2,3)}`$, $`P_{(3,2,1)}`$. Since the effect on the molecular wave function of these rotations must be the same of such cyclic permutations, we can compare the general transformation rules under rotations and under permutations. For the formers, the appropriate phase shift is given by
$$\mathrm{\Psi }_r\mathrm{\Psi }_r\mathrm{exp}^{iϵ2\pi K/3},$$
(38)
since the rotation is in the molecule’s plane, while for the permutations we have
$`\mathrm{\Psi }_r\mathrm{\Psi }_r\mathrm{and}\mathrm{\Psi }_r\mathrm{\Psi }_r\mathrm{exp}^{iϵ2\pi /3},`$ (39)
respectively for states in $`_+`$, $`_{}`$ and in $`^{}`$. Therefore, a state with $`K`$=3$`q`$, with $`q`$ integer, can be defined in $`_+`$ or $`_{}`$, while states with $`K`$=3$`q\pm `$1 can be defined only in $`^{}`$.
For the special case of a rotational state with $`K`$=0 we can consider also rotations about any of the three symmetry axes in the plane of the molecule by a angles $`\varphi `$=$`ϵ\pi `$. Such rotations are equivalent to the three exchanges $`P_{(1,2)}`$, $`P_{(2,3)}`$ and $`P_{(1,3)}`$, and also the correspondent phase shifts in the wave function must be identical. This class of rotations transform the wave function as
$$\mathrm{\Psi }_r\mathrm{\Psi }_r\mathrm{exp}^{iϵ\pi J},$$
(40)
according to Eq. (36), since the angular momentum J lays in the molecule’s plane. The exchanges are instead characterized by
$`\mathrm{\Psi }_r\mathrm{\Psi }_r\mathrm{and}\mathrm{\Psi }_r\mathrm{\Psi }_r\mathrm{exp}^{i\pi },`$ (41)
for states belonging respectively to $`_+`$ and $`_{}`$. As a result, the even-J states belong to $`_+`$, and the odd-J ones to $`_{}`$, as summarized in Tab. I.
From the results obtained above, it appears that all the rotational states with $`K`$=3$`q\pm `$1 are strictly forbidden by the SP, which requires any physical state to be defined in $`_+`$ or $`_{}`$. Moreover, the odd-$`J`$ states with $`K`$=0 are forbidden by the spin-statistics, since the identical nuclei are spin-0 particles. The latter case is very similar to that of symmetrical molecules with two identical nuclei . These consequences of both the SP and the spin statistics are well known in the field of molecular spectroscopy , since they lead to the absence of more than two thirds of the rotational states in molecules with the proper symmetry, as confirmed by the experimental investigation performed so far.
We now discuss a property of the SP-forbidden states, which can be defined in the unsymmetrical subspace $`^{}`$. Specifically, the states with $`K`$=3$`q`$-1 could be described by the vectors $`v_1,v_2`$, while the ones with $`K`$=3$`q`$+1 by the vectors $`v_3,v_4`$ to preserve the equality of the phase shifts reported in Eq. (38) and Eq. (39). As we noted, such pairs of vectors do not define invariant subspaces, and therefore the physical states have to be defined in the two invariant subspaces $`_1^{}`$ and $`_2^{}`$. Although in each of these subspaces the sign of the rotational phase shift is not defined (see Eq. (30)), also the sign of $`K`$ is not a physical observable, due to the degeneracy of the states with $`\pm K`$. In other words, even if in $`^{}`$ it is not possible to define a complete set of operators due to its multi-dimensionality, and in particular the sign of $`K`$ is undefined, no paradox arises, since we have no means to measure such a sign.
### B Three spin-1/2 identical nuclei
In case of a molecule containing three identical nuclei with non-zero spin, we should consider also the influence of the nuclear spin wave function; in the following we will discuss the special case of spin 1/2. Since the spin is non-zero, we can now define a non-trivial Hilbert space for the spin, in the same way we defined it for the spatial coordinates. The total Hilbert space will be the direct product of the rotational and spin spaces
$$=_S_R.$$
(42)
The general properties of the Hilbert space can be easily found, considered the discrete nature of the quantum labels associated to each particle. The generic three-particles state is now defined as
$$|S_{z1},S_{z2},S_{z3},$$
(43)
and there are two states with total spin $`I`$=1/2 and one with $`I`$=3/2 (two doublet and one quartet). Using the formalism of Section II it is possible to assign the quartet state to the symmetrical subspace $`_+`$, and the two doublet to the unsymmetrical $`^{}`$. We note that is not possible to build non-vanishing states completely antisymmetric under permutations of more than two spin-1/2 particles, since each quantum label can assume only two values.
Even in absence of an electronic or vibrational angular momentum, the nuclear spin can couple to the rotational angular momentum and give rise to a very small hyperfine splitting of the rotational states. The molecular states will therefore be eigenstates of the total angular momentum F=J+I. We want now to assign each molecular state to a particular subspace of the overall Hilbert space, as we did in the previous case. For such purpose we have to decompose the representation of the S<sub>3</sub> group on the direct product Hilbert space $`_S_R`$, into irreducible representations. By using general rules we easily obtain the decomposition shown in Tab. II.
It is then straightforward to evaluate the appropriate Hilbert subspace for the overall rotational and nuclear spin states, using the associations of Tab. I; the results are reported in Tab. III. If the hyperfine splitting of the rotational states is not resolved, as it is usually the case, the states with $`K`$ 0 are no more completely forbidden by the SP, while the even-$`J`$ states with $`K`$=0 are forbidden by both the SP and the spin-statistics. Specifically, the states defined in $`_+`$ will have $`I`$=3/2, and the ones in $`^{}`$ $`I`$=1/2. As a consequence of the multidimensionality of $`^{}`$, it is impossible to distinguish between the two doublet states, but this is not a problem, since they are rigorously degenerate. In fact the invariant subspaces of $`^{}`$ are eigenstates of $`I_z`$, and therefore also in presence of a magnetic field there is no fundamental principle other than the SP to forbid the existence of these unsymmetrical states.
We conclude by noting that in the case of a nuclear spin larger than 1/2 it is possible to build also antisymmetrical spin states in $`_{}`$, and therefore none of the rotational states with $`K`$=0 is forbidden.
### C Non-planar symmetrical molecules
We now extend the discussion to non-planar symmetrical molecules containing three spin-0 or spin-1/2 identical nuclei (point group $`C_{3v}`$: the molecule has additional non-identical nuclei out of the plane containing the identical nuclei). The main difference from the planar case stands in the transformation rule of the $`K`$=0 rotational states under rotations about a symmetry axis in the plane of the three identical nuclei. These rotations, by angles $`\varphi `$=$`ϵ\pi `$, are not identical to permutations of two identical nuclei, as it was for a planar molecule, since also the additional nuclei out of the plane has rotated. On the other hand, the combination of such a rotation with an inversion of the coordinates of all the nuclei with respect to the origin cannot be distinguished from a permutation of two nuclei.
The properties of the rotational states under space inversion are well known in the field of molecular spectroscopy , and here we give only the main result. If the non-rigidity of the molecule is taken into account, each rotational state is split into two states with definite symmetry under inversion (s\- and a-species, respectively symmetric and antisymmetric under inversion). The splitting of the degeneracy of these pairs of states is indicated as inversion splitting. The analogous of the transformation rule of Eq. (40) for the rotation plus inversion is therefore
$$\mathrm{\Psi }_r\pm \mathrm{\Psi }_r\mathrm{exp}^{iϵ\pi J},$$
(44)
where the positive (negative) sign refer to the s-species (a-species). If we compare this rule with the one for the exchange of two particles (Eq. (41)) we can assign each rotational state to the proper symmetry class, as summarized in Tab. IV for the cases of spin-0 and spin-1/2.
In the last Section we will indicate one molecular species in which the inversion splitting is sufficiently large to resolve the states forbidden by the SP from those allowed.
## IV Testing small symmetry violations
As we noted above, those rotational states forbidden by the SP can be defined consistently in the unsymmetrical subspace $`^{}`$, since they are not violating any basic principle of Quantum Mechanics. The discussion can be extended to particular classes of rotational sub-states in excited vibrational and electronic states; due to the different symmetry of the vibrational and electronic wave functions, different classes of $`K`$-levels and $`J`$-levels can be forbidden by the SP. In any case, electromagnetic transitions between rotational levels are allowed only if the involved states belong to the same symmetry type, in accordance with the superselection rule.
Therefore, a violation of the SP can be thought of only in terms of the appearance of a population of nuclei identical between themselves, whose quantum states belong to $`^{}`$, and whose properties (mass, spin, charge) are the same as those of the normal nuclei. The presence of a non-zero population in SP-forbidden rotational states of the molecules formed by such nuclei could be detected by exciting electromagnetic transitions towards other forbidden states, in the same way it was done in previous experiments for testing the statistics. For a matter of sensitivity, the experimental investigation should focus on states completely forbidden by the SP alone, or by both the SP and the spin-statistics, since we expect the violation, if present, to be very small. Therefore, the most interesting states to be probed in both $`D_{3h}`$ and $`C_{3v}`$ molecules are those with $`K`$=3$`q\pm `$1 in the case of $`S`$=0, and those with $`K`$=0 in the case of $`S`$=1/2.
We note also that all the statements above would not lose validity in case of breakdown of the Born-Oppenheimer approximation. Indeed, the symmetry character of the particles under permutation can only be assigned to a single class (i.e. Hilbert subspace), as stated by the super-selection rule, and no coherent superposition is allowed . Therefore, if the total wave function appears to be correctly described as symmetrical (unsymmetrical) in the frame of the Born-Oppenheimer approximation, it would be so for any degree of violation of the approximation itself. The same argument applies in case of presence of external perturbations, such as electric or magnetic fields: even if they can change the molecular wave function to a large extent in a continuous way, the symmetry character of the wave function will be locked to a single class.
## V Proposal for a high-sensitivity test of the symmetry on oxygen nuclei and protons
A simple tool to investigate for possible violations of the SP would be high sensitivity spectroscopy of a thermal sample of gas, in search for a non-zero absorption of light by symmetry-violating molecules, as performed in previous experiments on the spin-statistics. It is possible to define a few general criteria to choose the proper molecule for such an experimental search. The first important parameter is the weight of the molecule, which determines the spacing between adjacent rotational transitions (the rotational constants $`B`$ and $`C`$ are inversely proportional to the mass), and therefore also the capability of resolving the forbidden transitions from the allowed ones . This point is very important, since in non-linear molecules the number of transitions which can possibly interfere in the detection of SP-violations is increased with respect to linear molecules, due to the increased molteplicity of rotational states. Secondly, in order to have a high sensitivity in detecting a violation of the SP, one should probe strong transitions starting from low-laying rotational states in the ground electronic and vibrational state, which have the largest occupation probability. As usual, in this kind of molecules the strong electronic transitions are confined to the UV, where laser sources are not easily available, and therefore one should rely on the fundamental vibrational transitions in the infrared, which are somewhat weaker. It is therefore important to find a near coincidence between a vibrational band and the emission of a coherent laser source suitable for high sensitivity spectroscopy. The quantum-cascade semiconductor lasers are particularly interesting from this point of view, since they can be designed to emit almost everywhere in the mid-infrared, they are widely tunable and have low-noise characteristics .
In the case of spin-0 nuclei, the lightest molecule with the proper symmetry is the plane SO<sub>3</sub> (group $`D_{3h}`$), whose vibrational energies are reported in Tab. V. In principle the sensitivity for detection of SP-violating transitions can be expected to be comparable to those obtained in previous experiments on molecular systems. It could possibly be reduced by increased experimental difficulties (for example, SO<sub>3</sub> is not chemically stable in presence of oxygen, and therefore may not be easy to have the proper pressure and cell volume to optimize the sensitivity) or more fundamental limitations, such as the presence of weak hot-transitions close to the SP-forbidden ones.
An alternative approach could be to perform the spectroscopic investigation on a supersonic beam of SO<sub>3</sub>, instead of a thermal sample. With this technique, the problems connected with the high reactivity of SO<sub>3</sub> could possibly be solved, while the reduction in density of the sample could be compensated by the reduction of the rotational temperature and by an increase of the detection sensitivity. In addition, one would benefit from the accompanying increase in resolution, for a better identification of the detected spectrum of transitions.
As for the case of spin-1/2 nuclei, it is possible to find a relatively light plane molecule, BH<sub>3</sub> (group $`D_{3h}`$), on which to search violations of the SP and of the spin-statistics for $`K`$=0 rotational states. As we mentioned, the same symmetry properties can be found also for the corresponding rotational states of non-planar symmetrical molecules (group $`C_{3v}`$) with inversion splitting. One of the most interesting species of this kind is NH<sub>3</sub>, which is stable and particularly light, and has very strong vibrational absorption bands in the infrared spectral region covered by semiconductor lasers. Moreover, its spectrum has been the subject of extensive investigation, and therefore the assignment of observed lines can be performed in a relatively easy way. In addition, it is characterized by a very large inversion splitting of the rotational spectrum, and therefore the forbidden lines can be easily distinguished form the allowed ones. For such molecules, which are composed by spin-1/2 nuclei, we noted that for the $`K`$=0 states different hyperfine components are forbidden by the SP and by the spin-statistics. Since the hyperfine slitting of the rotational levels is usually negligible (at least for $`\mathrm{\Sigma }`$ electronic states ), we can expect to resolve the two classes of substates only thanks to Zeeman effects in strong magnetic fields.
To conclude, we note that in the case of BH<sub>3</sub> or NH<sub>3</sub>, one would be testing the SP and the spin-statistics for protons. Such an experimental test would be particularly interesting to compare the present bound for violation of the statistics for composed nucleons to a corresponding bound for fundamental nucleons, also in view of the recent theoretical predictions .
## VI Conclusions
We have shown how to define consistently particular classes of rotational levels of symmetrical molecules containing three identical nuclei, by allowing for violations of the symmetrization postulate. This violation, if existing, must be small, according to the results of the past experimental observations. However, it is possible to look with high sensitivity for the presence of tiny population in SP-forbidden states, using spectroscopic tools in a scheme similar to those of previous experiments on O<sub>2</sub> and CO<sub>2</sub>. If we allow for the presence of molecules composed by SP-violating nuclei in a sample of gas, which are identical to the normal molecules in all but the symmetry, then such a measurement can be interpreted as a test of the SP. This kind of experiment would represent a substantial improvement with respect to the past investigations on possible violations of the spin-statistics connection. We propose a few simple molecules containing spin-0 nuclei (SO<sub>3</sub>) or spin-1/2 nuclei (BH<sub>3</sub> and NH<sub>3</sub>) as candidates for a high sensitivity spectroscopic investigation to be performed on infrared vibrational transitions, with the help of semiconductor lasers. The latter species are particularly interesting, since they represent systems on which to test the SP for fundamental nucleons.
###### Acknowledgements.
We acknowledge stimulating discussions with G. M. Tino.
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# 1 Introduction
## 1 Introduction
This is the third paper where we study the theory of scattering and more precisely the existence of modified wave operators for a class of long range Hartree type equations
$$i_tu+\frac{1}{2}\mathrm{\Delta }u=\stackrel{~}{g}(|u|^2)u$$
$`(1.1)`$
where $`u`$ is a complex function defined in space time $`IR^{n+1}`$, $`\mathrm{\Delta }`$ is the Laplacian in $`IR^n`$, and
$$\stackrel{~}{g}(|u|^2)=\kappa t^{\mu \gamma }||^{\mu n}|u|^2$$
$`(1.2)`$
with $`||=(\mathrm{\Delta })^{1/2}`$, $`\kappa IR`$, $`0<\gamma 1`$ and $`0<\mu n`$. For $`\mu <n`$, the operator $`||^{\mu n}`$ can be represented by the convolution in $`IR^n`$
$$||^{\mu n}f=C_{n,\mu }|x|^\mu f$$
$`(1.3)`$
so that (1.2) is a Hartree type interaction with potential $`V(x)=C|x|^\mu `$. The more standard Hartree equation corresponds to the case $`\gamma =\mu `$. In that case, the nonlinearity $`\stackrel{~}{g}(|u|^2)`$ becomes
$$\stackrel{~}{g}(|u|^2)=V|u|^2=\kappa |x|^\gamma |u|^2$$
$`(1.4)`$
with a suitable redefinition of $`\kappa `$.
A large amount of work has been devoted to the theory of scattering for the Hartree equation (1.1) with nonlinearity (1.4) as well as with similar nonlinearities with more general potentials. As in the case of the linear Schrödinger equation, one must distinguish the short range case, corresponding to $`\gamma >1`$, from the long range case corresponding to $`\gamma 1`$. In the short range case, it is known that the (ordinary) wave operators exist in suitable function spaces for $`\gamma >1`$ . Furthermore for repulsive interactions, namely for $`\kappa 0`$, it is known that all solutions in suitable spaces admit asymptotic states in $`L^2`$ for $`\gamma >1`$, and that asymptotic completeness holds for $`\gamma >4/3`$ . In the long range case $`\gamma 1`$, the ordinary wave operators are known not to exist in any reasonable sense , and should be replaced by modified wave operators including a suitable phase in their definition, as is the case for the linear Schrödinger equation. A well developed theory of long range scattering exists for the latter. See for instance for a recent treatment and for an extensive bibliography. In contrast with that situation, only partial results are available for the Hartree equation. For small solutions (or equivalently small asymptotic states) the existence of modified wave operators has been proved in the critical case $`\gamma =1`$ . On the other hand, it has been shown, first in the critical case $`\gamma =1`$ and then in the whole range $`0<\gamma 1`$ that the global solutions of the Hartree equation (1.1) with nonlinearity (1.4) and with small initial data exhibit an asymptotic behaviour as $`t\pm \mathrm{}`$ of the expected scattering type characterized by scattering states $`u_\pm `$ and including suitable phase factors that are typical of long range scattering.
In the previous two papers of this series (hereafter referred to as I and II) we proved the existence of modified wave operators in suitable spaces for the equation (1.1) with nonlinearity (1.2), and we gave a description of the asymptotic behaviour in time of solutions in the ranges of those operators, with no size restriction on the data, first for $`1/2<\gamma <1`$ in I and then in the whole range $`0<\gamma 1`$ in II. The method is an extension of the energy method used in , and uses in particular an auxiliary system of equations introduced in to study the asymptotic behaviour of small solutions. The spaces of initial data, namely in the present case of asymptotic states, are Sobolev spaces of finite order. However there occurs a loss of derivatives in the auxiliary system, which has to be compensated for by the smoothing effect of the operator $`||^{\mu n}`$ in (1.2). This is done in the framework of Sobolev spaces at the expense of assuming $`\mu n2`$, which in particular restrict the space dimension to $`n3`$. In the present paper, we overcome that difficulty by treating the problem in Gevrey spaces , following and extending the method used in to treat the case of small solutions. This makes it possible to cover the whole range $`0<\mu n`$, and in particular the case of dimensions 1 and 2 and the case of cubic nonlinear Schrödinger (NLS) equations (with time dependent nonlinearity). More precisely we use Gevrey classes $`G^{1/\nu }`$ of order $`1/\nu `$ with $`0<\nu 1`$, and the method applies under the condition $`\mu n2+2\nu `$. In particular for cubic NLS equations we need $`\nu =1`$, namely spaces of analytic functions. The previous restriction on $`\mu `$ and $`\nu `$ can still be weakened and has been weakened in at the expense of introducing parabolic terms in the auxiliary system of equations. However those terms introduce a priviledged orientation of time, which is inconvenient for the study of scattering theory, where we like to go back and forth from finite to infinite time, and we shall not make use of that extension here. The origin of the derivative loss and the mechanism by which that loss is overcome in Gevrey spaces, which is the same as in , will be described in Section 3 after sufficient technical material has been introduced, namely after Lemma 3.4.
The construction of the modified wave operators is in its principle the same as in II and will be recalled in Section 2 below, which is mostly a summary of Section 2 of II. It involves the study of the same auxiliary system of equations as in II for an amplitude $`w`$ and a phase $`\phi `$ which replace the original function $`u`$, and the definition of the same modified asymptotic dynamics for that system as in II.
We now give a brief outline of the contents of this paper. A more technical description will be given at the end of Section 2. In Section 3 we define the relevant Gevrey spaces and derive the basic estimates in those spaces that are needed to study the auxiliary system. In Section 4 we prove the existence of the large time dynamics associated with that system and some preliminary asymptotic properties of that dynamics. In Section 5 we study the asymptotic dynamics and we prove the existence of asymptotic states for the previously constructed solutions of the auxiliary system. In Section 6 we construct the local wave operators at infinity for the auxiliary system by solving the Cauchy problem for that system with infinite initial time. We then come back from the auxiliary system to the original equation (1.1) for $`u`$ and construct the (local) wave operators (at infinity) for $`u`$ in Section 7, where the main result is stated as Proposition 7.5.
We have tried to make this paper self-contained and at the same time to keep duplication with I and II to a minimum. Duplication occurs in Section 2, as already mentioned, and in part of Section 7. The more technical sections 3 to 6 follow the same pattern as in II, but there is almost no duplication because the functional framework is significantly different.
We conclude this section by giving some general notation which will be used freely throughout this paper. We shall work mostly in Fourier space. We denote by $``$ the convolution in $`IR^n`$, by $`F`$ the Fourier transform, and by $`\widehat{u}=Fu`$ the Fourier transform of $`u`$. We denote by $`_r`$ the norm in $`L^rL^r(IR^n)`$ and by $`<,>`$ the scalar product in $`L^2`$. For any interval $`I`$ and any Banach space $`X`$, we denote by $`𝒞(I,X)`$ the space of strongly continuous functions from $`I`$ to $`X`$, by $`L^{\mathrm{}}(I,X)`$ (resp. $`L_{loc}^{\mathrm{}}(I,X))`$ the space of measurable essentially bounded (resp. locally essentially bounded) functions from $`I`$ to $`X`$, and by $`L^2(I,X)`$ (resp. $`L_{loc}^2(I,X)`$, resp. $`L_\rho ^2(I,X))`$ the space of measurable functions $`u`$ from $`I`$ to $`X`$ such that $`u();X`$ belongs to $`L^2(I)`$ (resp. $`L_{loc}^2(I)`$, resp. $`L_\rho ^2(I))`$, where $`L_\rho ^2(I)`$ is the weighted space $`L^2(I,\rho (t)dt)`$ for some positive function $`\rho `$. For real numbers $`a`$ and $`b`$, we use the notation $`ab=\mathrm{Max}(a,b)`$ and $`ab=\mathrm{Min}(a,b)`$. In the estimates of solutions of the relevant equations, we shall use the letter $`C`$ to denote constants, possibly different from an estimate to the next, depending on various parameters such as $`\gamma `$, but not on the solutions themselves or on their initial data. Those constants will be bounded in $`\gamma `$ for $`\gamma `$ away from zero. We shall use the notation $`A(a_1,a_2,\mathrm{})`$ for estimating functions, also possibly different from an estimate to the next, depending in addition on suitable norms $`a_1,a_2,\mathrm{}`$ of the solutions or of their initial data. If (p.q) is a double inequality, we denote by (p.qa) and (p.qb) the first and second inequality in (p.q). Finally, Item (p.q) of I or II will be referred to as Item (I.p.q.) or (II.p.q). Additional notation will be given when needed.
In all this paper, we assume that $`0<\mu n`$ and $`0<\gamma 1`$.
## 2 Heuristics
In this section, we describe in heuristic terms the construction of the modified wave operators for the equation (1.1). That construction is the same as that performed in II, and this section is mostly a summary of Section II.2, which we include in order to make this paper self-contained.
The problem that we address is that of classifying the possible asymptotic behaviours of the solutions of (1.1) by relating them to a set of model functions $`𝒱=\{v=v(u_+)\}`$ parametrized by some data $`u_+`$ and with suitably chosen and preferably simple asymptotic behaviour in time. For each $`v𝒱`$, one tries to construct a solution $`u`$ of (1.1) such that $`u(t)`$ behaves as $`v(t)`$ when $`t\mathrm{}`$ in a suitable sense. The map $`\mathrm{\Omega }:u_+u`$ thereby obtained classifies the asymptotic behaviours of solutions of (1.1) and is a preliminary version of the wave operator for positive time. A similar question can be asked for $`t\mathrm{}`$. From now on we restrict our attention to positive time.
In the short range case corresponding to $`\gamma >1`$ in (1.2), the previous scheme can be implemented by taking for $`𝒱`$ the set $`𝒱=\{v=U(t)u_+\}`$ of solutions of the equation
$$i_tv+\frac{1}{2}\mathrm{\Delta }v=0,$$
$`(2.1)`$
with $`U(t)`$ being the unitary group
$$U(t)=\mathrm{exp}\left(i(t/2)\mathrm{\Delta }\right).$$
$`(2.2)`$
The initial data $`u_+`$ for $`v`$ is called the asymptotic state for $`u`$.
In the long range case corresponding to $`\gamma 1`$ in (1.2), the previous set is known to be inadequate and has to be replaced by a better set of model functions obtained by modifying the previous ones by a suitable phase. The modification that we use requires additional structure of $`U(t)`$. In fact $`U(t)`$ can be written as
$$U(t)=M(t)D(t)FM(t)$$
$`(2.3)`$
where $`M(t)`$ is the operator of multiplication by the function
$$M(t)=\mathrm{exp}\left(ix^2/2t\right),$$
$`(2.4)`$
and $`D(t)`$ is the dilation operator defined by
$$\left(D(t)f\right)(x)=(it)^{n/2}f(x/t).$$
$`(2.5)`$
Let now $`\phi ^{(0)}=\phi ^{(0)}(x,t)`$ be a real function of space time and let $`z^{(0)}(x,t)=\mathrm{exp}(i\phi ^{(0)}(x,t))`$. We replace $`v(t)=U(t)u_+`$ by the modified free evolution
$$v(t)=M(t)D(t)z^{(0)}(t)w_+$$
$`(2.6)`$
where $`w_+=Fu_+`$. In order to allow for easy comparison of $`u`$ with $`v`$, it is then convenient to represent $`u`$ in terms of a phase factor $`z(t)=\mathrm{exp}(i\phi (t))`$ and of an amplitude $`w(t)`$ in such a way that asymptotically $`\phi (t)`$ behaves as $`\phi ^{(0)}(t)`$ and $`w(t)`$ tends to $`w_+`$. This is done by writing $`u`$ in the form
$$u(t)=M(t)D(t)z(t)w(t)\left(\mathrm{\Lambda }(w,\phi )\right)(t).$$
$`(2.7)`$
The construction of the wave operators for $`u`$ proceeds by first constructing the wave operators for the pair $`(w,\phi )`$ and then recovering the wave operators for $`u`$ therefrom by the use of (2.7). The evolution equation for $`(w,\phi )`$ is obtained by substituting (2.7) into the equation (1.1). One obtains the equation
$$\left(i_t+(2t^2)^1\mathrm{\Delta }D^{}\stackrel{~}{g}D\right)zw=0$$
$`(2.8)`$
for $`zw`$, with
$$\stackrel{~}{g}\stackrel{~}{g}\left(|u|^2\right)=\stackrel{~}{g}\left(|Dw|^2\right),$$
$`(2.9)`$
or equivalently, by expanding the derivatives in (2.8),
$$\left\{i_t+(2t^2)^1\mathrm{\Delta }i(2t^2)^1\left(2\phi +(\mathrm{\Delta }\phi )\right)\right\}w$$
$$+\left\{_t\phi (2t^2)^1|\phi |^2D^{}\stackrel{~}{g}D\right\}w=0.$$
$`(2.10)`$
We are now in the situation of a gauge theory. The equation (2.8) or (2.10) is invariant under the gauge transformation $`(w,\phi )(w\mathrm{exp}(i\omega ),\phi +\omega )`$, where $`\omega `$ is an arbitrary function of space time, and the original gauge invariant equation is not sufficient to provide evolution equations for the two gauge dependent quantities $`w`$ and $`\phi `$. At this point we arbitrarily add the Hamilton-Jacobi equation as a gauge condition. This yields a system of evolution equations for $`(w,\phi )`$, namely
$$\{\begin{array}{cc}_tw=i(2t^2)^1\mathrm{\Delta }w+(2t^2)^1\left(2\phi +(\mathrm{\Delta }\phi )\right)w\hfill & (2.11)\hfill \\ & \\ _t\phi =(2t^2)^1|\phi |^2+t^\gamma g_0(w,w)\hfill & (2.12)\hfill \end{array}$$
where we have defined
$$g_0(w_1,w_2)=\kappa \mathrm{Re}||^{\mu n}w_1\overline{w}_2$$
$`(2.13)`$
and rewritten the nonlinear interaction term in (2.10) as
$$D^{}\stackrel{~}{g}\left(|Dw|^2\right)D=t^\gamma g_0(w,w).$$
The gauge freedom in (2.11)-(2.12) is now reduced to that given by an arbitrary function of space only. It will be shown in Section 4 that the Cauchy problem for the system (2.11)-(2.12) is locally well-posed in a neighborhood of infinity in time. The solutions thereby obtained behave asymptotically as $`w(t)=O(1)`$ and $`\phi (t)O(t^{1\gamma })`$ as $`t\mathrm{}`$, a behaviour that is immediately seen to be compatible with (2.11)-(2.12).
We next study the asymptotic behaviour of the solutions of the auxiliary system (2.11)-(2.12) in more detail and try to construct wave operators for that system. For that purpose, we need to choose a set of model functions playing the role of $`v`$, in the spirit of (2.6). We proceed as follows. Let $`p0`$ be an integer. We write
$$\{\begin{array}{cc}w=\underset{0mp}{}w_m+q_{p+1}W_p+q_{p+1}\hfill & (2.14)\hfill \\ & \\ \phi =\underset{0mp}{}\phi _m+\psi _{p+1}\varphi _p+\psi _{p+1}\hfill & (2.15)\hfill \end{array}$$
with the understanding that asymptotically in $`t`$
$$w_m(t)=O\left(t^{m\gamma }\right),q_{p+1}(t)=o\left(t^{p\gamma }\right),$$
$`(2.16)`$
$$\phi _m(t)=O\left(t^{1(m+1)\gamma }\right),\psi _{p+1}(t)=o\left(t^{1(p+1)\gamma }\right).$$
$`(2.17)`$
Substituting (2.14)-(2.15) into (2.11)-(2.12) and identifying the various powers of $`t^\gamma `$ yields the following system of equations for $`(w_m,\phi _m)`$ :
$$\{\begin{array}{cc}_tw_{m+1}=\left(2t^2\right)^1\underset{0jm}{}\left(2\phi _j+(\mathrm{\Delta }\phi _j)\right)w_{mj}\hfill & (2.18)\hfill \\ & \\ _t\phi _{m+1}=\left(2t^2\right)^1\underset{0jm}{}\phi _j\phi _{mj}+t^\gamma \underset{0jm+1}{}g_0(w_j,w_{m+1j})\hfill & (2.19)\hfill \end{array}$$
for $`m+10`$. Here it is understood that $`w_j=0`$ and $`\phi _j=0`$ for $`j<0`$. We supplement that system with the initial conditions
$$\{\begin{array}{cc}w_0(\mathrm{})=w_+,w_m(\mathrm{})=0\text{for}m1\hfill & (2.20)\hfill \\ & \\ \phi _m(1)=0\text{for}0mp.\hfill & (2.21)\hfill \end{array}$$
The system (2.18)-(2.19) with the initial conditions (2.20)-(2.21) can be solved by successive integrations : knowing $`(w_j,\phi _j)`$ for $`0jm`$, one constructs successively $`w_{m+1}`$ by integrating (2.18) between $`t`$ and $`\mathrm{}`$, and then $`\phi _{m+1}`$ by integrating (2.19) between 1 and $`t`$.
If $`(p+1)\gamma <1`$, that method of resolution reproduces the asymptotic behaviour in time (2.16) (2.17) which was used in the first place to provide a heuristic derivation of the system (2.18)-(2.19). For sufficiently large $`p`$, $`\varphi _p`$ is a sufficiently good approximation for $`\phi `$ to ensure that $`\psi _{p+1}`$ has a limit as $`t\mathrm{}`$. In fact by comparing the system (2.18)-(2.19) with (2.11)-(2.12), one finds that $`_t\psi _{p+1}`$ is essentially of the same order in $`t`$ as $`_t\phi _{p+1}`$, namely $`_t\psi _{p+1}O(t^{(p+2)\gamma })`$, which is integrable at infinity for $`(p+2)\gamma >1`$. In this way every solution $`(w,\phi )`$ of the system (2.11)-(2.12) as obtained previously has asymptotic states consisting of $`w_+=\underset{t\mathrm{}}{lim}w(t)`$ and $`\psi _+=\underset{t\mathrm{}}{lim}\psi _{p+1}(t)`$.
Conversely, under the condition $`(p+2)\gamma >1`$, we shall be able to solve the system (2.11)-(2.12) by looking for solutions in the form (2.14)-(2.15) with the additional initial condition $`\psi _{p+1}(\mathrm{})=\psi _+`$, thereby getting a solution which is asymptotic to $`(W_p,\varphi _p+\psi _+)`$ with
$$wW_pO\left(t^{(p+1)\gamma }\right),\phi \varphi _p\psi _+O\left(t^{1(p+2)\gamma }\right).$$
$`(2.22)`$
This allows to define a map $`\mathrm{\Omega }_0:(w_+,\psi _+)(w,\phi )`$ which is essentially the wave operator for $`(w,\phi )`$.
We next discuss the gauge covariance properties of $`\mathrm{\Omega }_0`$. Two solutions $`(w,\phi )`$ and $`(w^{},\phi ^{})`$ of the system (2.11)-(2.12) will be said to be gauge equivalent if they give rise to the same $`u`$ through (2.7), namely if $`w\mathrm{exp}(i\phi )=w^{}\mathrm{exp}(i\phi ^{})`$. If $`(w,\phi )`$ and $`(w^{},\phi ^{})`$ are two gauge equivalent solutions, one can show easily that the difference $`\phi _{}=\phi ^{}\phi `$ has a limit $`\omega `$ when $`t\mathrm{}`$ and that $`w_+^{}=w_+\mathrm{exp}(i\omega )`$. Under that condition, it turns out that the phases $`\{\phi _j\}`$ and $`\varphi _p`$ (but not the amplitudes) obtained by solving (2.18)-(2.19) are gauge invariant, namely $`\phi _m=\phi _m^{}`$ for $`0mp`$ and therefore $`\varphi _p=\varphi _p^{}`$, so that $`\psi _+^{}=\psi _++\omega `$. It is then natural to define gauge equivalence of asymptotic states $`(w_+,\psi _+)`$ and $`(w_+^{},\psi _+^{})`$ by the condition $`w_+\mathrm{exp}(i\psi _+)=w_+^{}\mathrm{exp}(i\psi _+^{})`$ and the previous result can be rephrased as the statement that gauge equivalent solutions of (2.11)-(2.12) in $`(\mathrm{\Omega }_0)`$ have gauge equivalent asymptotic states. Conversely, we shall show that gauge equivalent asymptotic states have gauge equivalent images under $`\mathrm{\Omega }_0`$. Here however we meet with a technical problem coming from the construction of $`\mathrm{\Omega }_0`$ itself. For given $`(w_+,\psi _+)`$ we construct $`(w,\phi )`$ in practice as follows. We take a (large) finite time $`t_0`$ and we define a solution $`(w_{t_0},\phi _{t_0})`$ of the system (2.11)-(2.12) by imposing a suitable initial condition at $`t_0`$, depending on $`(w_+,\psi _+)`$, and solving the Cauchy problem with finite initial time. We then let $`t_0`$ tend to infinity and obtain $`(w,\phi )`$ as the limit of $`(w_{t_0},\phi _{t_0})`$. The simplest way to prove the gauge equivalence of two solutions $`(w,\phi )`$ and $`(w^{},\phi ^{})`$ obtained in this way from gauge equivalent $`(w_+,\psi _+)`$ and $`(w_+^{},\psi _+^{})`$ consists in using an initial condition at $`t_0`$ which already ensures that $`(w_{t_0},\phi _{t_0})`$ and $`(w_{t_0}^{},\phi _{t_0}^{})`$ are gauge equivalent. However the natural choice $`(w_{t_0}(t_0),\phi _{t_0}(t_0))=(W_p(t_0),\varphi _p(t_0)+\psi _+)`$ does not satisfy that requirement as soon as $`p1`$ because $`\varphi _p(t_0)`$ is gauge invariant while $`W_p(t_0)\mathrm{exp}(\psi _+)`$ is not. In order to overcome that difficulty, we introduce a new amplitude $`V`$ and a new phase $`\chi `$ defined by solving the transport equations
$$\{\begin{array}{cc}_tV=(2t^2)^1\left(2\varphi _{p1}+\left(\mathrm{\Delta }\varphi _{p1}\right)\right)V\hfill & (2.23)\hfill \\ & \\ _t\chi =t^2\varphi _{p1}\chi \hfill & (2.24)\hfill \end{array}$$
with initial condition
$$V(\mathrm{})=w_+,\chi (\mathrm{})=\psi _+.$$
$`(2.25)`$
It follows from (2.23) (2.24) that $`V\mathrm{exp}(i\chi )`$ satisfies the same transport equation as $`V`$, now with gauge invariant initial condition $`(V\mathrm{exp}(i\chi ))(\mathrm{})=w_+\mathrm{exp}(i\psi _+)`$, and is therefore gauge invariant. Furthermore, $`(V,\chi )`$ is a sufficiently good approximation of $`(W_p,\psi _+)`$ in the sense that
$$V(t)W_p(t)O\left(t^{(p+1)\gamma }\right),\chi (t)\psi _+O(t^\gamma ).$$
$`(2.26)`$
One then takes $`(w_{t_0}(t_0),\phi _{t_0}(t_0))=(V(t_0),\varphi _p(t_0)+\chi (t_0))`$ as an initial condition at time $`t_0`$, thereby ensuring that $`(w_{t_0},\phi _{t_0})`$ and $`(w_{t_0}^{},\phi _{t_0}^{})`$ are gauge equivalent. That equivalence is easily seen to be preserved in the limit $`t_0\mathrm{}`$. Furthermore, the estimates (2.26) ensure that the asymptotic properties (2.22) are preserved by the modified construction. As a consequence of the previous discussion, the map $`\mathrm{\Omega }_0`$ is gauge covariant, namely induces an injective map of gauge equivalence classes of asymptotic states $`(w_+,\psi _+)`$ to gauge equivalence classes of solutions $`(w,\phi )`$ of the system (2.11)-(2.12).
The wave operator for $`u`$ is obtained from $`\mathrm{\Omega }_0`$ just defined and from $`\mathrm{\Lambda }`$ defined by (2.7). From the previous discussion it follows that the map $`\mathrm{\Lambda }\mathrm{\Omega }_0:(w_+,\psi _+)u`$ is injective from gauge equivalence classes of asymptotic states $`(w_+,\psi _+)`$ to solutions of (1.1). In order to define a wave operator for $`u`$ involving only the asymptotic state $`u_+`$ but not an arbitrary phase $`\psi _+`$, we choose a representative in each equivalence class $`(w_+,\psi _+)`$, namely we define the wave operator for $`u`$ as the map $`\mathrm{\Omega }:u_+u=(\mathrm{\Lambda }\mathrm{\Omega }_0)(Fu_+,0)`$. Since each equivalence class of asymptotic states contains at most one element with $`\psi _+=0`$, the map $`\mathrm{\Omega }`$ is again injective.
The previous heuristic discussion was based in part on a number of time decay estimates in terms of negative powers of $`t`$. In practice however two complications occur, namely (i) for integer $`\gamma ^1`$, some of the estimates involve logarithmic factors in time, and (ii) the use of Gevrey spaces requires that of norms defined by integrals over time involving a convergence factor which eventually produces a small loss in the time decays. Both difficulties are handled by introducing suitable estimating functions of time, some of which are defined by integral representations and generalize in a natural way a similar family of functions defined in II.
In the same way as in I, the system (2.11)-(2.12) can be rewritten as a system of equations for $`w`$ and for $`s=\phi `$, from which $`\phi `$ can then be recovered by (2.12), thereby leading to a slightly more general theory since the system for $`(w,s)`$ can be studied without even assuming that $`s`$ is a gradient. For simplicity, and in the same way as in II, we shall not follow that track. However, we shall use systematically the notation $`s=\phi `$, and for the purposes of estimation, we shall supplement the system (2.11)-(2.12) with the equation satisfied by $`s`$, which is simply the gradient of (2.12), namely
$$_ts=t^2ss+t^\gamma g_0(w,w).$$
$`(2.27)`$
We are now in a position to describe in more detail the contents of the technical sections 3-7 of this paper. In Section 3, we introduce the relevant Gevrey spaces and derive the basic estimates in those spaces that are needed to study the system (2.11)-(2.12) (Lemmas 3.4-3.7), we explain in passing the mechanism by which those spaces make it possible to overcome the derivative loss in (2.11)-(2.12) for $`\mu >n2`$ (after Lemma 3.4) and finally we introduce the estimating functions of time mentioned above and obtain some estimates for them. In Section 4, we prove that the Cauchy problem for the system (2.11)-(2.12) is well-posed for large time, with large but finite initial time (Proposition 4.1), we prove the existence of a limit for $`w(t)`$ as $`t\mathrm{}`$ for the solutions thereby obtained (Proposition 4.2) and we derive a uniqueness result of solutions with prescribed asymptotic behaviour (Proposition 4.3). In Section 5, we study the asymptotic behaviour in time of the solutions obtained in Section 4. We derive a number of properties and estimates for the solutions of the asymptotic system (2.18)-(2.19), defined inductively (Proposition 5.2). We then obtain asymptotic estimates on the approximation of the solutions of the system (2.11)-(2.12) by the asymptotic functions $`(W_m,\varphi _m)`$ defined by (2.14)-(2.15), and in particular we complete the proof of the existence of asymptotic states for those solutions (Proposition 5.3). In Section 6, we study the Cauchy problem with infinite initial time, first for the transport equations (2.23) (Proposition 6.1) and (2.24) (Proposition 6.2), and then for the system (2.11)-(2.12). For a given solution $`(V,\chi )`$ of the system (2.23)-(2.24) and a given (large) $`t_0`$, we construct a solution $`(w_{t_0},\phi _{t_0})`$ of the system (2.11)-(2.12) which coincides with $`(V,\varphi _p+\chi )`$ at $`t_0`$ and we estimate it uniformly in $`t_0`$ (Proposition 6.4). We then prove that when $`t_0\mathrm{}`$, $`(w_{t_0},\phi _{t_0})`$ has a limit $`(w,\phi )`$ which is asymptotic both to $`(V,\varphi _p+\chi )`$ and to $`(W_p,\varphi _p+\psi _+)`$ (Proposition 6.5). In Section 7, we exploit the results of Section 6 to construct the wave operators for the equation (1.1) and to describe the asymptotic behaviour of solutions in their range. We first prove that the local wave operator at infinity for the system (2.11)-(2.12) defined through Proposition 6.5 in Definition 7.1 is gauge covariant in the sense of Definitions 7.2 and 7.3 in the best form that can be expected with the available regularity (Propositions 7.2 and 7.3). With the help of some information on the Cauchy problem for (1.1) at finite time (Proposition 7.1), we then define the wave operator $`\mathrm{\Omega }:u_+u`$ (Definition 7.4), and we prove that it is injective. We then collect all the available information on $`\mathrm{\Omega }`$ and on solutions of (1.1) in its range in Proposition 7.5, which contains the main results of this paper. Finally some side results relevant for the definition and properties of the Gevrey spaces used here are collected in two Appendices.
## 3 Gevrey spaces and preliminary estimates
In this section, we define the Gevrey spaces where we shall study the auxiliary system (2.11)-(2.12) and we derive a number of energy type estimates which hold in those spaces and play an essential role in that study. We then introduce some estimating functions of time generalizing those of II and we derive a number of estimates for them.
The relevant spaces will be defined with the help of the functions
$$f_0(\xi )=\mathrm{exp}\left(\rho |\xi |^\nu \right),f(\xi )=\mathrm{exp}\left(\rho (|\xi |^\nu 1)\right)$$
$`(3.1)`$
where $`0<\nu 1`$, $`\rho `$ is a positive parameter to be specified later, and $`\xi IR^n`$. The dependence of $`f`$ on $`\rho `$ will always be omitted in the notation.
In all this paper, one could use instead of $`f_0`$ the function
$$\stackrel{~}{f}(\xi )=\underset{j0}{}(j!)^{1/\nu }\rho ^{j/\nu }|\xi |^j$$
$`(3.2)`$
which satisfies the same basic estimates and would yield essentially the same results. The function $`\stackrel{~}{f}`$ is also convenient in order to relate the definition of the Gevrey spaces $`K_\rho ^k`$ and $`Y_\rho ^{\mathrm{}}`$ (see (3.8) (3.9) below) to more standard definitions. Those points are discussed in Appendix A.
The functions $`f_0`$ and $`f`$ satisfy the following estimates.
Lemma 3.1. Let $`\xi `$, $`\eta IR^n`$. Then $`f`$ satisfies the estimates :
$$f(\xi )f(\xi \eta )f(\eta )\text{for all }\xi \text{ , }\eta ,$$
$`(3.3)`$
$$f(\xi )f(\xi \eta )f_0(\eta )^\nu \text{for}|\xi ||\eta ||\xi \eta |,$$
$`(3.4)`$
$$|f(\xi )f(\eta )||\eta |^{1\nu }|\xi \eta |^{1\nu }f(\xi \eta )f(\eta )\text{for all }\xi \text{}\eta ,$$
$`(3.5)`$
$$|f(\xi )f(\eta )||\eta |^{1\nu }C|\xi \eta |^{1\nu }f_0(\xi \eta )^\nu f(\eta )\text{for}|\xi ||\xi \eta ||\eta |,$$
$`(3.6)`$
$$|f(\xi )f(\eta )||\eta |^{1\nu }C|\xi \eta |^{1\nu }f(\xi \eta )f_0(\eta )^\nu \text{for}|\xi ||\eta ||\xi \eta |.$$
$`(3.7)`$
In (3.6) and (3.7), one can take $`C=1`$, except in the region $`|\xi ||\xi \eta ||\eta |`$ where $`C=2^{1\nu }`$.
The function $`f_0`$ satisfies the same estimates as $`f`$.
Proof. We first prove the estimates for $`f_0`$
(3.3) follows from the fact that $`|\xi |^\nu |\xi \eta |^\nu +|\eta |^\nu `$.
(3.4) is obvious for $`|\xi ||\xi \eta |`$. For $`|\eta ||\xi \eta ||\xi |`$, we estimate
$$|\xi |^\nu |\xi \eta |^\nu +\nu |\xi \eta |^{\nu 1}|\eta ||\xi \eta |^\nu +\nu |\eta |^\nu .$$
(3.5) follows from (3.3) for $`|\eta ||\xi \eta |`$ and from (3.6) with $`C=1`$ for $`|\xi \eta ||\xi ||\eta |`$. For $`|\xi ||\xi \eta ||\eta |`$, we estimate
$$\left(f_0(\eta )f_0(\xi )\right)|\eta |^{1\nu }f_0(\eta )|\eta |^{1\nu }\left(1\mathrm{exp}(\rho |\eta |^\nu )\right)f_0(\eta )\rho |\eta |$$
$$f_0(\eta )\left(|\eta |/|\xi \eta |\right)|\xi \eta |^{1\nu }\rho |\xi \eta |^\nu $$
and the result follows from the fact that $`|\eta |2|\xi \eta |`$ and $`\rho |\xi \eta |^\nu e^1f_0(\xi \eta )`$.
(3.6) is obvious for $`|\xi ||\eta ||\xi \eta |`$, with $`C=1`$.
For $`|\xi ||\xi \eta ||\eta |`$, we estimate
$$\left(f_0(\eta )f_0(\xi )\right)|\eta |^{1\nu }2^{1\nu }|\xi \eta |^{1\nu }f_0(\eta )$$
which yields (3.6) with $`C=2^{1\nu }`$ since $`|\eta |2|\xi \eta |`$.
The really crucial case is the case $`|\xi \eta ||\xi ||\eta |`$.
For $`|\xi \eta ||\xi ||\eta |`$, we estimate
$$\left(f_0(\eta )f_0(\xi )\right)|\eta |^{1\nu }f_0(\eta )|\eta |^{1\nu }\left(1\mathrm{exp}(\rho \nu |\xi \eta ||\xi |^{\nu 1})\right)$$
where we have used
$$|\eta |^\nu |\xi |^\nu \nu |\xi |^{\nu 1}|\xi \eta |,$$
$$\mathrm{}f_0(\eta )\left(|\eta |/|\xi |\right)^{1\nu }\rho \nu |\xi \eta |$$
$$f_0(\eta )|\xi \eta |^{1\nu }2^{1\nu }e^1f_0(\xi \eta )^\nu $$
since $`|\eta |2|\xi |`$ and $`\rho \nu |\xi \eta |^\nu e^1f_0(\xi \eta )^\nu `$. This proves (3.6) with $`C=1`$ in that case.
For $`|\xi \eta ||\eta ||\xi |`$, we estimate similarly
$$\left(f_0(\xi )f_0(\eta )\right)|\eta |^{1\nu }f_0(\eta )\left\{|\eta |^{1\nu }\left(\mathrm{exp}(\rho \nu |\xi \eta ||\eta |^{\nu 1})1\right)\right\}.$$
Now for fixed $`|\xi \eta |`$, the last bracket is a decreasing function of $`|\eta |`$, and is therefore bounded by its value for $`|\eta |=|\xi \eta |`$, so that
$$\mathrm{}f_0(\eta )|\xi \eta |^{1\nu }\left(f_0(\xi \eta )^\nu 1\right)$$
which proves (3.6) with $`C=1`$ in that case.
(3.7) is obvious for $`|\xi ||\eta ||\xi \eta |`$ and follows from (3.4) with $`C=1`$ for $`|\eta ||\xi ||\xi \eta |`$. For $`|\xi ||\xi \eta ||\eta |`$, we estimate
$$\left(f_0(\eta )f_0(\xi )\right)|\eta |^{1\nu }2^{1\nu }|\xi \eta |^{1\nu }f_0(\eta )$$
and (3.7) with $`C=2^{1\nu }`$ follows from
$$|\eta |^\nu \nu |\eta |^\nu +(1\nu )2^\nu |\xi \eta |^\nu \nu |\eta |^\nu +|\xi \eta |^\nu .$$
The estimates for $`f`$ follow from those for $`f_0`$. This is obvious for (3.3) (3.4). For (3.5) (3.6) (3.7), it follows from the fact that for all $`\xi `$, $`\eta `$ and all $`a>0`$
$$|f_0(\xi )af_0(\eta )a||f_0(\xi )f_0(\eta )|.$$
$``$$``$
We now turn to the definition of the spaces where we shall solve the system (2.11)-(2.12). For any tempered distribution $`u`$ in $`IR^n`$ with $`\widehat{u}L_{loc}^1(IR^n)`$, we define $`u_{\genfrac{}{}{0pt}{}{>}{<}}`$ by $`\widehat{u}_>(\xi )=\widehat{u}(\xi )`$ for $`|\xi |>1`$, $`\widehat{u}_>(\xi )=0`$ for $`|\xi |1`$, $`\widehat{u}_<(\xi )=0`$ for $`|\xi |>1`$, $`\widehat{u}_<(\xi )=\widehat{u}(\xi )`$ for $`|\xi |1`$. Similarly, for $`\xi IR^n`$ and $`mIR`$, we define $`|\xi |_>^m`$ and $`|\xi |_<^m`$ to be equal to $`|\xi |^m`$ for $`|\xi |>1`$ and $`|\xi |1`$ respectively, and zero otherwise. Occasionally we shall make the separation between low and high $`|\xi |`$ at some value $`a1`$. In that case we shall denote by $`u_{<a}`$ and $`u_{>a}`$ the corresponding components of $`u`$.
Let now $`\rho >0`$, $`kIR`$, $`\mathrm{}IR`$ and $`0\mathrm{}_<<n/2`$. Starting from Lemma 3.4 below (but not until then), we shall assume in addition that $`\mathrm{}_<>n/2\mu `$. We define
$$K_\rho ^k=\{w:|w|_k^2w;K_\rho ^k^2|\xi |^kf(\xi )\widehat{w}_>(\xi )_2^2+f(\xi )\widehat{w}_<(\xi )_2^2<\mathrm{}\},$$
$`(3.8)`$
$$Y_\rho ^{\mathrm{}}=\{\phi :\widehat{\phi }L_{loc}^1(IR^n)\text{and}|\phi |_{\mathrm{}}^2\phi ;Y_\rho ^{\mathrm{}}^2|\xi |^{\mathrm{}+2}f(\xi )\widehat{\phi }_>(\xi )_2^2$$
$$+|\xi |^\mathrm{}_<f(\xi )\widehat{\phi }_<(\xi )_2^2<\mathrm{}\}.$$
$`(3.9)`$
The apparent ambiguity in the notation $`||_b`$ will be lifted by the fact that the symbol $`b`$ will always contain the letter $`k`$ when referring to $`K_\rho ^k`$ spaces and the letter $`\mathrm{}`$ when referring to $`Y_\rho ^{\mathrm{}}`$ spaces. The spaces $`K_\rho ^k`$ and $`Y_\rho ^{\mathrm{}}`$ are Hilbert spaces and satisfy the embeddings $`K_\rho ^kK_\rho ^k^{}`$ for $`kk^{}`$ and $`Y_\rho ^{\mathrm{}}Y_\rho ^{\mathrm{}^{}}`$ for $`\mathrm{}\mathrm{}^{}`$, $`K_\rho ^kK_\rho ^{}^k`$ and $`Y_\rho ^{\mathrm{}}Y_\rho ^{}^{\mathrm{}}`$, for $`\rho \rho ^{}`$.
Remark 3.1. The norms in the spaces $`K_\rho ^k`$ and $`Y_\rho ^{\mathrm{}}`$ are both of the form $`f_1f\widehat{u}_2`$ with
$$f_1=|\xi |_>^{k_>}+|\xi |_<^{k_<}$$
$`(3.10)`$
where $`(k_>,k_<)=(k,0)`$ for $`K_\rho ^k`$ and $`(k_>,k_<)=(\mathrm{}+2,\mathrm{}_<)`$ for $`Y_\rho ^{\mathrm{}}`$. In particular this implies that
$$\{\begin{array}{c}w;K_\rho ^k=F^1(f\widehat{w});K_0^k\hfill \\ \\ \phi ;Y_\rho ^{\mathrm{}}=F^1(f\widehat{\phi });Y_0^{\mathrm{}}.\hfill \end{array}$$
$`(3.11)`$
If $`k_>k_<`$, namely if $`k0`$ for $`K_\rho ^k`$ and if $`\mathrm{}+2\mathrm{}_<`$ for $`Y_\rho ^{\mathrm{}}`$ (which will always be the case in the applications), one can omit either or both of the upper and lower restrictions in the definition of $`K_\rho ^k`$ or $`Y_\rho ^{\mathrm{}}`$, thereby obtaining equivalent norms uniformly in $`\rho `$. In fact, in that case
$$|\xi |_>^{2k_>}+|\xi |_<^{2k_<}|\xi |^{2k_>}+|\xi |^{2k_<}2\left(|\xi |_>^{2k_>}+|\xi |_<^{2k_<}\right).$$
Furthermore, the relations (3.11) are preserved under thoses changes.
We shall use the spaces $`K_\rho ^k`$ and $`Y_\rho ^{\mathrm{}}`$ with a time dependent parameter $`\rho 𝒞^1(IR^+,IR^+)`$. The form (3.8)-(3.9) of the norms has been chosen so as to ensure that for fixed (time independent) $`w`$ and $`\phi `$, the following relations hold
$$\{\begin{array}{c}\frac{d}{dt}|w|_k^2=2\rho ^{}|w|_{k+\nu /2}^2,\hfill \\ \\ \frac{d}{dt}|\phi |_{\mathrm{}}^2=2\rho ^{}|\phi |_{\mathrm{}+\nu /2}^2,\hfill \end{array}$$
$`(3.12)`$
where $`\rho ^{}=d\rho /dt`$.
We shall look for $`w`$ as complex $`K_\rho ^k`$ valued functions of time and for $`\phi `$ as real $`Y_\rho ^{\mathrm{}}`$ valued functions of time. More precisely, we shall look for $`(w,\phi )`$ such that for some interval $`I[1,\mathrm{})`$
$$(w,\phi )𝒳_{\rho ,loc}^{k,\mathrm{}}(I)𝒞(I,K_\rho ^kY_\rho ^{\mathrm{}})L_{loc}^2(I,K_\rho ^{k+\nu /2}Y_\rho ^{\mathrm{}+\nu /2})$$
$`(3.13)`$
by which is meant, especially as regards continuity, that
$$(F^1f\widehat{w},F^1f\widehat{\phi })𝒳_{0,loc}^{k,\mathrm{}}(I)𝒞(I,K_0^kY_0^{\mathrm{}})L_{loc}^2(I,K_0^{k+\nu /2}Y_0^{\mathrm{}+\nu /2})$$
in the usual sense. In particular, when taking norms such as $`|w(t)|_k`$ or $`|\phi (t)|_{\mathrm{}}`$ with time dependent $`\rho `$, it will always be understood that $`\rho `$ in the definition of the relevant space is taken at the same value of the time as $`w`$ or $`\phi `$.
We shall also need global versions of the previous spaces, especially when the interval $`I`$ is unbounded. The definition of those global versions will require assumptions on $`\rho `$ that are irrelevant for the considerations of this section and will be postponed until the beginning of Section 4.
We shall need the following elementary estimates.
Lemma 3.2. Let $`mIR`$. The following estimates hold :
$$|\xi |^m\widehat{(u_1u_2)}_>_2C<\xi >^{k_1}\widehat{u}_1_2<\xi >^{k_2}\widehat{u}_2_2$$
$`(3.14)`$
for $`k_1`$, $`k_2m0`$ and $`k_1+k_2>m+n/2`$, where $`<>=(1+||^2)^{1/2}`$,
$$|\xi |^m\widehat{(u_1u_2)}_<_2Cu_1_2u_2_2$$
$`(3.15)`$
for $`m>n/2`$.
Proof. (3.14) for $`m0`$ follows from (3.14) (3.15) for $`m=0`$ and from
$$|\xi |^m\widehat{(u_1u_2)}_2C\left((|\xi |^m|\widehat{u}_1|)|\widehat{u}_2|_2+|\widehat{u}_1|(|\xi |^m|\widehat{u}_2|)_2\right).$$
For $`m<0`$, we estimate by the Hölder and Young inequalities
$$|\xi |^m\widehat{(u_1u_2)}_>_2C|\xi |_>^m_s\widehat{u}_1_{\overline{r}_1}\widehat{u}_2_{\overline{r}_2}$$
with $`1/s+1/\overline{r}_1+1/\overline{r}_2=3/2`$ , $`n/s<|m|`$ and $`1\overline{r}_1`$, $`\overline{r}_22`$. The last two norms are estimated by $`<\xi >^{k_i}\widehat{u}_i_2`$ provided $`k_i>n/\overline{r}_in/2`$. One can find $`\overline{r}_1`$, $`\overline{r}_2`$ satisfying all previous conditions under the assumptions made on $`k_i`$.
For $`m=0`$, we apply the same argument with $`s=\mathrm{}`$.
For the proof of (3.15) we estimate by the Schwarz and Young inequalities
$$|\xi |^m\widehat{(u_1u_2)}_<_2|\xi |_<^m_2u_1_2u_2_2$$
for $`m>n/2`$.
$``$$``$
For future reference, we also state the following elementary inequalities
$$|\xi |^m\widehat{\phi }_<_1C|\xi |^\mathrm{}_<\widehat{\phi }_<_2\text{for all}m0,$$
$`(3.16)`$
$$|\xi |^m\widehat{\phi }_>_1C|\xi |^{\mathrm{}+2}\widehat{\phi }_>_2\text{for }\mathrm{}+2>m+n/2.$$
$`(3.17)`$
In what follows, we shall repeatedly estimate norms such as $`|\xi |^mf\widehat{u_1u_2}_2`$ with $`m0`$. For that purpose, using (3.3), we shall write
$$|\xi |^mf|\widehat{u_1u_2}||\xi |^m𝑑\eta f(\xi )|\widehat{u}_1(\xi \eta )||\widehat{u}_2(\eta )|$$
$$2^m𝑑\eta \left(|\xi \eta |^m+|\eta |^m\right)f(\xi \eta )f(\eta )|\widehat{u}_1(\xi \eta )||\widehat{u}_2(\eta )|$$
$$=2^m\{(||^mf|\widehat{u}_1|)(f|\widehat{u}_2|)+(12)\}.$$
$`(3.18)`$
That inequality will be often combined with restrictions to low or high values of $`|\xi |`$, either in the product $`(u_1u_2)`$ or in $`u_1`$ or $`u_2`$ separately.
The next lemma states that under suitable assumptions on $`k`$ and $`\mathrm{}`$, $`Y_\rho ^{\mathrm{}}`$ is an algebra under ordinary multiplication and acts boundedly on $`K_\rho ^k`$ by multiplication.
Lemma 3.3. Let $`\mathrm{}+2>n/2`$ and $`0k\mathrm{}+2`$. Then there exist constants $`C_1`$ and $`C_2`$, independent of $`\rho `$, such that
$$|\phi \psi |_{\mathrm{}}C_1|\phi |_{\mathrm{}}|\psi |_{\mathrm{}}\text{for all }\phi \text{,}\psi Y_\rho ^{\mathrm{}},$$
$`(3.19)`$
$$|\phi w|_kC_2|\phi |_{\mathrm{}}|w|_k\text{for all }\phi Y_\rho ^{\mathrm{}}\text{,}wK_\rho ^k.$$
$`(3.20)`$
In particular
$$|(\mathrm{exp}(i\phi )1)w|_kC_2C_1^1\left(\mathrm{exp}(C_1|\phi |_{\mathrm{}})1\right)|w|_k$$
$`(3.21)`$
for all $`\phi Y_\rho ^{\mathrm{}}`$, $`wK_\rho ^k`$.
Proof. From the definitions (3.8) (3.9) of the norms and in particular from (3.11), and from (3.3), it follows that (3.19) (3.20) for $`\rho =0`$ imply the same estimates for arbitrary $`\rho >0`$ with the same constants. We therefore restrict our attention to the case $`\rho =0`$. In that case (3.19) (3.20) are almost standard properties of Sobolev spaces, except for the presence of $`|\xi |_<^\mathrm{}_<`$ with possibly $`\mathrm{}_<>0`$ in (3.9). We give a proof for completeness.
(3.19). From the definition (3.9) with $`\rho =0`$, it follows that it is sufficient to estimate $`|\xi |^{\mathrm{}+2}(\widehat{\phi \psi })_{>2}_2`$ and $`|\xi |^\mathrm{}_<\widehat{\phi \psi }_2`$. We estimate
$$|\xi |^{\mathrm{}+2}(\widehat{\phi \psi })_{>2}_2C\{(|\xi |^{\mathrm{}+2}|\widehat{\phi }_>|)|\widehat{\psi }|_2$$
$$+|\widehat{\phi }_<||\widehat{\psi }_>|_2+(\phi \psi )\}$$
$$C\{|\xi |^{\mathrm{}+2}\widehat{\phi }_>_2\widehat{\psi }_1+\widehat{\phi }_<_1\widehat{\psi }_>_2+(\phi \psi )\}$$
$$C|\phi |_{\mathrm{}}|\psi |_{\mathrm{}}$$
by (3.3), by the Young inequality and by (3.16) (3.17). The lower restriction $`|\xi |>2`$ in $`\phi \psi `$ implies that there is no $`\phi _<\psi _<`$ contribution.
On the other hand
$$|\xi |^\mathrm{}<\widehat{\phi \psi }_2C\{|\xi |^\mathrm{}_<\widehat{\phi }_2\widehat{\psi }_1+(\phi \psi )\}C|\phi |_{\mathrm{}}|\psi |_{\mathrm{}}$$
by the Young inequality and (3.16) (3.17).
(3.20). We estimate similarly by (3.3) and the Young inequality
$$<\xi >^k\widehat{\phi w}_2C\{\widehat{\phi }_1<\xi >^k\widehat{w}_2+\widehat{\phi }_<_1\widehat{w}_2$$
$$+(|\xi |^k|\widehat{\phi }_>|)|\widehat{w}|_2\}$$
and the result follows from (3.16) (3.17) and from Lemma 3.2 with $`m=0`$, $`k_1=\mathrm{}+2k`$ and $`k_2=k`$.
(3.21) follows immediately from a repeated application of (3.19) (3.20).
$``$$``$
Remark 3.2. In Lemma 3.3 we have used only (3.3) from Lemma 3.1. For $`\nu <1`$, by using (3.4), one can obtain more general results. In particular $`K_\rho ^k`$ and $`Y_\rho ^{\mathrm{}}`$ are algebras under multiplication for any $`k`$ and $`\mathrm{}`$ in $`IR^+`$ (see Appendix B). Similarly in what follows, we shall use only (3.3) and (3.5). For $`\nu <1`$, by using in addition (3.4) and (3.6), one could generalize some of the results by weakening the assumptions made on $`k`$ and $`\mathrm{}`$. That extension however would not hold uniformly in $`\nu `$ for $`\nu 1`$, and we shall therefore refrain from following that track.
We now turn to the derivation of the basic estimates needed to study the auxiliary system (2.11)-(2.12). For that purpose, we shall use a regularization. Let $`j𝒞^1(IR^n,IR)`$ with $`0j1`$ and $`j(0)=1`$. We denote by $`j_\epsilon `$ both the function $`j_\epsilon (\xi )=j(\epsilon \xi )`$ and the operator of multiplication by that function in Fourier space variables, and by $`J_\epsilon `$ the operator $`F^{}j_\epsilon F`$.
We recall that $`g_0`$ is defined by (2.13). We shall use systematically the notation $`s=\phi `$ and whenever convenient, start from the equation (2.27) satisfied by $`s`$ instead of the equation (2.12) satisfied by $`\phi `$.
We first state the basic estimates.
Lemma 3.4. Let $`m0`$. The following estimates hold :
$$|\mathrm{Re}<j_\epsilon |\xi |^kf\widehat{w},j_\epsilon |\xi |^kf(\widehat{sw})>|+|\mathrm{Re}<j_\epsilon f\widehat{w},j_\epsilon f(\widehat{sw})>|C|w|_{k+\nu /2}^2|\phi |_{\mathrm{}}$$
$`(3.22)`$
uniformly in $`\epsilon `$, for $`\mathrm{}>n/2\nu `$ , $`k\nu /2`$ , $`\mathrm{}+1k\nu /2`$.
$$|\xi |^mf(\widehat{sw})_2C|\phi |_{\mathrm{}}|w|_k$$
$`(3.23)`$
for $`k+\mathrm{}>m+n/2`$ , $`km+1`$ , $`\mathrm{}+1m`$.
$$|\xi |^mf((\widehat{s)w})_2C|\phi |_{\mathrm{}}|w|_k$$
$`(3.24)`$
for $`k+\mathrm{}>m+n/2`$ , $`km`$ , $`\mathrm{}m`$.
$$|\mathrm{Re}<j_\epsilon |\xi |^{\mathrm{}^{}+1}f\widehat{s_>^{}},j_\epsilon |\xi |^{\mathrm{}^{}+1}f(\widehat{ss^{}})_>>|C|\phi ^{}|_{\mathrm{}^{}+\nu /2}^2|\phi |_{\mathrm{}}$$
$`(3.25)`$
uniformly in $`\epsilon `$, for $`\mathrm{}>n/2\nu `$ , $`\mathrm{}^{}+1\nu /2`$ , $`\mathrm{}\mathrm{}^{}\nu /2`$.
$$|\xi |^mf(\widehat{ss^{}})_>_2C|\phi |_{\mathrm{}}|\phi ^{}|_{\mathrm{}^{}}$$
$`(3.26)`$
for $`\mathrm{}+\mathrm{}^{}>m1+n/2`$ , $`\mathrm{}+1m`$ , $`\mathrm{}^{}m`$.
$$|\xi |^\mathrm{}_<f(\widehat{ss^{}})_2C|\phi |_{\mathrm{}}|\phi ^{}|_{\mathrm{}^{}}$$
$`(3.27)`$
for $`\mathrm{}`$, $`\mathrm{}^{}>n/22`$.
$$|\xi |^mf\widehat{g_0(w_1}w_2)_>_2C\{|w_1|_{k_1}|w_2|_{k_2}+|w_1|_{k_1^{}}|w_2|_{k_2^{}}\}$$
$`(3.28)`$
for $`(k_1+k_2)(k_1^{}+k_2^{})>\beta +n/2`$ , $`k_1k_2^{}\beta 0`$ , $`k_2k_1^{}0`$, where $`\beta =m+\mu n`$.
$$|\xi |^\mathrm{}_<\widehat{g_0(w_1}w_2)_<_2Cw_1_2w_2_2$$
$`(3.29)`$
for $`\mathrm{}_<>n/2\mu `$.
The constants $`C`$ in (3.22)-(3.29) can be taken independent of $`\rho `$.
Proof. (3.22). We have to estimate
$$\mathrm{Im}𝑑\xi 𝑑\eta j_\epsilon (\xi )|\xi |^kf(\xi )\overline{\widehat{w}}(\xi )j_\epsilon (\xi )|\xi |^kf(\xi )\widehat{s}(\xi \eta )\eta \widehat{w}(\eta )$$
$`(3.30)`$
and a similar expression with $`k=0`$. We consider only the former one. The proof for the latter is similar and simpler. We split the domain of integration into three regions, namely
$$|\xi \eta ||\xi ||\eta |,|\xi ||\xi \eta ||\eta |\text{and}|\eta ||\xi ||\xi \eta |$$
and correspondingly the integral (3.30) is written as the sum $`I_1+I_2+I_3`$ of three terms which we estimate successively.
Region $`|\xi \eta ||\xi ||\eta |`$, estimate of $`I_1`$.
In this region we decompose the integrand according to the identity
$$j_\epsilon (\xi )|\xi |^kf(\xi )\eta =j_\epsilon (\xi )|\xi |^k\left(f(\xi )f(\eta )\right)\eta +\left(j_\epsilon (\xi )|\xi |^kj_\epsilon (\eta )|\eta |^k\right)f(\eta )\eta +j_\epsilon (\eta )|\eta |^kf(\eta )\eta $$
$`(3.31)`$
and correspondingly $`I_1`$ is written as the sum $`I_{1,1}+I_{1,2}+I_{1,3}`$ of three terms which we estimate successively.
Estimate of $`I_{1,1}`$. From (3.5) of Lemma 3.1, we obtain
$$\left|j_\epsilon (\xi )|\xi |^{2k}(f(\xi )f(\eta ))\eta \right|C|\xi |^{k+\nu /2}|\xi \eta |^{1\nu }|\eta |^{k+\nu /2}f(\xi \eta )f(\eta )$$
and therefore by the Schwarz and Young inequalities
$$|I_{1,1}|C|w|_{k+\nu /2}^2|\xi |^{1\nu }f\widehat{s}_1$$
so that by (3.16) (3.17)
$$|I_{1,1}|C|w|_{k+\nu /2}^2|\phi |_{\mathrm{}}.$$
$`(3.32)`$
Estimate of $`I_{1,2}`$. We rewrite
$$\left(j_\epsilon (\xi )|\xi |^kj_\epsilon (\eta )|\eta |^k\right)\eta =j_\epsilon (\xi )\left(|\xi |^k|\eta |^k\right)\eta +|\eta |^k\left(j_\epsilon (\xi )j_\epsilon (\eta )\right)\eta .$$
$`(3.33)`$
We estimate
$$\left||\xi |^k|\eta |^k\right||\eta |k2^{|k1|}|\xi \eta ||\eta |^k$$
$`(3.34)`$
and we rewrite
$$\left(j_\epsilon (\xi )j_\epsilon (\eta )\right)\eta =j_\epsilon (\xi )\xi j_\epsilon (\eta )\eta j_\epsilon (\xi )(\xi \eta )$$
$$=_0^1𝑑\theta \left\{(\xi \eta )j_\epsilon (\xi _\theta )\xi _\theta +(j_\epsilon (\xi _\theta )j_\epsilon (\xi ))(\xi \eta )\right\}$$
$`(3.35)`$
with $`\xi _\theta =\theta \xi +(1\theta )\eta `$, so that
$$|j_\epsilon (\xi )j_\epsilon (\eta )||\eta |(||j_{\mathrm{}}+2)|\xi \eta |.$$
$`(3.36)`$
Comparing (3.33) (3.34) (3.36), we obtain
$$\left|j_\epsilon (\xi )|\xi |^kj_\epsilon (\eta )|\eta |^k\right||\eta |C|\xi \eta ||\eta |^kC|\xi |^{\nu /2}|\xi \eta |^{1\nu }|\eta |^{k+\nu /2}$$
$`(3.37)`$
from which we obtain as previously
$$|I_{1,2}|C|w|_{k+\nu /2}^2|\phi |_{\mathrm{}}.$$
$`(3.38)`$
Estimate of $`I_{1,3}`$. By reality and symmetry, which is the Fourier space version of integration by parts, $`I_{1,3}`$ can be rewritten as
$$I_{1,3}=(i/2)_{|\xi \eta ||\xi ||\eta |}𝑑\xi 𝑑\eta j_\epsilon (\xi )|\xi |^kf(\xi )\overline{\widehat{w}}(\xi )(\xi \eta )\widehat{s}(\xi \eta )j_\epsilon (\eta )|\eta |^kf(\eta )\widehat{w}(\eta ).$$
Using the second inequality in (3.37) we obtain as previously
$$|I_{1,3}|C|w|_{k+\nu /2}^2|s|_{\mathrm{}}.$$
$`(3.39)`$
We now turn to the contribution of the region $`|\xi ||\xi \eta ||\eta |`$.
Estimate of $`I_2`$. Using (3.3) and
$$|\xi |^{2k}|\eta |C|\xi |^{k+\nu /2}|\xi \eta |^{1\nu }|\eta |^{k+\nu /2}$$
$`(3.40)`$
we obtain as previously
$$|I_2|C|w|_{k+\nu /2}^2|\phi |_{\mathrm{}}.$$
$`(3.41)`$
We finally consider the region $`|\eta ||\xi ||\xi \eta |`$.
Estimate of $`|I_3|`$. Using (3.3) and
$$|\xi |^{2k}|\eta |C|\xi |^{k+\nu /2}|\xi \eta |^{k\nu /2+1\theta }|\eta |^\theta $$
with $`0\theta 1`$ and decomposing $`s=s_>+s_<`$, we estimate
$$|I_3|C|w|_{k+\nu /2}\{(|\xi |^{k\nu /2+1}f|\widehat{s}_<|)(f|\widehat{w}|)_2$$
$$+(|\xi |^{k\nu /2+1\theta }f|\widehat{s}_>|)(|\xi |^\theta f|\widehat{w}|)_2\}.$$
$`(3.42)`$
Using the Young inequality for the term in $`s_<`$ and Lemma 3.2 for the term in $`s_>`$, we obtain
$$|I_3|C|w|_{k+\nu /2}\left\{|\xi |^{k\nu /2+1}f\widehat{s}_<_1f\widehat{w}_2+|\xi |^{\mathrm{}+1}f\widehat{s}_>_2|\xi |^{k+\nu /2}f\widehat{w}_2\right\}$$
with $`\mathrm{}>n/2\nu `$ , $`k+\nu /2\theta `$ and $`\mathrm{}k\nu /2\theta `$. We choose $`\theta =1(k+\nu /2)`$ and the last two conditions reduce to $`\mathrm{}+1k\nu /2`$. Using (3.16), we obtain
$$|I_3|C|w|_{k+\nu /2}^2|\phi |_{\mathrm{}}.$$
$`(3.43)`$
Collecting (3.32) (3.38) (3.39) (3.41) (3.43) yields (3.22).
The constants $`C`$ appearing in the proof of (3.22) are independent of $`\rho `$. This can be checked explicitly on each case. More generally, it is a consequence of the following facts. We treat $`f`$ only through the inequalities (3.3) and (3.5) of Lemma 3.1, until we end up with integrals of the type (3.30) with however $`f`$ appearing only through the product $`f(\xi )f(\xi \eta )f(\eta )`$ (possibly after adding one missing $`f`$ and using the fact that $`f1`$). It then follows from (3.11) that an estimate of the type (3.22) of such a quantity for $`\rho =0`$ implies the same estimate for any $`\rho >0`$ with the same constant. The same argument applies to all subsequent estimates of Lemma 3.4 that contain $`f`$. Therefore from now on and in the same way as in the proof of Lemma 3.3, we omit $`f`$ in the proofs.
(3.23). We estimate
$$|\xi |^m(\widehat{sw})_2C\left\{|\widehat{s}|(|\xi |^{m+1}|\widehat{w}|)_2+(|\xi |^m|\widehat{s}|)(|\xi ||\widehat{w}|)_2\right\}.$$
We decompose $`s=s_<+s_>`$, we estimate the contribution of $`s_<`$ by the Young inequality and the contribution of $`s_>`$ by Lemma 3.2, thereby obtaining
$$\mathrm{}C\left\{\widehat{s}_<_1\left(|\xi |^{m+1}\widehat{w}_2+|\xi |\widehat{w}_2\right)+|\xi |^{\mathrm{}+1}\widehat{s}_>_2|\xi |^k\widehat{w}_2\right\}$$
under the condition stated on $`k`$, $`\mathrm{}`$, $`m`$. (3.23) then follows by (3.16).
(3.24). We estimate
$$|\xi |^m(\widehat{(s)w})_2C\left\{(|\xi |^{m+1}|\widehat{s}|)|\widehat{w}|_2+(|\xi ||\widehat{s}|)(|\xi |^m|\widehat{w}|)_2\right\}.$$
Proceeding as above, we obtain
$$\mathrm{}C\left\{\widehat{s}_<_1\left(|\xi |^m\widehat{w}_2+\widehat{w}_2\right)+|\xi |^{\mathrm{}+1}\widehat{s}_>_2|\xi |^k\widehat{w}_2\right\}$$
from which (3.24) follows by the same argument as above.
(3.25) and (3.26). We decompose $`s^{}=s_>^{}+s_<^{}`$. The contribution of $`s_>^{}`$ to (3.25) and (3.26) is estimated by (3.22) and (3.24) respectively, by replacing $`w`$ by $`s_>^{}`$ and $`k`$ by $`\mathrm{}^{}+1`$. In order to complete the proof it remains to estimate $`|\xi |_>^m(\widehat{ss_<^{}})_2`$ with $`m=\mathrm{}^{}+1\nu /2`$ for (3.25) and with general $`m`$ for (3.26), namely to estimate the $`L^2`$ norm in $`\xi `$ of the integral
$$J=|\xi |_>^mf(\xi )𝑑\eta \widehat{s}(\xi \eta )\eta \widehat{s}_<^{}(\eta ).$$
For that purpose, we decompose $`s=s_{>1/2}+s_{<1/2}`$ and correspondingly $`J=J_>+J_<`$. In $`J_>`$ we have $`|\eta |12|\xi \eta |`$ and therefore $`|\xi |3|\xi \eta |`$, so that
$$\begin{array}{cc}J_>_2\hfill & C(|\xi |^m|\widehat{s}_{>1/2}|)(|\eta ||\widehat{s}_<^{}|)_2\hfill \\ & \\ & C|s|_{\mathrm{}}\widehat{s}_<^{}_1\hfill \end{array}$$
$`(3.44)`$
by the Young inequality, provided $`m\mathrm{}+1`$, a condition which appears explicitly in (3.25) and which reduces to $`\mathrm{}\mathrm{}^{}\nu /2`$ in (3.26) for $`m=\mathrm{}^{}+1\nu /2`$.
In $`J_<`$, we have $`|\xi \eta |1/2`$, $`|\xi |1`$ and $`|\eta |1`$, and therefore $`1/2|\eta |1`$ and $`|\xi |3/2`$, so that
$$J_<_2C\widehat{s}_{<1/2}_1s_{>1/2}^{}_2.$$
$`(3.45)`$
(3.25) and (3.26) now follow from (3.22) (3.24), (3.44) (3.16) and (3.45).
(3.27) follows immediately from
$$|\xi |^\mathrm{}_<(\widehat{ss^{}})_2C\{(|\xi |^\mathrm{}<|\widehat{s}|)|\widehat{s}^{}|_2+(ss^{})\}$$
$$C\{|\xi |^\mathrm{}_<\widehat{s}_2\widehat{s}^{}_1+(ss^{})\}$$
and from (3.16) (3.17).
(3.28) follows from (3.14) either directly for $`\beta 0`$, or through the inequality
$$|\xi |^m\widehat{g_0(w_1}w_2)_>_2C\{(|\xi |^\beta |\widehat{w}_1|)(|\widehat{w}_2|)_2+(12)\}$$
for $`\beta 0`$.
(3.29) follows from (3.15) with $`m=\beta `$.
$``$$``$
We now explain the origin of the derivative loss in the system (2.11)-(2.12) for$`\lambda \mu n+2>0`$ and the mechanism by which that loss is overcome through the use of the spaces $`𝒳_{\rho ,loc}^{k,\mathrm{}}`$ defined by (3.13). If we try to solve the system (2.11)-(2.12) by the energy method in a space like $`𝒞(I,H^k\dot{H}^{\mathrm{}+2}`$), we have to estimate in particular
$$\{\begin{array}{c}_t^kw_2^2=2\mathrm{R}\mathrm{e}<^kw,^k_tw>\hfill \\ \\ _t^{\mathrm{}+2}\phi _2^2=2<^{\mathrm{}+2}\phi ,^{\mathrm{}+2}_t\phi >.\hfill \end{array}$$
$`(3.46)`$
The term with $`\mathrm{\Delta }\varphi `$ from $`_tw`$ forces us to apply $`k+2`$ derivatives to $`\phi `$ and requires therefore $`\mathrm{}k`$, while the term with $`g_0`$ from $`_t\phi `$ forces us to apply $`\mathrm{}+2`$ derivatives to $`g_0`$ or equivalently $`\mathrm{}+\lambda `$ derivatives to $`|w|^2`$ and requires therefore $`k\mathrm{}+\lambda `$. The terms with $`\phi w`$ from $`_tw`$ and with $`|\phi |^2`$ from $`_t\phi `$ can be handled essentially under the same assumptions, possibly after an integration by parts. The method therefore applies only if $`\lambda 0`$, namely $`\mu n2`$, which is the case treated in I and II.
If we try instead to solve the same problem in the space $`𝒳_{\rho ,loc}^{k,\mathrm{}}`$ with time dependent $`\rho `$, we have by (3.12)
$$\{\begin{array}{c}_t|w|_k^2=2\rho ^{}|w|_{k+\nu /2}^2+2\mathrm{R}\mathrm{e}<w,_tw>_k\hfill \\ \\ _t|\phi |_{\mathrm{}}^2=2\rho ^{}|\phi |_{\mathrm{}+\nu /2}^2+2<\phi ,_t\phi >_{\mathrm{}}\hfill \end{array}$$
$`(3.47)`$
where $`<,>_k`$ and $`<,>_{\mathrm{}}`$ denote the scalar products in $`K_\rho ^k`$ and $`Y_\rho ^{\mathrm{}}`$. If $`\rho ^{}`$ has a favourable sign, namely if $`\rho `$ decreases away from the initial time, the terms containing $`\rho ^{}`$ provide a control of the norm in $`L^2(K_\rho ^{k+\nu /2}Y_\rho ^{\mathrm{}+\nu /2})`$, and it suffices to control the scalar products at most quadratically in terms of the norms in $`K_\rho ^{k+\nu /2}`$ and $`Y_\rho ^{\mathrm{}+\nu /2}`$. In the term with $`\mathrm{\Delta }\varphi `$ from $`_tw`$, after distributing the function $`f`$ with the help of (3.3) and shifting $`\nu /2`$ derivatives on the first vector of the scalar product, it suffices to apply $`k\nu /2`$ derivatives on $`\mathrm{\Delta }\varphi `$ while one is allowed to use the norm $`|\phi |_{\mathrm{}+\nu /2}`$. This requires only $`\mathrm{}k\nu `$. Similarly in the term with $`g_0`$ from $`_t\phi `$, it suffices to apply $`\mathrm{}+2\nu /2`$ derivatives to $`g_0`$ or equivalently $`\mathrm{}+\lambda \nu /2`$ derivatives to $`|w|^2`$, while one is allowed to use $`|w|_{k+\nu /2}`$. This requires only $`k\mathrm{}+\lambda \nu `$. The two conditions on $`(k,\mathrm{})`$ are compatible provided $`\lambda 2\nu `$, which allows for $`\mu n`$ under that condition. It remains to estimate the terms $`\phi w`$ from $`_tw`$ and $`|\phi |^2`$ from $`_t\phi `$, which in the Sobolev case requires the integration by part of one derivative. Here however, by the same argument as above, it suffices to integrate by parts $`1\nu `$ derivative. Now it turns out that integration by parts of $`1\nu `$ derivative is exactly what is allowed by the inequality (3.5), which is exploited through (3.31) to derive the estimates (3.22) (3.25) where that integration by parts occurs. Actually the inequality (3.5) is optimal in the dangerous part $`|\xi \eta ||\xi ||\eta |`$ of the region where it is used, and more precisely when $`|\eta |\mathrm{}`$ for fixed $`|\xi \eta |`$. The conditions $`\mathrm{}k\nu `$ and $`k\mathrm{}+\lambda \nu `$ and therefore their consequence $`\lambda 2\nu `$ will appear from the next lemma onward as the most important part of the condition (3.48) and will propagate throughout this paper (except in Section 5) up to the main and final results of Propositions 6.5 and 7.5.
We now exploit Lemma 3.4 to derive energy like estimates for the solutions of the auxiliary system (2.11)-(2.12). In the following three lemmas, $`I`$ is an interval contained in $`[1,\mathrm{})`$, $`\rho `$ is a nonnegative continuous and piecewise $`𝒞^1`$ function defined in $`I`$. We shall be interested in solutions $`(w,\phi )`$ in spaces of the type $`𝒳_{\rho ,loc}^{k,\mathrm{}}(I)`$ for suitable values of $`k`$ and $`\mathrm{}`$. The estimates will hold in integrated form in any compact subinterval of $`I`$ under the available regularity, but will be stated in differential form for brevity.
Lemma 3.5. Let $`k`$, $`\mathrm{}`$ satisfy
$$\{\begin{array}{c}\mathrm{}>n/2\nu ,k\nu /2,\mathrm{}k\nu ,\hfill \\ \\ k\mathrm{}+\lambda \nu ,2k>\mathrm{}+\lambda \nu +n/2,\hfill \end{array}$$
$`(3.48)`$
where $`\lambda =\mu n+2`$.
Let $`(w,\phi )𝒳_{\rho ,loc}^{k,\mathrm{}}(I)`$ be a solution of the system (2.11)-(2.12). Then the following estimates hold :
$$\left|_t|w|_k^22\rho ^{}|w|_{k+\nu /2}^2\right|Ct^2\left\{|w|_{k+\nu /2}^2|\phi |_{\mathrm{}}+|w|_{k+\nu /2}|\phi |_{\mathrm{}+\nu /2}|w|_k\right\},$$
$`(3.49)`$
$$\left|_t|\phi |_{\mathrm{}}^22\rho ^{}|\phi |_{\mathrm{}+\nu /2}^2\right|Ct^2|\phi |_{\mathrm{}+\nu /2}^2|\phi |_{\mathrm{}}+Ct^\gamma |\phi |_{\mathrm{}+\nu /2}|w|_{k+\nu /2}|w|_k.$$
$`(3.50)`$
Proof. We use the same regularization as in the proof of Lemma 3.4. From (3.12) we obtain
$$_t|J_\epsilon w|_k^22\rho ^{}|J_\epsilon w|_{k+\nu /2}^2=2\mathrm{R}\mathrm{e}(<j_\epsilon |\xi |^kf\widehat{w}_>,j_\epsilon |\xi |^kf_t\widehat{w}_>>+<j_\epsilon f\widehat{w}_<,j_\epsilon f_t\widehat{w}_<>).$$
$`(3.51)`$
We substitute $`_tw`$ from (2.11) and we estimate the various terms successively. The term with $`\mathrm{\Delta }w`$ does not contribute. The term $`sw`$ is estimated by (3.22) under the conditions
$$\mathrm{}>n/2\nu ,k\nu /2,\mathrm{}+1k\nu /2.$$
The term $`(s)w`$ is estimated by (3.24) with $`m=k\nu /2`$ or $`m=0`$, and $`\mathrm{}`$ replaced by $`\mathrm{}+\nu /2`$, under the conditions
$$\mathrm{}>n/2\nu ,\mathrm{}k\nu .$$
Substituting those two estimates into (3.51), integrating over time and taking the limit $`\epsilon 0`$ yields the integrated form of (3.49). The required conditions on $`k`$, $`\mathrm{}`$ are implied by (3.48).
Similarly, we obtain from (3.12)
$$_t|J_\epsilon \phi |_{\mathrm{}}^22\rho ^{}|J_\epsilon \phi |_{\mathrm{}+\nu /2}^2=2<j_\epsilon |\xi |^{\mathrm{}+1}f\widehat{s}_>,j_\epsilon |\xi |^{\mathrm{}+1}f_t\widehat{s}_>>+2<j_\epsilon |\xi |^\mathrm{}_<f\widehat{\phi }_<,j_\epsilon |\xi |^\mathrm{}_<f_t\widehat{\phi }_<>.$$
$`(3.52)`$
We substitute $`_ts`$ and $`_t\phi `$ from (2.27) (2.12) into (3.52) and we estimate the various terms successively. The term $`ss`$ from $`_ts`$ is estimated by (3.25) with $`s^{}=s`$, $`\mathrm{}^{}=\mathrm{}`$ under the conditions
$$\mathrm{}>n/2\nu ,\mathrm{}+1\nu /2.$$
The term $`g_0`$ from $`_ts`$ is estimated by (3.28) with $`m=\mathrm{}+2\nu /2`$, $`w_1=w_2=w`$, $`k_1(=k_2^{})=k+\nu /2`$ , $`k_2(=k_1^{})=k`$ , under the conditions
$$2k>\mathrm{}+\lambda \nu +n/2,k\mathrm{}+\lambda \nu .$$
The terms $`|s|^2`$ and $`g_0`$ from $`_t\phi _<`$ are estimated by (3.27) and (3.29) respectively. Using those estimates yields (3.50) in the same way as above. The required conditions on $`k`$, $`\mathrm{}`$ are implied by (3.48). In particular the condition $`\mathrm{}+1\nu /2`$ follows from $`\mathrm{}>n/2\nu `$.
$``$$``$
The next lemma is a regularity result.
Lemma 3.6. Let $`k`$, $`\mathrm{}`$ satisfy (3.48) and let $`\overline{k}`$, $`\overline{\mathrm{}}`$ satisfy
$$\overline{k}k=\overline{\mathrm{}}\mathrm{}0.$$
$`(3.53)`$
Let $`(w,\phi )𝒳_{\rho ,loc}^{\overline{k},\overline{\mathrm{}}}(I)`$ be a solution of the system (2.11)-(2.12). Then the following estimates hold :
$$\left|_t|w|_{\overline{k}}^22\rho ^{}|w|_{\overline{k}+\nu /2}^2\right|Ct^2\left\{|w|_{\overline{k}+\nu /2}^2|\phi |_{\mathrm{}}+|w|_{\overline{k}+\nu /2}|\phi |_{\overline{\mathrm{}}+\nu /2}|w|_k\right\},$$
$`(3.54)`$
$$\left|_t|\phi |_\overline{\mathrm{}}^22\rho ^{}|\phi |_{\overline{\mathrm{}}+\nu /2}^2\right|Ct^2|\phi |_{\overline{\mathrm{}}+\nu /2}^2|\phi |_{\mathrm{}}+Ct^\gamma |\phi |_{\overline{\mathrm{}}+\nu /2}|w|_{\overline{k}+\nu /2}|w|_k.$$
$`(3.55)`$
Proof. The proof follows the same pattern as that of Lemma 3.5 and we concentrate on the differences, which bear on the estimates connected with Lemma 3.4, omitting the regularization for brevity. We estimate only the contribution of the high $`|\xi |`$ region, since that of the low $`|\xi |`$ region is already estimated by Lemma 3.5.
We first estimate $`_t|w|\frac{2}{k}`$, starting from (3.51) with $`k`$ replaced by $`\overline{k}`$ and with $`J_\epsilon `$ omitted and we estimate successively the contribution of the various terms form (2.11).
The contribution of $`sw`$, namely
$$2\mathrm{R}\mathrm{e}<|\xi |^{\overline{k}}f\widehat{w},|\xi |^{\overline{k}}f(\widehat{sw})>$$
is estimated in the same way as in the proof of (3.22). The contributions $`I_1`$ and $`I_2`$ of the regions $`|\xi \eta ||\xi ||\eta |`$ and $`|\xi ||\xi \eta ||\eta |`$ are estimated as in the latter with $`k`$ replaced by $`\overline{k}`$ under the conditions $`\overline{k}\nu /2`$ , $`\mathrm{}>n/2\nu `$.
The contribution of the region $`|\eta ||\xi ||\xi \eta |`$ is estimated by (3.42) with $`k`$ replaced by $`\overline{k}`$ and with $`\theta =0`$, or equivalently
$$|I_3|C|w|_{\overline{k}+\nu /2}\left\{f\widehat{s}_<_1f\widehat{w}_2+(|\xi |^{\overline{k}\nu /2+1}f|\widehat{s}_>|)(f|\widehat{w}|)_2\right\}.$$
The last norm is then estimated by Lemma 3.2, thereby yielding
$$|I_3|C|w|_{\overline{k}+\nu /2}|\phi |_{\overline{\mathrm{}}+\nu /2}|w|_k$$
under the conditions
$$k+\overline{\mathrm{}}>\overline{k}+n/2\nu ,\overline{\mathrm{}}\overline{k}\nu $$
which for $`\overline{k}k=\overline{\mathrm{}}\mathrm{}`$ reduce to
$$\mathrm{}>n/2\nu ,\mathrm{}k\nu .$$
The contribution of $`(s)w`$ is estimated by
$$|w|_{\overline{k}+\nu /2}|\xi |^{\overline{k}\nu /2}f(\widehat{(s)w})_2$$
and subsequently in a way similar to the proof of (3.24) with $`m=\overline{k}\nu /2`$. Using the elementary inequality
$$|\xi |^m|\xi \eta |2^m\left(|\xi \eta |^{m+1}+|\eta |^{m+\theta }|\xi \eta |^{1\theta }\right)$$
$`(3.56)`$
valid for $`m0`$ and $`0\theta 1`$, we estimate the last norm by
$$\mathrm{}C\left\{(|\xi |^{m+1}f|\widehat{s}|)(f|\widehat{w}|)_2+(|\xi |^{1\theta }f|\widehat{s}|)(|\xi |^{m+\theta }f|\widehat{w}|)_2\right\}.$$
We then estimate the first term by Lemma 3.2 and the second term with $`\theta =\nu `$ by the Young inequality and (3.16) (3.17), thereby obtaining
$$\mathrm{}C\left(|\phi |_{\overline{\mathrm{}}+\nu /2}|w|_k+|\phi |_{\mathrm{}}|w|_{\overline{k}+\nu /2}\right)$$
under the same conditions as before, namely
$$k+\overline{\mathrm{}}>\overline{k}+n/2\nu ,\overline{\mathrm{}}\overline{k}\nu .$$
This completes the proof of (3.54).
We next estimate $`_t|\phi |_\overline{\mathrm{}}^2`$ by substituting similarly the various terms from (2.12) (2.27) into the right-hand side of (3.52) with $`\mathrm{}`$ replaced by $`\overline{\mathrm{}}`$ and $`J_\epsilon `$ omitted. The contribution of $`ss`$ to the high $`|\xi |`$ part of the norm, namely
$$<|\xi |^{\overline{\mathrm{}}+1}f\widehat{s}_>,|\xi |^{\overline{\mathrm{}}+1}f(\widehat{ss})_>>$$
is estimated by modifying the proof of (3.25) along the same lines as that of (3.22) above.
The contribution of $`g_0(w,w)_>`$ is estimated by
$$|\phi |_{\overline{\mathrm{}}+\nu /2}|\xi |^{\overline{\mathrm{}}+1\nu /2}f\widehat{g_0(w,}w)_>_2$$
followed by (3.28) with $`m=\overline{\mathrm{}}+2\nu /2`$ , $`k_1(=k_2^{})=\overline{k}+\nu /2`$ , $`k_2(=k_1^{})=k`$ under the conditions
$$k+\overline{k}>\overline{\mathrm{}}+\lambda \nu +n/2,\overline{k}\overline{\mathrm{}}+\lambda \nu $$
which reduce to the same conditions with $`(\overline{k},\overline{\mathrm{}})`$ replaced by $`(k,\mathrm{})`$ under the condition $`\overline{k}k=\overline{\mathrm{}}\mathrm{}`$.
This completes the proof of (3.55).
$``$$``$
We next estimate the difference between two solutions of the system (2.11)-(2.12).
Lemma 3.7. Let $`k`$, $`\mathrm{}`$ satisfy (3.48) and let $`k^{}`$, $`\mathrm{}^{}`$ satisfy
$$k^{}\nu /2,kk^{}=\mathrm{}\mathrm{}^{}1\nu .$$
$`(3.57)`$
Let $`(w_1,\phi _1)`$ and $`(w_2,\phi _2)𝒳_{\rho ,loc}^{k,\mathrm{}}(I)`$ be two solutions of the system (2.11)-(2.12) and let $`w_\pm =w_1\pm w_2`$, $`\phi _\pm =\phi _1\pm \phi _2`$. Then the following estimates hold :
$$|_t|w_{}|_k^{}^22\rho ^{}|w_{}|_{k^{}+\nu /2}^2|Ct^2\{|w_{}|_{k^{}+\nu /2}^2|\phi _+|_{\mathrm{}}$$
$$+|w_{}|_{k^{}+\nu /2}(|w_+|_{k+\nu /2}|\phi _{}|_{\mathrm{}^{}}+|\phi _+|_{\mathrm{}+\nu /2}|w_{}|_k^{}+|\phi _{}|_{\mathrm{}^{}+\nu /2}|w_+|_k)\},$$
$`(3.58)`$
$$\left|_t|\phi _{}|_{\mathrm{}^{}}^22\rho ^{}|\phi _{}|_{\mathrm{}^{}+\nu /2}^2\right|Ct^2\left\{|\phi _{}|_{\mathrm{}^{}+\nu /2}^2|\phi _+|_{\mathrm{}}+|\phi _{}|_{\mathrm{}^{}+\nu /2}|\phi _{}|_{\mathrm{}^{}}|\phi _+|_{\mathrm{}+\nu /2}\right\}$$
$$+Ct^\gamma |\phi _{}|_{\mathrm{}^{}+\nu /2}\left\{|w_+|_{k+\nu /2}|w_{}|_k^{}+|w_{}|_{k^{}+\nu /2}|w_+|_k\right\}.$$
$`(3.59)`$
Proof. The proof follows the same pattern as that of Lemma 3.5, using the estimates of Lemma 3.4. We omit again the regularization for brevity. The equations satisfied by $`(w_{},\phi _{})`$ are
$$_tw_{}=i(2t^2)^1\mathrm{\Delta }w_{}+(2t)^2\left\{2s_+w_{}+2s_{}w_++(s_+)w_{}+(s_{})w_+\right\},$$
$`(3.60)`$
$$_t\phi _{}=(2t^2)^1(s_+s_{})+t^\gamma g_0(w_+,w_{}),$$
$`(3.61)`$
and in the same way as in the proof of Lemma 3.5, we shall also use the equation for $`s_{}`$ obtained by taking the gradient of (3.61), namely
$$_ts_{}=(2t^2)^1\left(s_+s_{}+s_{}s_+\right)+t^\gamma g_0(w_+,w_{}).$$
$`(3.62)`$
From (3.12), we obtain
$$|_t|w_{}]_k^{}^22\rho ^{}|w_{}|_{k^{}+\nu /2}^2|=2\mathrm{R}\mathrm{e}(<|\xi |^k^{}f\widehat{w}_>,|\xi |^k^{}f_t\widehat{w}_>>+<f\widehat{w}_<,f_t\widehat{w}_<>).$$
$`(3.63)`$
We substitute $`_tw_{}`$ from (3.60) into (3.63) and we estimate the various terms successively. The term $`\mathrm{\Delta }w_{}`$ does not contribute. We consider only the contribution of the high $`|\xi |`$ region. The low $`|\xi |`$ region is treated in a similar and simpler way.
Term with $`s_+w_{}`$. We apply (3.22) with $`k`$ replaced by $`k^{}`$, thereby obtaining a contribution
$$|w_{}|_{k^{}+\nu /2}^2|\phi _+|_{\mathrm{}}$$
under conditions which follows from (3.48) and from $`\nu /2k^{}k`$.
Term with $`s_{}w_+`$. We apply (3.23) with $`m=k^{}\nu /2`$ , $`\mathrm{}`$ replaced by $`\mathrm{}^{}`$ and $`k`$ replaced by $`k+\nu /2`$, thereby obtaining a contribution
$$|w_{}|_{k^{}+\nu /2}|w_+|_{k+\nu /2}|\phi _{}|_{\mathrm{}^{}}$$
under the conditions
$$k+\mathrm{}^{}>k^{}+n/2\nu ,k+\nu k^{}+1,\mathrm{}+1k^{}\nu /2,$$
which follow from (3.48) (3.57).
Term with $`(s_+)w_{}`$. We apply (3.24) with $`m=k^{}\nu /2`$, $`k`$ replaced by $`k^{}`$ and $`\mathrm{}`$ replaced by $`\mathrm{}+\nu /2`$, thereby obtaining a contribution
$$|w_{}|_{k^{}+\nu /2}|\phi _+|_{\mathrm{}+\nu /2}|w_{}|_k^{}$$
under the conditions
$$\mathrm{}>n/2\nu ,\mathrm{}k^{}\nu ,$$
which follow from (3.48) and from $`k^{}k`$.
Term with $`(s_{})w_+`$. We apply (3.24) with $`m=k^{}\nu /2`$, and $`\mathrm{}`$ replaced by $`\mathrm{}^{}+\nu /2`$, thereby obtaining a contribution
$$|w_{}|_{k^{}+\nu /2}|\phi _{}|_{\mathrm{}^{}+\nu /2}|w_+|_k$$
under the conditions
$$k+\mathrm{}^{}>k^{}+n/2\nu ,\mathrm{}^{}k^{}\nu ,kk^{}\nu /2,$$
which follow from (3.48) (3.57).
Collecting the previous four estimates together with the contribution of the low $`|\xi |`$ region, yields (3.58).
We now turn to the estimate of $`\phi _{}`$. From (3.12) we obtain
$$_t|\phi _{}|_{\mathrm{}^{}}^22\rho ^{}|\phi _{}|_{\mathrm{}^{}+\nu /2}^2=2<|\xi |^{\mathrm{}^{}+1}f\widehat{s}_>,|\xi |^{\mathrm{}^{}+1}f_t\widehat{s}_>>+2<|\xi |^\mathrm{}_<f\widehat{\phi }_<,|\xi |^\mathrm{}_<f\widehat{\phi }_<>.$$
$`(3.64)`$
We substitute $`_ts_{}`$ and $`_t\phi _{}`$ from (3.62) and (3.61) into (3.64) and we estimate the various terms successively.
Term with $`s_+s_{}`$. We apply (3.25) with $`s^{}=s_{}`$ and obtain a contribution
$$|\phi _{}|_{\mathrm{}^{}+\nu /2}^2|\phi _+|_{\mathrm{}}$$
under the conditions
$$\mathrm{}>n/2\nu ,\mathrm{}^{}+1\nu /2,\mathrm{}\mathrm{}^{}\nu /2,$$
which follow from (3.48), from $`\mathrm{}^{}+1k^{}+1\nu 1\nu /2\nu /2`$ and from $`\mathrm{}\mathrm{}^{}`$.
Term with $`s_{}s_+`$. We apply (3.26) with $`m=\mathrm{}^{}+1\nu /2`$ , $`\mathrm{}`$ replaced by $`\mathrm{}^{}`$ and $`\mathrm{}^{}`$ replaced by $`\mathrm{}+\nu /2`$, thereby obtaining a contribution
$$|\phi _{}|_{\mathrm{}^{}+\nu /2}|\phi _{}|_{\mathrm{}^{}}|\phi _+|_{\mathrm{}+\nu /2}$$
under the conditions
$$\mathrm{}>n/2\nu ,\mathrm{}\mathrm{}^{}+1\nu ,$$
which follow from (3.48) (3.57).
Term with $`g_0(w_+,w_{})_>`$. We apply (3.28) with $`m=\mathrm{}^{}+1\nu /2`$, $`w_1=w_+`$, $`w_2=w_{}`$, $`k_1=k+\nu /2`$ , $`k_2=k^{}`$ , $`k_1^{}=k`$ , $`k_2^{}=k^{}+\nu /2`$, thereby obtaining a contribution
$$|\phi _{}|_{\mathrm{}^{}+\nu /2}\left\{|w_+|_{k+\nu /2}|w_{}|_k^{}+|w_{}|_{k^{}+\nu /2}|w_+|_k\right\}$$
under the conditions
$$k+k^{}+\nu >\mathrm{}^{}+\mu n+2+n/2,k^{}+\nu \mathrm{}^{}+\mu n+2$$
which follow from (3.48) (3.57).
The terms with $`(s_+,s_{})_<`$ and $`g_0(w_+w_{})_<`$ are treated by the use of (3.27) and (3.29) as in the proof of Lemma 3.5.
Collecting the previous estimates yields (3.59).
$``$$``$
We conclude this section by introducing a number of estimating functions of time generalizing those introduced in Section II.3 and by deriving a number of estimates for them. Those functions will be defined in terms of the derivative $`h_0^{}`$ of a given function $`h_0`$ on which we make the following assumptions
$$h_0𝒞^1([1,\mathrm{}),IR^+),h_0^{}0,t^2h_0(t)L^1([1,\mathrm{})),t^1h_0^{}(t)L^1([1,\mathrm{})).$$
$`(3.65)`$
From the relation
$$t_2^1h_0(t_2)t_1^1h_0(t_1)=_{t_1}^{t_2}𝑑tt^1h_0^{}(t)_{t_1}^{t_2}𝑑tt^2h_0(t)$$
$`(3.66)`$
it follows that the last condition on $`h_0`$ in (3.65) can be replaced by the condition that $`t^1h_0(t)`$ tends to zero when $`t\mathrm{}`$. A typical example for $`h_0^{}`$ is that considered in Section II.3, namely $`h_0^{}(t)=t^\gamma `$.
The first and basic estimating function $`h`$ is defined by
$$h(t)=_1^{\mathrm{}}𝑑t_1(tt_1)^1h_0^{}(t_1)$$
$`(3.67)`$
from which it follows that $`h(t)`$ is decreasing in $`t`$ and tends to zero when $`t\mathrm{}`$, while $`th(t)`$ is increasing in $`t`$. We next define for any $`m0`$
$$N_m(t)=_1^t𝑑t_1h_0^{}(t_1)h^m(t_1),$$
$`(3.68)`$
$$Q_m(t)=_1^{\mathrm{}}𝑑t_1(tt_1)^1h_0^{}(t_1)h^m(t_1),$$
$`(3.69)`$
where the integral in (3.69) is convergent since $`t^1h_0^{}(t)L^1([1,\mathrm{}))`$ and since $`h(t)`$ is decreasing in $`t`$. It follows from (3.68) that $`N_m(t)`$ is increasing in $`t`$ and from (3.69) that $`Q_m(t)`$ is decreasing in $`t`$ and tends to zero when $`t\mathrm{}`$, while $`tQ_m(t)`$ is increasing in $`t`$, so that $`Q_m(t)t^1Q_m(1)`$. Moreover, for any nonnegative integers $`i`$ and $`j`$
$$N_{i+j}(t)h(1)^iN_j(t)h(1)^{i+j}N_0(t),$$
$`(3.70)`$
$$Q_{i+j}(t)h(1)^iQ_j(t)h(1)^{i+j}h(t)h(1)^{i+j+1}.$$
$`(3.71)`$
Clearly $`N_0(t)=h_0(t)h_0(1)`$ and $`Q_0=h`$. It will be convenient to introduce the notation $`Q_1=1`$.
Finally we set
$$P_m(t)=_1^{\mathrm{}}𝑑t_1h(tt_1)h_0^{}(t_1)h^m(t_1)$$
$`(3.72)`$
which is well defined provided
$$P_m(1)=_1^{\mathrm{}}𝑑t_1h_0^{}(t_1)h^{m+1}(t_1)<\mathrm{}.$$
$`(3.73)`$
It follows from (3.72) that $`P_m(t)`$ is decreasing in $`t`$ and tends to zero when $`t\mathrm{}`$ while $`h(t)^1P_m(t)`$ is increasing in $`t`$, so that
$$P_m(1)h(t)h(1)P_m(t)h(1)P_m(1).$$
$`(3.74)`$
We now collect a number of identities and inequalities satisfied by the previous estimating functions.
Lemma 3.8 Let $`i`$, $`j`$ and $`m`$ be nonnegative integers, let $`1ab`$ and $`t1`$. Then the following identities and inequalities hold :
$$_t^{\mathrm{}}𝑑t_1t_1^2N_m(t_1)=Q_m(t_1)$$
$`(3.75)`$
$$_1^t𝑑t_1t_1^2N_0(t_1)N_m(t_1)=N_{m+1}(t)h(t)N_m(t)N_{m+1}(t)$$
$`(3.76)`$
$$_t^{\mathrm{}}𝑑t_1t_1^2N_0(t_1)N_m(t_1)=P_m(t)\text{if (3.73) holds}$$
$`(3.77)`$
$$N_i(t)N_j(t)N_0(t)N_{i+j}(t)$$
$`(3.78)`$
$$N_i(t)Q_j(t)h(t)N_{i+j}(t)N_{i+j+1}(t)$$
$`(3.79)`$
$$Q_i(t)Q_j(t)h(t)Q_{i+j}(t)2Q_{i+j+1}(t)$$
$`(3.80)`$
$$_t^{\mathrm{}}𝑑t_1h_0^{}(t_1)h(t_1)Q_{m1}(t_1)_t^{\mathrm{}}𝑑t_1h_0^{}(t_1)Q_m(t_1)$$
$`(3.81)`$
$$_t^{\mathrm{}}𝑑t_1h_0^{}(t_1)Q_m(t_1)P_m(t_1)\text{if (3.73) holds}$$
$`(3.82)`$
$$_1^t𝑑t_1h_0^{}(t_1)h(t_1)Q_{m1}(t_1)N_{m+1}(t_1)$$
$`(3.83)`$
$$_1^t𝑑t_1h_0^{}(t_1)Q_m(t_1)N_{m+1}(t)$$
$`(3.84)`$
$$_a^b𝑑th_0^{}(t)Q_m(t)Q_m(a)\left(h_0(b)h_0(a)\right)$$
$`(3.85)`$
$$_a^b𝑑th_0^{}(t)h(t)Q_{m1}(t)2Q_m(a)\left(h_0(b)h_0(a)\right).$$
$`(3.86)`$
Proof. This lemma is a generalization of Lemma II.3.6 and most of the proofs are obtained by manipulations of integrals similar to those of the corresponding integrals in Lemma II.3.6, after the replacement of $`t^\gamma `$ by $`h_0^{}(t)`$. This applies to (3.75) (3.76) (3.77) (3.81) (3.82) (3.83) (3.84) and (3.80b). (The latter is the generalization of (II.3.49)). The estimates (3.78) and (3.80a) follow from the Hölder inequality. The estimate (3.79a) is the pointwise counterpart of (II.3.39) and is proved in the same way. The estimate (3.79b) follows from the decrease of $`h`$, (3.85) follows from the decrease of $`Q_m`$, and (3.86) follows from (3.80b) and from (3.85).
$``$$``$
## 4 Cauchy problem and preliminary asymptotics for the auxiliary system
In this section, we study the existence of solutions in a neighborhood of infinity in time for the auxiliary system (2.11)-(2.12) and we derive some preliminary results on the asymptotic behaviour in time of the solutions thereby obtained. We first introduce some notation.
We choose once for all a strictly positive $`𝒞^1`$ function defined in $`[1,\mathrm{})`$, which we call $`|\rho ^{}|`$ for reasons that will soon become obvious. We assume that $`|\rho ^{}|`$ satisfies the following properties
(i) $`|\rho ^{}|L^1([1,\mathrm{}))`$,
(ii) the function $`t^\gamma |\rho ^{}|^1`$ is nondecreasing (and therefore $`|\rho ^{}|`$ is nonincreasing), and the function $`t^2|\rho ^{}|^1`$ is nonincreasing and tends to zero at infinity.
Typical examples of suitable functions $`|\rho ^{}|`$ are $`t^{1\epsilon }`$ for $`\epsilon `$ sufficiently small, depending on $`\gamma `$, or $`t^1(\alpha +\mathrm{}nt)^\alpha `$ for $`\alpha >1`$. It will be useful to keep those examples in mind in order to understand the time decay implied by the subsequent estimates.
Let now $`I[1,\mathrm{})`$ be an interval, possibly unbounded, and let $`t_0I`$ (or $`t_0=\mathrm{}`$ if $`I`$ is unbounded). We define a function $`\rho `$ in $`I`$ by
$$\rho (t)=\rho (t_0)\left|_{t_0}^t𝑑t_1|\rho ^{}(t_1)|\right|$$
$`(4.1)`$
so that $`\rho `$ is increasing (resp. decreasing) for $`tt_0`$ (resp. $`tt_0`$) and has $`|\rho ^{}|`$ as the absolute value of its derivative, which justifies the notation. We take $`\rho (t_0)`$ sufficiently large so that $`\rho `$ is nonnegative in $`I`$. All subsequent estimates will be independent of $`\rho (t_0)`$ (they will however depend on $`|\rho ^{}|`$). The previous choice of $`\rho `$ will be used in this section without further comment unless otherwise stated. We now define the global version of the fundamental spaces, corresponding to the local version (3.13). We define
$$𝒳_\rho ^{k,\mathrm{}}(I)\left(𝒞L^{\mathrm{}}\right)(I,K_\rho ^kY_\rho ^{\mathrm{}})L_{|\rho ^{}|}^2(I,K_\rho ^{k+\nu /2}Y_\rho ^{\mathrm{}+\nu /2})$$
$`(4.2)`$
where $`L_{|\rho ^{}|}^2`$ denotes weighted $`L^2`$ in time with weight $`|\rho ^{}|`$. More precisely, especially as regards continuity, $`(w,\phi )𝒳_\rho ^{k,\mathrm{}}(I)`$ is understood to mean that
$$(F^1f\widehat{w},F^1f\widehat{\phi })𝒳_0^{k,\mathrm{}}(I)\left(𝒞L^{\mathrm{}}\right)(I,K_0^kY_0^{\mathrm{}})L_{|\rho ^{}|}^2(I,K_0^{k+\nu /2}Y_0^{\mathrm{}+\nu /2}).$$
$`(4.3)`$
Note that in (4.3) we keep the weight $`|\rho ^{}|`$ in the $`L^2`$ part.
The norm of $`(w,\phi )`$ in $`𝒳_\rho ^{k,\mathrm{}}`$ is made up of several pieces for which we introduce additional notation which will be most helpful in the derivation of the estimates. Let $`H`$ be a continuous strictly positive function defined in $`I`$. We define
$$y(w;I,H,k)=\underset{tI}{Sup}H^1(t)|w(t)|_k,$$
$`(4.4)`$
$$y_1(w;I,H,k)=\left\{_I𝑑t\right|\rho ^{}(t)|H^2(t)|w(t)|_{k+\nu /2}^2\}^{1/2},$$
$`(4.5)`$
$$Y(w;I,H,k)=(yy_1)(w;I,H,k),$$
$`(4.6)`$
$$z(\phi ;I,H,\mathrm{})=\underset{tI}{Sup}H^1(t)|\phi (t)|_{\mathrm{}},$$
$`(4.7)`$
$$z_1(w;I,H,\mathrm{})=\left\{_I𝑑t\right|\rho ^{}(t)|H^2(t)|\phi (t)|_{\mathrm{}+\nu /2}^2\}^{1/2},$$
$`(4.8)`$
$$Z(\phi ;I,H,\mathrm{})=(zz_1)(\phi ,I,H,\mathrm{}).$$
$`(4.9)`$
We then take
$$(w,\phi );𝒳_\rho ^{k,\mathrm{}}(I)=Y(w;I,1,k)+Z(\phi ;I,1,\mathrm{}).$$
$`(4.10)`$
In order to study the Cauchy problem for the auxiliary system (2.11)-(2.12) we need additional a priori estimates of solutions of that system, which are a continuation of those of Lemmas 3.5, 3.6 and 3.7 in which we now take into account the dependence on time in the framework of the spaces $`𝒳_\rho ^{k,\mathrm{}}`$ just defined.
Lemma 4.1. Let $`k`$, $`\mathrm{}`$ satisfy (3.48). Let $`1Tt_0<\mathrm{}`$ and let $`(w,\phi )𝒳_{\rho ,loc}^{k,\mathrm{}}([T,\mathrm{}))`$ be a solution of the system (2.11)-(2.12). Let $`h_0`$ and $`h_1`$ be $`𝒞^1`$ positive functions defined in $`[T,\mathrm{})`$, with $`h_0`$ nondecreasing, $`h_1`$ nonincreasing, $`h_0t^\gamma |\rho ^{}|^1`$ and $`h_1t^2|\rho ^{}|^1h_0`$. Let $`a_0=|w(t_0)|_k`$ and $`b_0=h_0(t_0)^1|\phi (t_0)|_{\mathrm{}}`$.
(1) There exist constants $`c`$ and $`C`$ such that if
$$\left(b_0+a_0^2\right)h_1(t_0)c,$$
$`(4.11)`$
then $`(w,h_0^1\phi )𝒳_\rho ^{k,\mathrm{}}([t_0,\mathrm{}))`$ and $`(w,\phi )`$ satisfies the estimates
$$\{\begin{array}{c}Y(w;[t_0,\mathrm{}),1,k)Ca_0,\hfill \\ \\ Z(\phi ;[t_0,\mathrm{}],h_0,\mathrm{})C\left(b_0+a_0^2\right).\hfill \end{array}$$
$`(4.12)`$
(2) There exist constants $`c`$ and $`C`$ such that if
$$\left(b_0+a_0^2\right)T^2|\rho ^{}(T)|^1h_0(t_0)c,$$
$`(4.13)`$
then $`(w,\phi )`$ satisfies the estimates
$$\{\begin{array}{c}Y(w;[T,t_0],1,k)Ca_0,\hfill \\ \\ Z(\phi ;[T,t_0],1,\mathrm{})C\left(b_0+a_0^2\right)h_0(t_0).\hfill \end{array}$$
$`(4.14)`$
Proof. The proof requires the same regularization procedure as that of (the integrated form of) the estimates of Lemma 3.5, but we omit it for brevity. We begin the proof by treating simultaneously the cases $`tt_0`$ and $`tt_0`$. Let $`H`$ be a $`𝒞^1`$ positive function of time, increasing for $`tt_0`$ and decreasing for $`tt_0`$ and let $`\stackrel{~}{\phi }=H^1\phi `$. From Lemma 3.5 we obtain
$$_t|w|_k^2\underset{>}{<}2\rho ^{}|w|_{k+\nu /2}^2\pm Ct^2\left\{|w|_{k+\nu /2}^2|\phi |_{\mathrm{}}+|w|_{k+\nu /2}|\phi |_{\mathrm{}+\nu /2}|w|_k\right\}$$
$`(4.15)`$
$$_t|\stackrel{~}{\phi }|_{\mathrm{}}^2\underset{>}{<}2\rho ^{}|\stackrel{~}{\phi }|_{\mathrm{}+\nu /2}^2\pm Ct^2|\stackrel{~}{\phi }|_{\mathrm{}+\nu /2}^2|\phi |_{\mathrm{}}\pm Ct^\gamma H^1|\stackrel{~}{\phi }|_{\mathrm{}+\nu /2}|w|_{k+\nu /2}|w|_k$$
$`(4.16)`$
for $`t\underset{<}{>}t_0`$. In (4.16), we have dropped the term $`2H^{}H^1|\stackrel{~}{\phi }|_{\mathrm{}}^2`$ coming from the derivative of $`H`$. Let now $`y_{(1)}=y_{(1)}(w;I,1,k)`$ and $`z_{(1)}=z_{(1)}(\phi ;I,H,\mathrm{})`$ where $`I=[t_0,t]`$ for $`tt_0`$ and $`I=[t,t_0]`$ for $`tt_0`$. Integrating (4.15) (4.16) over time and using the Schwarz inequality, we obtain
$$|w(t)|_k^2+y_1^2y_0^2+C\{\underset{tI}{Sup}t^2|\rho ^{}|^1H\}(y_1^2z+y_1z_1y)$$
$`(4.17)`$
$$|\stackrel{~}{\phi }(t)|_{\mathrm{}}^2+z_1^2z_0^2+C\left\{\underset{tI}{Sup}(t^2|\rho ^{}|^1H)\right\}z_1^2z+C\{\underset{tI}{Sup}t^\gamma |\rho ^{}|^1H^1\}z_1y_1y$$
$`(4.18)`$
where $`y_0=|w(t_0)|_k=a_0`$ and $`z_0=|\stackrel{~}{\phi }(t_0)|_{\mathrm{}}`$, and since the RHS of (4.17) (4.18) is increasing in $`|tt_0|`$,
$$y^2y_1^2\text{RHS of (4.17)}$$
$$z^2z_1^2\text{RHS of (4.18)}$$
so that $`Y=yy_1`$ and $`Z=zz_1`$ satisfy
$$\{\begin{array}{c}Y^2y_0^2+m_1Y^2Z\hfill \\ \\ Z^2z_0^2+m_1Z^3+m_2ZY^2\hfill \end{array}$$
with
$$m_1=C\underset{tI}{Sup}t^2|\rho ^{}|^1H,m_2=C\underset{tI}{Sup}(t^\gamma |\rho ^{}|^1H^1).$$
If we can arrange that $`m_1Z1/2`$, then we obtain
$$\{\begin{array}{c}Y^22y_0^2\hfill \\ \\ Z^22z_0^2+2m_2ZY^22z_0^2+4m_2y_0^2Z\hfill \end{array}$$
so that
$$\{\begin{array}{c}Y2a_0\hfill \\ \\ Z2z_0+4m_2a_0^2\hfill \end{array}$$
$`(4.19)`$
which will yield the final estimates, while the condition $`m_1Z1/2`$ is implied by
$$4m_1\left(z_0+2m_2a_0^2\right)1.$$
$`(4.20)`$
We now consider separately the cases $`tt_0`$ and $`tt_0`$.
For $`tt_0`$, we choose $`H=h_0`$, so that $`m_21`$ and $`m_1h_1(t_0)`$. Furthermore $`z_0=b_0`$. Then (4.19) reduces to (4.12) and (4.20) reduces to (4.11).
For $`tt_0`$, we choose $`H=1`$, so that $`m_2=t_0^\gamma |\rho ^{}(t_0)|^1h_0(t_0)`$, $`z_0=b_0h_0(t_0)`$ and $`m_1=t^2|\rho ^{}(t)|^1`$. Then (4.19) reduces to (4.14) and (4.20) reduces to (4.13) with $`t`$ replaced by $`T`$.
$``$$``$
Remark 4.1. The optimal time decay results in Lemma 4.1 are obtained by saturating the conditions on $`h_0`$ and $`h_1`$, namely by taking $`h_0=t^\gamma |\rho ^{}|^1`$, which we have already assumed to be nondecreasing, and $`h_1=t^2|\rho ^{}|^1h_0=t^{2\gamma }|\rho ^{}|^2`$, in so far as that function is nonincreasing, a property that we could have (but have not) assumed. We have stated Lemma 4.1 with inequalities instead of the previous special choices, in order not to hide the flexibility allowed by the proof, and we shall proceed in the same way in all subsequent similar estimates. In the special case $`|\rho ^{}|=t^{1\epsilon }`$, the optimal decays obtained with the previous special choices are $`h_0=t^{1\gamma +\epsilon }`$ and $`h_1=t^{\gamma +2\epsilon }`$.
We now turn to the extension of the regularity result of Lemma 3.6.
Lemma 4.2. Let $`k`$, $`\mathrm{}`$ satisfy (3.48) and let $`\overline{k}`$, $`\overline{\mathrm{}}`$ satisfy $`\overline{k}k=\overline{\mathrm{}}\mathrm{}0`$. Let $`1Tt_0<\mathrm{}`$ and let $`(w,\phi )𝒳_{\rho ,loc}^{\overline{k},\overline{\mathrm{}}}([T,\mathrm{}))`$ be a solution of the system (2.11)-(2.12). Let $`h_0`$ and $`h_1`$ be as in Lemma 4.1 and assume that $`(w,\phi )`$ satisfies
$$|w(t)|_ka,|\phi (t)|_{\mathrm{}}bh_0(t)\text{or}|\phi (t)|_{\mathrm{}}bh_0(tt_0).$$
$`(4.21)`$
(1) There exist constants $`c`$ and $`C`$ such that if
$$(b+a^2)h_1(t_0)c,$$
$`(4.22)`$
then $`(w,h_0^1\phi )𝒳_\rho ^{\overline{k},\overline{\mathrm{}}}([t_0,\mathrm{}))`$ and $`(w,\phi )`$ satisfies the estimates
$$\{\begin{array}{c}Y(w;[t_0,\mathrm{}),1,\overline{k})C\left\{|w(t_0)|_{\overline{k}}+ah_1(t_0)h_0(t_0)^1|\phi (t_0)|_\overline{\mathrm{}}\right\},\hfill \\ \\ Z(\phi ;[t_0,\mathrm{}),h_0,\overline{\mathrm{}})C\left\{h_0(t_0)^1\right|\phi (t_0)|_\overline{\mathrm{}}+a|w(t_0)|_{\overline{k}}\}.\hfill \end{array}$$
$`(4.23)`$
(2) In the case where $`|\phi (t)|_{\mathrm{}}bh_0(tt_0)`$, there exist constants $`c`$ and $`C`$ such that if
$$(b+a^2)T^2|\rho ^{}(T)|^1h_0(t_0)c,$$
$`(4.24)`$
then $`(w,\phi )`$ satisfies the estimates
$$\{\begin{array}{c}Y(w;[t,t_0],1,\overline{k})C\left\{|w(t_0)|_{\overline{k}}+at^2|\rho ^{}|^1|\phi (t_0)|_\overline{\mathrm{}}\right\}\text{for all}t[T,t_0],\hfill \\ \\ Z(\phi ;[T,t_0],1,\overline{\mathrm{}})C\left\{|\phi (t_0)|_\overline{\mathrm{}}+ah_0(t_0)|w(t_0)|_{\overline{k}}\right\}.\hfill \end{array}$$
$`(4.25)`$
(3) In the case where $`|\phi (t)|_{\mathrm{}}bh_0(t)`$, there exist constants $`c`$ and $`C`$ such that if
$$(b+a^2)h_1(T)c,$$
$`(4.26)`$
then $`(w,\phi )`$ satisfies the estimates
$$\{\begin{array}{c}Y(w;[t,t_0],h_0^1,\overline{k})C\left\{h_0(t_0)\right|w(t_0)|_{\overline{k}}+ah_1(t)|\phi (t_0)|_\overline{\mathrm{}}\}\text{for all}t[T,t_0],\hfill \\ \\ Z(\phi ;[T,t_0],1,\overline{\mathrm{}})C\left\{|\phi (t_0)|_\overline{\mathrm{}}+ah_0(t_0)|w(t_0)|_{\overline{k}}\right\}.\hfill \end{array}$$
$`(4.27)`$
Proof. The proof follows closely that of Lemma 4.1. Let $`h_2`$ and $`h_3`$ be $`𝒞^1`$ positive functions of time, increasing for $`tt_0`$ and decreasing for $`tt_0`$, and let $`\stackrel{~}{w}=h_2^1w`$, $`\stackrel{~}{\phi }=h_3^1\phi `$. From Lemma 3.6 we obtain for $`t\underset{<}{>}t_0`$
$$_t|\stackrel{~}{w}|_{\overline{k}}^2\underset{>}{<}2\rho ^{}|\stackrel{~}{w}|_{\overline{k}+\nu /2}^2\pm Ct^2\left\{|\stackrel{~}{w}|_{\overline{k}+\nu /2}^2|\phi |_{\mathrm{}}+h_3h_2^1|\stackrel{~}{w}|_{\overline{k}+\nu /2}|\stackrel{~}{\phi }|_{\overline{\mathrm{}}+\nu /2}|w|_k\right\},$$
$`(4.28)`$
$$_t|\stackrel{~}{\phi }|_\overline{\mathrm{}}^2\underset{>}{<}2\rho ^{}|\stackrel{~}{\phi }|_{\overline{\mathrm{}}+\nu /2}^2\pm Ct^2|\stackrel{~}{\phi }|_{\overline{\mathrm{}}+\nu /2}^2|\phi |_{\mathrm{}}\pm Ct^\gamma h_3^1h_2|\stackrel{~}{\phi }|_{\overline{\mathrm{}}+\nu /2}|\stackrel{~}{w}|_{\overline{k}+\nu /2}|w|_k$$
$`(4.29)`$
where we have omitted the terms containing $`h_2^{}`$ and $`h_3^{}`$. We define $`y_{(1)}=y_{(1)}(w;I,h_2,\overline{k})`$ and $`z_{(1)}=z_{(1)}(\phi ;I,h_3,\overline{\mathrm{}})`$ where $`I=[t_0,t]`$ for $`tt_0`$ and $`I=[t,t_0]`$ for $`tt_0`$. Proceeding as in the proof of Lemma 4.1, we obtain from (4.28) (4.29)
$$\{\begin{array}{c}y^2y_1^2y_0^2+m_0y_1^2+m_1ay_1z_1\hfill \\ \\ z^2z_1^2z_0^2+m_0z_1^2+m_2ay_1z_1\hfill \end{array}$$
$`(4.30)`$
where $`y_0=|\stackrel{~}{w}(t_0)|_{\overline{k}}`$, $`z_0=|\stackrel{~}{\phi }(t_0)|_\overline{\mathrm{}}`$,
$$m_0=C\underset{tI}{Sup}t^2|\rho ^{}|^1|\phi (t)|_{\mathrm{}},$$
$$m_1=C\underset{tI}{Sup}t^2|\rho ^{}|^1h_3h_2^1,m_2=C\underset{tI}{Sup}t^\gamma |\rho ^{}|^1h_2h_3^1.$$
$`(4.31)`$
If we can arrange that $`m_01/2`$, then $`Y=yy_1`$ and $`Z=zz_1`$ satisfy
$$\{\begin{array}{c}Y^22y_0^2+2m_1aYZ\hfill \\ \\ Z^22z_0^2+2m_2aYZ.\hfill \end{array}$$
$`(4.32)`$
By an elementary computation, one obtains from (4.32) the estimates
$$\{\begin{array}{c}Y4\left(y_0+2am_1z_0\right)\hfill \\ \\ Z4\left(z_0+2am_2y_0\right)\hfill \end{array}$$
$`(4.33)`$
under the condition
$$8a^2m_1m_21.$$
$`(4.34)`$
We now consider separately the various cases at hand.
For $`tt_0`$, we take $`h_2=1`$ and $`h_3=h_0`$, so that
$$y_0=|w(t_0)|_{\overline{k}},z_0=h_0(t_0)^1|\phi (t_0)|_\overline{\mathrm{}},$$
$$m_0Cbh_1(t_0),m_1Ch_1(t_0),m_2C.$$
The estimate (4.33) then reduces to (4.23), while the conditions $`m_01/2`$ and (4.34) recombine to yield (4.22).
For $`tt_0`$, in the case where $`|\phi (t)|_{\mathrm{}}bh_0(t_0)`$, we take $`h_2=h_3`$ = 1, so that
$$y_0=|w(t_0)|_{\overline{k}},z_0=|\phi (t_0)|_\overline{\mathrm{}},$$
$$m_0Cbt^2\rho ^1h_0(t_0),m_1Ct^2\rho ^1\text{and}m_2Ch_0(t_0),$$
thereby obtaining (4.25) under the condition (4.24).
For $`tt_0`$, in the case where $`|\phi (t)|_{\mathrm{}}bh_0(t)`$, we take $`h_2=h_0^1`$, $`h_3=1`$, so that
$$y_0=h_0(t_0)|w(t_0)|_{\overline{k}},z_0=|\phi (t_0)|_\overline{\mathrm{}},$$
$$m_0Cbh_1(t),m_1Ch_1(t)\text{and}m_2C,$$
thereby obtaining (4.27) under the condition (4.26).
$``$$``$
We now turn to the estimate of the difference of two solutions of the system (2.11)-(2.12).
Lemma 4.3. Let $`k`$, $`\mathrm{}`$ satisfy (3.48) and let $`k^{}`$, $`\mathrm{}^{}`$ satisfy (3.57). Let $`1Tt_0<\mathrm{}`$. Let $`h_0`$ and $`h_1`$ be as in Lemma 4.1 and let $`(w_i,\phi _i)`$, $`i=1,2`$, be two solutions of the system (2.11)-(2.12) such that $`(w_i,h_0^1\phi _i)𝒳_\rho ^{k,\mathrm{}}([T,\mathrm{}))`$ and such that $`(w_i,\phi _i)`$ satisfy the estimates
$$Y(w_i,[T,\mathrm{}),1,k)a,$$
$`(4.35)`$
$$Z(\phi _i,[t_0,\mathrm{}),h_0,\mathrm{})b,$$
$`(4.36)`$
and either
$$Z(\phi _i,[T,t_0),1,\mathrm{})bh_0(t_0)$$
$`(4.37)`$
or
$$Z(\phi _i,[T,t_0),h_0,\mathrm{})b.$$
$`(4.38)`$
Let $`w_\pm =w_1\pm w_2`$ and $`\phi _\pm =\phi _1\pm \phi _2`$.
(1) There exist constants $`c`$ and $`C`$ such that under the condition (4.22), $`(w_{},\phi _{})`$ satisfies the estimates
$$\{\begin{array}{c}Y(w_{};[t_0,\mathrm{}),1,k^{})C\left\{|w_{}(t_0)|_k^{}+ah_1(t_0)h_0(t_0)^1|\phi _{}(t_0)|_{\mathrm{}^{}}\right\},\hfill \\ \\ Z(\phi _{};[t_0,\mathrm{}),h_0,\mathrm{}^{})C\left\{h_0(t_0)^1\right|\phi _{}(t_0)|_{\mathrm{}^{}}+a|w_{}(t_0)|_k^{}\}.\hfill \end{array}$$
$`(4.39)`$
(2) In the case where $`\phi _i`$ satisfy (4.37), there exist constants $`c`$ and $`C`$ such that under the condition (4.24), $`(w_{},\phi _{})`$ satisfies the estimates
$$\{\begin{array}{c}Y(w_{};[t,t_0,],1,k^{})C\{|w_{}(t_0)|_k^{}+at^2|\rho ^{}|^1|\phi _{}(t_0)|_{\mathrm{}^{}}\}\text{for all}t[T,t_0],\hfill \\ \\ Z(\phi _{};[T,t_0],1,\mathrm{}^{})C\left\{|\phi _{}(t_0)|_{\mathrm{}^{}}+ah_0(t_0)|w_{}(t_0)|_k^{}\right\}.\hfill \end{array}$$
$`(4.40)`$
(3) In the case where $`\phi _i`$ satisfy (4.38), there exist constants $`c`$ and $`C`$ such that under the condition (4.26) $`(w_{},\phi _{})`$ satisfies the estimates
$$\{\begin{array}{c}Y(w_{},[t,t_0,],h_0^1,k^{})C\{h_0(t_0)|w_{}(t_0)|_k^{}+ah_1(t)|\phi _{}(t_0)|_{\mathrm{}^{}}\}\text{for all}t[T,t_0],\hfill \\ \\ Z(\phi _{},[T,t_0],1,\mathrm{}^{})C\left\{|\phi _{}(t_0)|_{\mathrm{}^{}}+ah_0(t_0)|w_{}(t_0)|_k^{}\right\}.\hfill \end{array}$$
$`(4.41)`$
(The estimates (4.39) (4.40) and (4.41) are obtained from (4.23) (4.25) and (4.27) by replacing $`(w,\phi )`$ by $`(w_{},\phi _{})`$ and $`(\overline{k},\overline{\mathrm{}})`$ by $`(k^{},\mathrm{}^{}`$).
Proof. The proof follows closely that of Lemma 4.2. Let $`h_2`$ and $`h_3`$ be $`𝒞^1`$ positive functions of time, increasing for $`tt_0`$ and decreasing for $`tt_0`$ and let $`\stackrel{~}{w}_{}=h_2^1w_{}`$, $`\stackrel{~}{\phi }_{}=h_3^1\phi _{}`$. Let $`H(t)=h_0(t)`$ in cases where (4.36) (4.38) are relevant, $`H(t)=h_0(t_0)`$ in the case where (4.37) is relevant. From Lemma 3.7 we obtain for $`t\underset{<}{>}t_0`$
$$_t|\stackrel{~}{w}_{}|_k^{}^2\underset{>}{<}2\rho ^{}|\stackrel{~}{w}_{}|_{k^{}+\nu /2}^2\pm Ct^2\left\{|\stackrel{~}{w}_{}|_{k^{}+\nu /2}^2|\phi _+|_{\mathrm{}}+|\stackrel{~}{w}_{}|_{k^{}+\nu /2}|\stackrel{~}{w}_{}|_k^{}|\phi _+|_{\mathrm{}+\nu /2}\right\}$$
$$\pm Ct^2h_3h_2^1|\stackrel{~}{w}_{}|_{k^{}+\nu /2}\left\{|\stackrel{~}{\phi }_{}|_{\mathrm{}^{}}|w_+|_{k+\nu /2}+|\stackrel{~}{\phi }_{}|_{\mathrm{}^{}+\nu /2}|w_+|_k\right\},$$
$`(4.42)`$
$$_t|\stackrel{~}{\phi }_{}|_{\mathrm{}^{}}^2\underset{>}{<}2\rho ^{}|\stackrel{~}{\phi }_{}|_{\mathrm{}^{}+\nu /2}^2\pm Ct^2\left\{|\stackrel{~}{\phi }_{}|_{\mathrm{}^{}+\nu /2}^2|\phi _+|_{\mathrm{}}+|\stackrel{~}{\phi }_{}|_{\mathrm{}^{}+\nu /2}|\stackrel{~}{\phi }_{}|_{\mathrm{}^{}}|\phi _+|_{\mathrm{}+\nu /2}\right\}$$
$$+Ct^\gamma h_2h_3^1|\stackrel{~}{\phi }_{}|_{\mathrm{}^{}+\nu /2}\left\{|\stackrel{~}{w}_{}|_k^{}|w_+|_{k+\nu /2}+|\stackrel{~}{w}_{}|_{k^{}+\nu /2}|w_+|_k\right\},$$
$`(4.43)`$
where we have omitted the terms containing $`h_2^{}`$ and $`h_3^{}`$. We define $`y_{(1)}=y_{(1)}(w_{};I,h_2,k^{})`$ and $`z_{(1)}=z_{(1)}(\phi _{};I,h_3,\mathrm{}^{})`$ where $`I=[t_0,t]`$ for $`tt_0`$ and $`I=[t,t_0]`$ for $`tt_0`$. Proceding as in the proof of Lemmas 4.1 and 4.2, we obtain from (4.42) (4.43) supplemented by (4.35)-(4.38) applied to $`(w_+,\phi _+)`$
$$\{\begin{array}{c}y^2y_1^2y_0^2+m_0by_1(y+y_1)+m_1ay_1(z+z_1)\hfill \\ \\ z^2z_1^2z_0^2+m_0bz_1(z+z_1)+m_2az_1(y+y_1)\hfill \end{array}$$
$`(4.44)`$
where
$$y_0=|\stackrel{~}{w}_{}(t_0)|_k^{},z_0=|\stackrel{~}{\phi }_{}(t_0)|_{\mathrm{}^{}},$$
$$m_0=C\underset{tI}{Sup}\left(t^2|\rho ^{}|^1H(t)\right)$$
and $`m_1`$, $`m_2`$ are defined by (4.31). From there on, the proof is identical with that of Lemma 4.2, with (4.30) replaced by (4.44).
$``$$``$
Remark 4.2. It is an unfortunate feature of Lemmas 4.2 and 4.3 that the derivation of regularity and of difference estimates requires a large time restriction (see (4.22) (4.24) (4.26)) whereas one would expect those estimates to hold for all times where the solution is a priori defined, since those estimates are linear in the higher or difference norm and are expected to follow from some kind of Gronwall’s inequality. The reason for that fact is the occurrence of integral norms in the definition of the spaces, for which we obtain algebraic inequalities which require some kind of smallness condition in order to enable us to conclude.
In practice, the conditions (4.22) (4.24) (4.26) required for those estimates to hold have the same form and the same dependence on basic parameters such as $`a_0`$, $`b_0`$ as the conditions that will be needed anyway in order to derive the a priori estimates on one single solution that are needed to solve the Cauchy problem. We shall impose all such conditions together, without any significant limitation on the range of validity of the results (see the proof of Proposition 4.1 below).
We now turn to the Cauchy problem for large time for the system (2.11)-(2.12).
Proposition 4.1. Let $`(k,\mathrm{})`$ satisfy (3.48) and in addition $`k1\nu /2`$. Let $`h_0`$ and $`h_1`$ be $`𝒞^1`$ positive functions defined in $`[1,\mathrm{})`$ with $`h_0`$ nondecreasing, $`h_1`$ nonincreasing and tending to zero at infinity, $`h_0t^\gamma |\rho ^{}|^1`$ and $`h_1t^2|\rho ^{}|^1h_0`$. Let $`a_0>0`$, $`b_0>0`$.
(1) There exists $`T_0<\mathrm{}`$, depending on $`a_0`$, $`b_0`$, such that for all $`t_0T_0`$, there exists $`Tt_0`$, depending on $`a_0`$, $`b_0`$ and $`t_0`$, such that for any $`(w_0,\phi _0)K_{\rho __0}^kY_{\rho __0}^{\mathrm{}}`$, where $`\rho _0=\rho (t_0)`$, satisfying $`|w_0|_ka_0`$, $`|\phi _0|_{\mathrm{}}h_0(t_0)^1b_0`$, the system (2.11)-(2.12) has a unique solution in the interval $`[T,\mathrm{})`$ with $`w(t_0)=w_0`$, $`\phi (t_0)=\phi _0`$, such that $`(w,h_0^1\phi )𝒳_\rho ^{k,\mathrm{}}([T,\mathrm{}))`$. One can define $`T_0`$ and $`T`$ by
$$\left(b_0+a_0^2\right)h_1(T_0)=c$$
$`(4.45)`$
$$T^2|\rho ^{}(T)|^1h_0(t_0)h_1(T_0)^1=1$$
$`(4.46)`$
and the solution $`(w,\phi )`$ satisfies the estimates
$$Y(w;[T,\mathrm{}),1,k)Ca_0,$$
$`(4.47)`$
$$Z(\phi ;[t_0,\mathrm{}),h_0,\mathrm{})C\left(b_0+a_0^2\right),$$
$`(4.48)`$
$$Z(\phi ;[T,t_0],1,\mathrm{})C\left(b_0+a_0^2\right)h_0(t_0).$$
$`(4.49)`$
(2) If $`(w_0,\phi _0)K_{\rho _0}^{\overline{k}}Y_{\rho _0}^\overline{\mathrm{}}`$ for some $`\overline{k}`$, $`\overline{\mathrm{}}`$ with $`\overline{k}k=\overline{\mathrm{}}\mathrm{}>0`$, then $`(w,h_0^1\phi )𝒳_\rho ^{\overline{k},\overline{\mathrm{}}}([T,\mathrm{}))`$, possibly after changing the constant $`c`$ in (4.45) (see Remark 4.2).
(3) The map $`(w_0,\phi _0)(w,\phi )`$ is norm continuous on the bounded sets of $`K_{\rho __0}^kY_{\rho __0}^{\mathrm{}}`$ from the norm of $`(w_0,\phi _0)`$ in $`K_{\rho __0}^k^{}Y_{\rho __0}^{\mathrm{}^{}}`$ to the norm of $`(w,h_0^1\phi )`$ in $`𝒳_\rho ^{k^{\prime \prime },\mathrm{}^{\prime \prime }}([T,\mathrm{}))`$ for $`k^{}\nu /2`$, $`kk^{}=\mathrm{}\mathrm{}^{}1\nu `$, $`kk^{\prime \prime }=\mathrm{}\mathrm{}^{\prime \prime }>0`$. Furthermore the same map is continuous from the same topology on $`(w_0,\phi _0)`$ to the weak-$``$ topology of $`(w,h_0^1\phi )`$ in $`𝒳_\rho ^{k,\mathrm{}}([T,\mathrm{}))`$.
Proof. Part (1). The proof proceeds in several steps using a parabolic regularization and a limiting procedure. We consider first the case $`tt_0`$. We shall then indicate briefly the modifications needed in the case $`tt_0`$. We shall need the function $`\stackrel{~}{w}`$ defined by $`\stackrel{~}{w}(t)=U(1/t)w(t)`$. Since the operator $`U(1/t)`$ is unitary in $`K_\rho ^k`$ for all $`k`$ and $`\rho `$, all subsequent norm estimates for $`\stackrel{~}{w}`$ in $`K_\rho ^k`$ are identical with the same estimates for $`w`$.
Step 1. Parabolic regularization and local resolution.
We rewrite the system (2.11)-(2.12) in the equivalent form
$$\{\begin{array}{c}_t\stackrel{~}{w}=\left(2t^2\right)^1U(1/t)\left(2s+(s)\right)U(1/t)^{}\stackrel{~}{w}G_1(\stackrel{~}{w},s),\hfill \\ \\ _t\phi =\left(2t^2\right)^1|s|^2+t^\gamma g_0\left(U(1/t)^{}\stackrel{~}{w}\right)G_2(\stackrel{~}{w},s).\hfill \end{array}$$
$`(4.50)`$
We introduce a parabolic regularization and consider the regularized system
$$\{\begin{array}{c}_t\stackrel{~}{w}=\theta \mathrm{\Delta }\stackrel{~}{w}+G_1(\stackrel{~}{w},s)\hfill \\ \\ _t\phi =\theta \mathrm{\Delta }\phi +G_2(\stackrel{~}{w},s)\hfill \end{array}$$
$`(4.51)`$
with $`\theta >0`$. We also regularize the initial data $`(\stackrel{~}{w}_0,\phi _0)`$ at time $`t_0`$ to $`(\overline{\stackrel{~}{w}}_0,\overline{\phi }_0)X_{\rho __0}^{\overline{k},\overline{\mathrm{}}}K_{\rho __0}^{\overline{k}}Y_{\rho _0}^\overline{\mathrm{}}`$ with $`\rho _0=\rho (t_0)`$ , $`\overline{k}k(1+\nu /2)`$ , $`\overline{\mathrm{}}\mathrm{}=\overline{k}k`$. For the purpose of proving Part (1), one can take equality in the previous condition, namely $`\overline{k}=k(1+\nu /2)`$ and in particular that second regularization is unnecessary if $`k1+\nu /2`$. We continue the argument with general $`(\overline{k},\overline{\mathrm{}})`$ because it will be useful for the proof of Part (2).
We rewrite the Cauchy problem for the system (4.51) in the integral form
$$\left(\genfrac{}{}{0pt}{}{\stackrel{~}{w}}{\phi }\right)(t)=V_\theta (tt_0)\left(\genfrac{}{}{0pt}{}{\overline{\stackrel{~}{w}}_0}{\overline{\phi }_0}\right)+_{t_0}^t𝑑t^{}V_\theta (tt^{})\left(\genfrac{}{}{0pt}{}{G_1(\stackrel{~}{w},s)}{G_2(\stackrel{~}{w},s)}\right)(t^{})$$
$`(4.52)`$
where $`V_\theta (t)\mathrm{exp}(\theta t\mathrm{\Delta })`$ is a contraction in $`X_{\rho __0}^{\overline{k},\overline{\mathrm{}}}`$ and satisfies the bound
$$V_\theta (t);(X_{\rho __0}^{\overline{k},\overline{\mathrm{}}})C(\theta t)^{1/2}.$$
$`(4.53)`$
By (3.23) (3.24) (3.26) (3.27) (3.28) (3.29) of Lemma 3.4, we estimate
$$\{\begin{array}{c}|G_1(\stackrel{~}{w},s)|_{\overline{k}1}Ct^2\left(|\stackrel{~}{w}|_{\overline{k}}|\phi |_{\overline{\mathrm{}}1}+|\stackrel{~}{w}|_{\overline{k}1}|\phi |_\overline{\mathrm{}}\right)\hfill \\ \\ |G_2(\stackrel{~}{w},s)|_{\overline{\mathrm{}}1}Ct^2|\phi |_\overline{\mathrm{}}|\phi |_{\overline{\mathrm{}}1}+Ct^\gamma |\stackrel{~}{w}|_{\overline{k}}|\stackrel{~}{w}|_{\overline{k}1}\hfill \end{array}$$
$`(4.54)`$
under the conditions (which follow from (3.48))
$$\overline{\mathrm{}}>n/2,\overline{\mathrm{}}\overline{k}10,\overline{k}\overline{\mathrm{}}+\lambda 1,2\overline{k}>\overline{\mathrm{}}+\lambda +n/2,$$
where $`\lambda =\mu n+2`$. In (4.54) the various norms are taken with constant $`\rho (t)\rho _0`$. By a standard contraction argument, the system (4.52) has a unique solution
$$(\stackrel{~}{w},\phi )𝒞([t_0,t_0+\tau )],X_{\rho __0}^{\overline{k},\overline{\mathrm{}}})$$
$`(4.55)`$
where one can take
$$\tau =C\theta \left(t_0^2h_0(t_0)+t_0^\gamma h_0(t_0)^1\right)^2\left(\overline{a}_0+\overline{b}_0\right)^2$$
$`(4.56)`$
with
$$\overline{a}_0=|\overline{\stackrel{~}{w}}_0|_{\overline{k}},\overline{b}_0=h_0(t_0)^1|\overline{\phi }_0|_\overline{\mathrm{}}.$$
Furthermore, from the estimates
$$_t|\stackrel{~}{w}|_{\overline{k}}^2+2\theta |\stackrel{~}{w}|_{\overline{k}}^2=2\mathrm{R}\mathrm{e}<|\xi |^{\overline{k}+1}f\widehat{\stackrel{~}{w}},|\xi |^{\overline{k}1}f\widehat{G_1}(\stackrel{~}{w},\phi )>+\text{lower order terms}$$
$$|\stackrel{~}{w}|_{\overline{k}+1}A_1(t,|\stackrel{~}{w}|_{\overline{k}},|\phi |_\overline{\mathrm{}}),$$
$`(4.57)`$
$$_t|\phi |_\overline{\mathrm{}}^2+2\theta |\phi |_\overline{\mathrm{}}^2=2<|\xi |^{\overline{\mathrm{}}+3}f\widehat{\phi },|\xi |^{\overline{\mathrm{}}+1}f\widehat{G_2}(\stackrel{~}{w},\phi )>+\text{lower order terms}$$
$$|\phi |_{\overline{\mathrm{}}+1}A_2(t,|\stackrel{~}{w}|_{\overline{k}},|\phi |_\overline{\mathrm{}})$$
$`(4.58)`$
for some estimating functions $`A_1`$ and $`A_2`$, it follows that
$$(\stackrel{~}{w},\phi )L^2([t_0,t_0+\tau ],X_{\rho __0}^{\overline{k}+1,\overline{\mathrm{}}+1}).$$
$`(4.59)`$
The estimates (4.57) (4.58) are derived with the help of a regularization $`j_\epsilon `$, which we have omitted for brevity, and of the estimate (4.54).
Step 2. Uniform estimates and globalisation.
¿From the regularity conditions (4.55) (4.59), from the fact that $`\rho `$ is decreasing and from Lemmas 4.1 and 4.2, especially (4.12) and (4.23), it follows that $`(\stackrel{~}{w},\phi )𝒳_\rho ^{\overline{k},\overline{\mathrm{}}}([t_0,t_0+\tau ])`$ and that $`(\stackrel{~}{w},\phi )`$ satisfies the estimates
$$\{\begin{array}{c}Y(\stackrel{~}{w};I,1,k)Ca_0\hfill \\ \\ Z(\phi ;I,h_0,\mathrm{})C\left(b_0+a_0^2\right),\hfill \end{array}$$
$`(4.60)`$
$$\{\begin{array}{c}Y(\stackrel{~}{w};I,1,\overline{k})C\left(\overline{a}_0+h_1(t_0)a_0\overline{b}_0\right)\hfill \\ \\ Z(\phi ;I,h_0,\overline{\mathrm{}})C\left(\overline{b}_0+a_0\overline{a}_0\right),\hfill \end{array}$$
$`(4.61)`$
for $`I=[t_0,t_0+\tau ]`$ and $`t_0T_0`$, under the condition (4.45) which ensures (4.11) (4.22). The estimates (4.60) (4.61) are uniform in $`\theta `$. From (4.56), from (4.60) (4.61) for general $`I`$ and from the fact that $`\rho `$ is decreasing, it follows by a minor modification of a standard globalisation argument that $`(\stackrel{~}{w},\phi )`$ can be continued to a solution of the system (4.51) such that $`(\stackrel{~}{w},h_0^1\phi )`$ belongs to $`𝒳_\rho ^{\overline{k},\overline{\mathrm{}}}([t_0,\mathrm{}))`$ and that $`(\stackrel{~}{w},\phi )`$ satisfies (4.60) (4.61) with $`I=[t_0,\mathrm{})`$.
Step 3. Limiting procedure.
We now take the limits $`\theta 0`$ and $`(\overline{\stackrel{~}{w}}_0,\overline{\phi }_0)(\stackrel{~}{w}_0,\phi _0)`$ in that order. We first keep $`(\overline{\stackrel{~}{w}}_0,\overline{\phi }_0)`$ fixed and consider two solutions $`(\stackrel{~}{w}_1,\phi _1)`$ and $`(\stackrel{~}{w}_2,\phi _2)`$ with $`(\stackrel{~}{w}_i,h_0^1\phi _i)𝒳_{\rho __0}^{\overline{k},\overline{\mathrm{}}}([t_0,\mathrm{}))`$ as obtained in Step 2, corresponding to two values $`\theta _1`$ and $`\theta _2`$. We estimate the difference $`(\stackrel{~}{w}_{},\phi _{})=(\stackrel{~}{w}_1\stackrel{~}{w}_2,\phi _1\phi _2)`$ by a minor variation of Lemma 4.3 with $`(k,\mathrm{})`$ replaced by $`(\overline{k},\overline{\mathrm{}})`$ and $`(k^{},\mathrm{}^{})=(\overline{k}1,\overline{\mathrm{}}1)`$, under the condition (4.22) which follows from $`t_0T_0`$ and from (4.45) (4.60), possibly after changing the constant $`c`$. More precisely in the proof of (4.39), we take the initial condition $`\stackrel{~}{w}_{}(t_0)=0`$, $`\phi _{}(t_0)=0`$, but we have an additional term coming from the parabolic regularization in the analogue of (4.42) (4.43), namely
$$\{\begin{array}{c}_t|\stackrel{~}{w}_{}|_{\overline{k}1}^22|\stackrel{~}{w}_{}|_{\overline{k}}\{\theta _1|\stackrel{~}{w}_1|_{\overline{k}}+\theta _2|\stackrel{~}{w}_2|_{\overline{k}})+\text{previous terms},\hfill \\ \\ _t|\phi _{}|_{\overline{\mathrm{}}1}^22|\phi _{}|_\overline{\mathrm{}}\left(\theta _1|\phi _1|_\overline{\mathrm{}}+\theta _2|\phi _2|_\overline{\mathrm{}}\right)+\text{previous terms}.\hfill \end{array}$$
$`(4.62)`$
(The tildas in (4.62) and in (4.42) (4.43) have different meanings, but this has no implication on the argument). From (4.61) and (4.62), for any $`t_1`$ with $`t_0<t_1<\mathrm{}`$, we obtain estimates of the type (4.44) for the quantities $`y_{(1)}=y_{(1)}(w_{};[t_0,t_1],1,\overline{k}1)`$ and $`z_{(1)}=z_{(1)}(\phi _{};[t_0,t_1],h_0,\overline{\mathrm{}}1)`$ where now
$$y_0^2z_0^2|t_1t_0|(\theta _1+\theta _2)A(t_0,\overline{a}_0,\overline{b}_0).$$
$`(4.63)`$
Therefore, by the same argument as in the proof of (4.39)
$$\{\begin{array}{c}Y(\stackrel{~}{w}_{};[t_0,t_1],1,\overline{k}1)C\left(y_0+a_0h_1(t_0)z_0\right)\hfill \\ \\ Z(\phi _{};[t_0,t_1],h_0,\overline{\mathrm{}}1)C\left(z_0+a_0y_0\right),\hfill \end{array}$$
$`(4.64)`$
which implies that the solution $`(\stackrel{~}{w}_\theta ,\phi _\theta )`$ associated with $`\theta `$ converges in norm in $`𝒳_\rho ^{\overline{k}1,\overline{\mathrm{}}1}([t_0,t_1])`$ when $`\theta 0`$ for all $`t_1t_0`$. Furthermore, the limit $`(\stackrel{~}{w},\phi )`$ is such that $`(\stackrel{~}{w},h_0^1\phi )𝒳_\rho ^{\overline{k},\overline{\mathrm{}}}([t_0,\mathrm{}))`$ and $`(\stackrel{~}{w},\phi )`$ satisfies (4.47) (4.48). This follows from the bound (4.61) with $`I=[t_0,\mathrm{})`$, which is uniform in $`\theta `$, and from the previous convergence by standard compactness arguments, except for the strong continuity in time. The latter follows from the weak continuity, which also follows from a compactness argument, and from the fact that the $`K_\rho ^{\overline{k}}Y_\rho ^\overline{\mathrm{}}`$ norm of $`(\stackrel{~}{w},\phi )`$ is (absolutely) continuous in $`t`$ by Lemma 3.5 with $`k`$, $`\mathrm{}`$ replaced by $`\overline{k}`$, $`\overline{\mathrm{}}`$. The limit obviously satisfies the system (4.50).
We let now $`(\overline{\stackrel{~}{w}}_0,\overline{\phi }_0)`$ tend to $`(\stackrel{~}{w}_0,\phi _0)`$ in $`K_{\rho __0}^kY_{\rho __0}^{\mathrm{}}`$ (that step is not needed if $`\overline{k}=k1+\nu /2`$). Let $`(\overline{\stackrel{~}{w}}_{0i},\overline{\phi }_{0i})K_{\rho __0}^{\overline{k}}Y_{\rho __0}^\overline{\mathrm{}}`$, $`i=1,2`$, be two sets of regularized initial conditions and let $`(\stackrel{~}{w}_i,\phi _i)`$ be the solutions of the system (4.50) obtained previously. The difference $`(\stackrel{~}{w}_{},\phi _{})=(\stackrel{~}{w}_1\stackrel{~}{w}_2,\phi _1\phi _2)`$ is then estimated by (4.39) with $`a`$ replaced by $`a_0`$ as follows from (4.60), under the condition $`t_0T_0`$ and (4.45) as previously. One can (but need not) take $`k^{}=k1+\nu `$, $`\mathrm{}^{}=\mathrm{}1+\nu `$. This implies that the solution $`(\overline{\stackrel{~}{w}},\overline{\phi })`$ associated with $`(\overline{\stackrel{~}{w}}_0,\overline{\phi }_0)`$ converges in the norm of $`(\overline{\stackrel{~}{w}},h^1\overline{\phi })`$ in $`𝒳_\rho ^{k^{},\mathrm{}^{}}([t_0,\mathrm{}))`$ when $`(\overline{\stackrel{~}{w}}_0,\overline{\phi }_0)`$ converges to $`(\stackrel{~}{w}_0,\phi _0)`$ in the norm of $`K_{\rho __0}^k^{}Y_{\rho __0}^{\mathrm{}^{}}`$ on the bounded sets of $`K_{\rho __0}^kY_{\rho __0}^{\mathrm{}}`$ (and a fortiori in the norm of $`K_{\rho __0}^kY_{\rho __0}^{\mathrm{}}`$). Let $`(\stackrel{~}{w},\phi )`$ be the limit. By the same arguments as above, $`(\stackrel{~}{w},h_0^1\phi )𝒳_\rho ^{k,\mathrm{}}([t_0,\mathrm{}))`$ and $`(\stackrel{~}{w},\phi )`$ satisfies the system (4.50) and the estimates (4.47) (4.48).
Step 4. Uniqueness follows immediately from Lemma 4.3, part (1).
We now turn to the case $`tt_0`$. The proof proceeds exactly in the same way, with Parts 1 of Lemmas 4.1, 4.2 and 4.3 replaced by Parts 2 of the same Lemmas. In the same way as before, (4.14) implies that (4.24) follows from (4.13), possibly after a change of constant $`c`$. With $`T_0`$ defined by (4.45) and possibly with another change of constant $`c`$, the condition (4.13) follows from (4.46).
Part (2). If $`\overline{k}1+\nu /2`$, the result follows from the proof of Part (1) with the second limiting procedure omitted. If $`\overline{k}<1+\nu /2`$, the result follows from the proof of Part (1) with $`\overline{k}`$ replaced by $`1+\nu /2`$ and with the second limiting procedure going down from $`1+\nu /2`$ to $`\overline{k}`$ instead of going down from $`1+\nu /2`$ to $`k`$.
Note that in the previous proof, the constant $`c`$ in (4.45) comes from a successive application of Lemma 4.1, especially from (4.11) (4.13), and of Lemmas 4.2 and 4.3, especially from (4.22) (4.24). In particular that constant depends on $`(k,\mathrm{})`$ and on $`(\overline{k},\overline{\mathrm{}})`$ (the pair $`(k^{},\mathrm{}^{})`$ in the applications of Lemma 4.3 is chosen as a function of $`(k,\mathrm{})`$ or $`(\overline{k},\overline{\mathrm{}}))`$. In the proof of Part (1), one can in addition choose $`\overline{k}=k(1+\nu /2)`$, $`\overline{k}k=\overline{\mathrm{}}\mathrm{}`$, and the constant $`c`$ therefore depends only on $`(k,\mathrm{})`$. In contrast with that, in the proof of Part (2), the pairs $`(k,\mathrm{})`$ and $`(\overline{k},\overline{\mathrm{}})`$ are independent, and the constant $`c`$ in (4.45) needed for Part (2) to hold may (and is expected to) depend on both $`(k,\mathrm{})`$ and $`(\overline{k},\overline{\mathrm{}})`$. The crucial information contained in Part (2) is the fact that (4.45) does not involve $`|w_0|_{\overline{k}}`$, $`|\phi _0|_\overline{\mathrm{}}`$, but only $`|w_0|_k`$ and $`|\phi _0|_{\mathrm{}}`$.
Part (3). Continuity with respect to initial data follows from Lemma 4.3 parts 1 and 2, and from the a priori estimates (4.47) (4.48) (4.49) by interpolation and compactness arguments.
$``$$``$
We conclude this section with two properties of the behaviour at infinity of solutions of the system (2.11)-(2.12) in spaces $`𝒳_\rho ^{k,\mathrm{}}`$. There is no initial time involved in those properties and $`\rho `$ is only required to satisfy a suitable monotony condition at infinity. The first property is the existence of a limit for $`w(t)`$ as $`t\mathrm{}`$. It applies in particular to the solutions constructed in Proposition 4.1. There is a large flexibility on the assumptions under which such a limit exist. We shall give another example in the next section in a different context.
Proposition 4.2. Let $`T1`$, $`\rho _{\mathrm{}}0`$ and
$$\rho (t)=\rho _{\mathrm{}}+_t^{\mathrm{}}𝑑t_1|\rho ^{}(t_1)|$$
for $`tT`$. Let $`k`$, $`\mathrm{}`$ satisfy
$$\mathrm{}>n/21,1\nu /2k\mathrm{}+1.$$
$`(4.65)`$
Let $`h_0`$ and $`h_1`$ be as in Proposition 4.1. Let $`(w,\phi )`$ be a solution of the system (2.11)-(2.12) such that $`(w,h_0^1\phi )𝒳_\rho ^{k,\mathrm{}}([T,\mathrm{}))`$. Let
$$a=Y(w;[T,\mathrm{}),1,k),b=Z(\phi ;[T,\mathrm{}),h_0,\mathrm{}).$$
Then there exists $`w_+K_\rho _{\mathrm{}}^k`$ such that $`w(t)`$ tends to $`w_+`$ strongly in $`K_\rho _{\mathrm{}}^k^{}`$ for $`k^{}<k`$ and weakly in $`K_\rho _{\mathrm{}}^k`$ when $`t\mathrm{}`$. Furthermore, the following estimates hold :
$$|w_+|_kw_+;K_\rho _{\mathrm{}}^ka,$$
$`(4.66)`$
$$Y(\stackrel{~}{w}\stackrel{~}{w}(t_1);[t_1,\mathrm{}),1,k^{})Cabh_1(t_1),$$
$`(4.67)`$
$$\stackrel{~}{w}(t)w_+;K_\rho _{\mathrm{}}^k^{}Cabh_1(t),$$
$`(4.68)`$
for $`k^{}=k1+\nu `$, for all $`t`$, $`t_1T`$, and with $`\stackrel{~}{w}(t)=U(1/t)w(t)`$.
Proof. We first prove (4.67). Let $`tt_1`$ and $`w_1(t)=\stackrel{~}{w}(t)\stackrel{~}{w}(t_1)`$ so that by (4.50)
$$_tw_1=(2t^2)^1U(1/t)(2s+(s))U(1/t)^{}\stackrel{~}{w}.$$
Using (3.12) and (3.23) (3.24) from Lemma 3.4 with $`k^{}=k1+\nu `$ , $`m=k^{}\nu /2`$, we obtain
$$_t|w_1|_k^{}^22\rho ^{}|w_1|_{k^{}+\nu /2}^2Ct^2|w_1|_{k^{}+\nu /2}\left\{|w|_{k+\nu /2}|\phi |_{\mathrm{}}+|\phi |_{\mathrm{}+\nu /2}|w|_k\right\}$$
$`(4.69)`$
under the conditions $`k^{}\nu /2`$ , $`\mathrm{}>n/21`$ , $`\mathrm{}k1`$, which reduce to (4.65). Let $`y_{(1)}=y_{(1)}(w_1;[t_1,t],1,k^{})`$. Integrating (4.69) over time in the same way as in Lemma 4.1, we obtain
$$y^2y_1^2Caby_1\left\{\underset{t^{}[t_1,t]}{Sup}t^2|\rho ^{}(t^{})|^1h_0(t^{})\right\}Caby_1h_1(t_1)$$
from which (4.67) follows. From (4.67) and from the fact that $`\rho `$ is decreasing, it follows that $`w(t)`$ has a strong limit $`w_+K_\rho _{_{\mathrm{}}}^k^{}`$ and that (4.68) holds. By standard compactness and interpolation arguments, $`w_+K_\rho _{\mathrm{}}^k`$, $`w_+`$ satisfies (4.66) and $`w(t)`$ converges to $`w_+`$ weakly in $`K_\rho _{_{\mathrm{}}}^k`$ and strongly in $`K_\rho _{_{\mathrm{}}}^k^{}`$ for $`k^{}<k`$.
$``$$``$
The second property is a uniqueness property for solutions with suitable restrictions on the behaviour at infinity. It will be used in Section 6. Since it corresponds to a situation with infinite initial time, it requires $`\rho `$ to be increasing.
Proposition 4.3. Let $`T1`$, $`\rho _{\mathrm{}}>0`$ and
$$\rho (t)=\rho _{\mathrm{}}_t^{\mathrm{}}𝑑t_1|\rho ^{}(t_1)|$$
with $`\rho (t)0`$ for $`tT`$. Let $`(k,\mathrm{})`$ satisfy (3.48) and $`k1\nu /2`$. Let $`h_0`$ and $`h_1`$ be as in Proposition 4.1. Let $`(w_i,\phi _i)`$, $`i=1,2`$ be two solutions of the system (2.11)-(2.12) such that $`(w_i,h_0^1\phi _i)𝒳_\rho ^{k,\mathrm{}}([T,\mathrm{}))`$ and such that
$$|w_1(t)w_2(t)|_k^{}h_0(t)0,|\phi _1(t)\phi _2(t)|_{\mathrm{}^{}}0$$
when $`t\mathrm{}`$, for some $`(k^{},\mathrm{}^{})`$ such that $`\nu /2k^{}k1+\nu `$, $`\mathrm{}\mathrm{}^{}=kk^{}`$. Let
$$a=\underset{i}{Max}Y(w_i;[T,\mathrm{}),1,k),b=\underset{i}{Max}Z(\phi _i;[T,\mathrm{}),h_0,\mathrm{}).$$
Then there exists a constant $`c`$ such that if (4.26) holds, then $`(w_1,\phi _1)=(w_2,\phi _2)`$.
Proof. The result follows immediately from Lemma 4.3, part (3) by taking the limit $`t_0\mathrm{}`$ in (4.41).
$``$$``$
## 5 Asymptotics of solutions of the auxiliary system
In this section, we continue the study of the asymptotic properties of the solutions of the auxiliary system (2.11)-(2.12) obtained in Proposition 4.1. We have already proved the existence of a limit $`w_+`$ for $`w(t)`$ as $`t\mathrm{}`$ for such solutions. Under suitable additional regularity assumptions in the form of stronger lower bounds on $`(k,\mathrm{})`$, we shall obtain estimates on the asymptotic behaviour in time of the asymptotic functions $`(w_m,\phi _m)`$ and $`(W_m,\varphi _m)`$ defined by (2.14)-(2.15) and estimates on the remainders $`(q_{p+1},\psi _{p+1})=(wW_p,\phi \varphi _p)`$, also defined by (2.14)-(2.15), eventually leading to the existence of an asymptotic state for the phase $`\phi `$ in the form of a limit $`\psi _+`$ for $`\psi _{p+1}`$ as $`t\mathrm{}`$ for sufficiently large $`t`$.
At a technical level however, the situation here differs significantly from that in Section 4. We are no longer trying to solve the system (2.11)-(2.12), but instead we assume a solution of that system to be given, and we estimate successively the asymptotic functions $`(w_m,\phi _m)`$ and the remainders $`(q_m,\psi _m)`$, which are defined by triangular systems of equations. As a consequence, there is no need to control the loss of derivatives as in (2.11)-(2.12), and we can simply let that loss accumulate in the solution of the triangular system. Therefore we no longer need a time dependent $`\rho `$ (see (3.12)), integral norms in $`𝒳_\rho ^{k,\mathrm{}}`$, and integration by parts (see (3.22) (3.25)). We do not even need Gevrey spaces, and we could use instead ordinary Sobolev spaces, in the same way as in II. We shall of course nevertheless keep Gevrey spaces in order to make contact with Section 4, but in all this section we shall take $`\rho `$ to be constant. Instead of the spaces $`𝒳_\rho ^{k,\mathrm{}}`$ defined by (4.2), we shall use the simpler spaces
$$𝒴_\rho ^{k,\mathrm{}}(I)=(𝒞L^{\mathrm{}})(I,K_\rho ^kY_\rho ^{\mathrm{}}).$$
$`(5.1)`$
The contact with Section 4, in particular the applicability of the results of this section to the solutions obtained in Proposition 4.1, will be achieved through the fact that if $`\rho `$ is defined in $`[t_0,\mathrm{})`$ by (4.1) or equivalently by
$$\rho (t)=\rho _{\mathrm{}}+_t^{\mathrm{}}𝑑t_1|\rho ^{}(t_1)|,$$
$`(5.2)`$
then
$$𝒳_\rho ^{k,\mathrm{}}([t_0,\mathrm{}))𝒴_\rho _{\mathrm{}}^{k,\mathrm{}}([t_0,\mathrm{})).$$
$`(5.3)`$
In all the estimates of this section, the function $`f`$ occurring in the definition (3.8) (3.9) of the spaces plays no role whatsoever, and is consistently eliminated from the proofs by using the submultiplicativity property (3.3). As a consequence all the estimates are uniform in (actually independent of) $`\rho `$ and $`\nu `$, and no assumption is made connecting $`\nu `$ to other parameters such as $`\mu `$ or $`(k,\mathrm{})`$ : we only assume $`0\nu 1`$ and $`\rho 0`$.
As a preliminary result, we give an existence result of the limit $`w_+`$ of $`w(t)`$ as $`t\mathrm{}`$ for a solution of the system (2.11)-(2.12). That result is a simplified version of Proposition 4.2 appropriate to the present context.
Proposition 5.1. Let $`(k,\mathrm{})`$ satisfy
$$k1,\mathrm{}0,k+\mathrm{}>n/2.$$
$`(5.4)`$
Let $`1T<\mathrm{}`$ and let $`h_0`$ be a $`𝒞^1`$ positive nondecreasing function defined in $`[T,\mathrm{})`$ with $`t^2h_0L^1([T,\mathrm{}))`$. Let $`(w,\phi )`$ be a solution of the system (2.11)-(2.12) such that $`(w,h_0^1\phi )𝒴_\rho ^{k,\mathrm{}}([T,\mathrm{}))`$. Then there exists $`w_+K_\rho ^k`$ such that $`w(t)`$ tends to $`w_+`$ strongly in $`K_\rho ^k^{}`$ for $`k^{}<k`$ and weakly in $`K_\rho ^k`$. Furthermore, the following estimate holds :
$$|\stackrel{~}{w}(t)w_+|_k^{}Cabh(t)$$
$`(5.5)`$
for $`0k^{}k1`$, $`k^{}\mathrm{}`$, $`k^{}<k+\mathrm{}n/2`$, where
$$a=\underset{t}{Sup}|w(t)|_k,b=\underset{t}{Sup}h_0^1(t)|\phi (t)|_{\mathrm{}},$$
$`(5.6)`$
$`h(t)=_t^{\mathrm{}}𝑑t_1t_1^2h_0(t_1)`$ and $`\stackrel{~}{w}(t)=U(1/t)w(t)`$.
Proof. By (3.23) (3.24), we estimate
$$|_t(\stackrel{~}{w}(t)\stackrel{~}{w}(t_1))|_k^{}Ct^2|w(t)|_k|\phi (t)|_{\mathrm{}}Ct^2h_0(t)ab$$
under the conditions stated on $`k^{}`$, and therefore by integration
$$|\stackrel{~}{w}(t)\stackrel{~}{w}(t_1)|_k^{}Cabh(tt_1)$$
which implies the existence of $`w_+K_\rho ^k^{}`$ satisfying the estimate (5.5). By standard compactness arguments, $`w_+K_\rho ^k`$, $`|w_+|_ka`$ and $`w(t)`$ tends weakly to $`w_+`$ in $`K_\rho ^k`$ when $`t\mathrm{}`$. The other convergences follow by interpolation.
$``$$``$
We now estimate the asymptotic functions $`(w_m,\phi _m)`$ defined by successive integrations from the system (2.18)-(2.19) with initial condition (2.20)-(2.21). For that purpose, we need the function
$$\overline{h}_0(t)=_1^t𝑑t_1t_1^\gamma $$
$`(5.7)`$
and the associated estimating functions defined by (3.68) (3.69) (3.72), which we denote by $`\overline{N}_m`$, $`\overline{Q}_m`$ and $`\overline{P}_m`$ for that special choice of $`h_0`$. (Those functions appeared already in II where they were called $`h_0`$, $`N_m`$, $`Q_m`$ and $`P_m`$, see (II.3.19) (II.3.25) (II.3.26) (II.3.31)). We recall that $`\lambda =\mu n+2`$ and we define $`\overline{\lambda }=\lambda 1`$. The following proposition is the extension to the present context of Proposition II.5.1.
Proposition 5.2. Let $`p0`$ be an integer, let $`k`$ satisfy
$$k>n/2,k(p+2)\overline{\lambda }1,$$
and let
$$k_m=km\overline{\lambda },\mathrm{}_m=km\overline{\lambda }\lambda ,0mp+1.$$
$`(5.8)`$
Let $`w_+K_\rho ^k`$ and let $`a=|w_+|_k`$. Let $`\{w_0=w_+,w_{m+1}\}`$ and $`\{\phi _m\}`$, $`0mp`$ be the solution of the system (2.18)-(2.19) with initial conditions (2.20)-(2.21). Then
(1) $`w_{m+1}𝒞([1,\mathrm{})`$, $`K_\rho ^{k_{m+1}})`$, $`\phi _m𝒞([1,\mathrm{}),Y_\rho ^\mathrm{}_m)`$ and the following estimates hold for all $`t1`$ :
$$|w_{m+1}(t)|_{k_{m+1}}A(a)\overline{Q}_m(t),$$
$`(5.9)`$
$$|\phi _m(t)|_\mathrm{}_mA(a)\overline{N}_m(t),$$
$`(5.10)`$
for some estimating function $`A(a)`$.
If in addition $`(p+2)\gamma >1`$ and if we define $`\phi _{p+1}`$ by (2.19) with initial condition $`\phi _{p+1}(\mathrm{})=0`$, then $`\phi _{p+1}𝒞([1,\mathrm{}),Y^{\mathrm{}_{p+1}})`$ and the following estimate holds :
$$|\phi _{p+1}(t)|_{\mathrm{}_{p+1}}A(a)\overline{P}_p(t).$$
$`(5.11)`$
(2) The functions $`\{\phi _m\}`$ are gauge invariant, namely if $`w_+^{}=w_+\mathrm{exp}(i\omega )`$ for some real valued function $`\omega `$ and if $`w_+^{}`$ gives rise to $`\{\phi _m^{}\}`$, then $`\phi _m^{}=\phi _m`$ for $`0mp+1`$.
(3) The map $`w_+\{w_{m+1},\phi _m\}`$ is uniformly Lipschitz continuous on the bounded sets from the norm topology of $`w_+`$ in $`K_\rho ^k`$ to the norms $`\overline{Q}_m^1w_{m+1};L^{\mathrm{}}([1,\mathrm{}),K_\rho ^{k_{m+1}})`$ and $`\overline{N}_m^1\phi _m;L^{\mathrm{}}([1,\mathrm{}),Y_\rho ^\mathrm{}_m)`$, $`0mp`$. A similar continuity holds for $`\phi _{p+1}`$.
Proof. The proof is essentially the same as that of Proposition II.5.1.
Part (1). We proceed by induction on $`m`$, starting from the assumption on $`w_+`$. We assume the results to hold for $`(w_j,\phi _j)`$ for $`jm`$ and we prove them for $`w_{m+1}`$ and $`\phi _{m+1}`$. We first consider $`w_{m+1}`$. Letting the exponents $`(k_m,\mathrm{}_m)`$ be undefined for the moment except for being nonincreasing in $`m`$, we obtain from (2.18) and from (3.23) (3.24)
$$|_tw_{m+1}|_{k_{m+1}}A(a)t^2\left\{\underset{0jm1}{}\overline{N}_j(t)\overline{Q}_{mj1}(t)+\overline{N}_m(t)\right\}$$
$`(5.12)`$
under the conditions
$$\{\begin{array}{c}k_{m+1}0,k_{m+1}(k_m1)\mathrm{}_m,\hfill \\ \\ k_{m+1}+n/2<k_j+\mathrm{}_{mj},0jm.\hfill \end{array}$$
$`(5.13)`$
Integrating (5.12) between $`t`$ and $`\mathrm{}`$ with $`w_{m+1}(\mathrm{})=0`$ and using (3.79) (3.75) yields the result for $`w_{m+1}`$.
We next consider $`\phi _{m+1}`$. From (2.19) and from (3.26) (3.27) (3.28) (3.29) we obtain
$$|_t\phi _{m+1}|_{\mathrm{}_{m+1}}A(a)\left\{t^2\underset{0jm}{}\overline{N}_j(t)\overline{N}_{mj}(t)+t^\gamma \left(\underset{0jm1}{}\overline{Q}_j(t)\overline{Q}_{m1j}(t)+\overline{Q}_m(t)\right)\right\}$$
$`(5.14)`$
under the conditions
$$\{\begin{array}{c}\mathrm{}_{m+1}+10,\mathrm{}_{m+1}(\mathrm{}_m1)(k_{m+1}\lambda ),\hfill \\ \\ \mathrm{}_{m+1}+\frac{n}{2}<\mathrm{}_j+\mathrm{}_{mj},\mathrm{}_{m+1}+\lambda +n/2<k_j+k_{m+1j},0jm,\hfill \end{array}$$
$`(5.15)`$
and for $`\phi _0`$
$$\mathrm{}_0+\lambda k,\mathrm{}_0+\lambda +n/2<2k.$$
$`(5.16)`$
Integrating (5.14) between 1 and $`t`$ with $`\phi _m(1)=0`$ and using (3.78) (3.76) and (3.80) (3.84) yields the result for $`\phi _{m+1}`$ (and similarly for $`\phi _0`$).
We saturate the nonstrict part of (5.13) (5.15) (5.16) by the choice (5.8), where in addition we optimize (maximize) $`\{\mathrm{}_m\}`$ for given $`k`$. The strict conditions then reduce to $`k>n/2`$, while the condition $`k(p+2)\overline{\lambda }1`$ is simply the condition $`k_{p+1}(\mathrm{}_{p+1}+1)0`$.
Finally, if $`(p+2)\gamma >1`$, we integrate (5.14) with $`m=p`$ between $`t`$ and $`\mathrm{}`$ and use (3.77) (3.78) and (3.80) (3.82) with $`m=p`$.
Part (2). The proof is identical with that of Proposition II.5.1 and will be omitted.
Part (3). From the fact that the RHS of (2.18)-(2.19) are bilinear, it follows by induction of $`m`$ that the difference between two solutions $`\{w_m,\phi _m\}`$ and $`\{w_m^{},\phi _m^{}\}`$ associated with $`w_+`$ and $`w_+^{}`$ is estimated by
$$|w_{m+1}w_{m+1}^{}|_{k_{m+1}}A(a)|w_+w_+^{}|_k\overline{Q}_m(t),$$
$`(5.17)`$
$$|\phi _m\phi _m^{}|_\mathrm{}_mA(a)|w_+w_+^{}|_k\overline{N}_m(t)$$
$`(5.18)`$
for $`0mp`$, and if $`(p+2)\gamma >1`$,
$$|\phi _{p+1}\phi _{p+1}^{}|_{\mathrm{}_{p+1}}A(a)|w_+w_+^{}|_k\overline{P}_p(t),$$
$`(5.19)`$
where $`a=|w_+|_k|w_+^{}|_k`$. The continuity stated in Part (3) follows from those estimates.
$``$$``$
We now turn to the main result of this section, namely to the proof of existence of asymptotic states $`(w_+,\psi _+)`$ for solutions of the auxiliary system (2.11)-(2.12). That result relies heavily on suitable estimates of the remainders
$$q_{m+1}(t)=w(t)W_m(t)$$
$`(5.20)`$
$$\psi _{m+1}(t)=\phi (t)\varphi _m(t)$$
$`(5.21)`$
where $`W_m`$ and $`\varphi _m`$ are defined (see (2.14) (2.15)) by
$$W_m=\underset{0jm}{}w_j,\varphi _m=\underset{0jm}{}\phi _j.$$
$`(5.22)`$
In view of Proposition 4.1, we shall consider solutions $`(w,\phi )`$ of the system (2.11)-(2.12) such that $`(w,h_0^1\phi )𝒴_\rho ^{k,\mathrm{}}([T,\mathrm{}))`$, where $`h_0`$ is a suitable $`𝒞^1`$ positive increasing function of time. We shall assume in addition that $`t^\gamma ch_0^{}`$, a property which occurs naturally in the interesting examples relevant for Section 4. In addition to the estimating functions $`\overline{N}_m`$, $`\overline{Q}_m`$ and $`\overline{P}_m`$ associated with $`\overline{h}_0`$, we shall also need the estimating functions $`N_m`$, $`Q_m`$ and $`P_m`$ associated with $`h_0`$, defined by (3.68) (3.69) (3.72). From the relation $`t^\gamma ch_0^{}`$, it follows that $`\overline{N}_m`$, $`\overline{Q}_m`$ and $`\overline{P}_m`$ are estimated by $`N_m`$, $`Q_m`$, and $`P_m`$, and more precisely
$$\overline{N}_mc^{m+1}N_m,\overline{Q}_mc^{m+1}Q_m,\overline{P}_mc^{m+2}P_m.$$
$`(5.23)`$
We can now state the main result of this section, which is the extension of Proposition II.6.1 to the present context.
Proposition 5.3. Let $`p0`$ be an integer. Let $`(k,\mathrm{},k_0)`$ satisfy
$$k>n/2+(p2)\overline{\lambda }+\lambda 0,kk_0\overline{\lambda }+2,\mathrm{}k_0\lambda ,$$
$`(5.24)`$
$$k_0>n/2,k_0(p+2)\overline{\lambda }1,$$
$`(5.25)`$
and let
$$k_m=k_0m\overline{\lambda },\mathrm{}_m=k_0\lambda m\overline{\lambda },0mp+1.$$
$`(5.26)`$
Let $`h_0`$ be a $`𝒞^1`$ positive increasing function defined in $`[1,\mathrm{})`$, such that $`t^2h_0`$, $`t^1h_0^{}L^1([1,\mathrm{}))`$ and $`t^\gamma ch_0^{}`$. Let $`T2`$, let $`(w,\phi )`$ be a solution of the system (2.11)-(2.12) such that $`(w,h_0^1\phi )𝒴_\rho ^{k,\mathrm{}}([T,\mathrm{}))`$ and define $`a`$, $`b`$ by (5.6). Let $`w_+=\underset{t\mathrm{}}{lim}w(t)K_\rho ^k`$ be defined by Proposition 5.1, so that in particular
$$|w(t)w_+|_{k_1}0\text{when}t\mathrm{}.$$
$`(5.27)`$
Let $`(w_{m+1},\phi _m)`$, $`0mp`$ be defined by Proposition 5.2 and let $`(W_m,\varphi _m)`$, $`0mp`$, be defined by (5.22). Then the following estimates hold for all $`t[T,\mathrm{})`$ :
$$|w(t)W_m(t)|_{k_{m+1}}A(a,b)Q_m(t),$$
$`(5.28)`$
$$|\phi (t)\varphi _m(t)|_{\mathrm{}_{m+1}}A(a,b)N_{m+1}(t).$$
$`(5.29)`$
for $`0mp`$, and for some estimating function $`A(a,b)`$.
If in addition $`(p+2)\gamma >1`$ and if $`P_p(1)<\mathrm{}`$, then the following limit exists
$$\underset{t\mathrm{}}{lim}\left(\phi (t)\varphi _p(t)\right)=\psi _+$$
$`(5.30)`$
as a strong limit in $`Y_\rho ^{\mathrm{}_{p+1}}`$, and the following estimate holds
$$|\phi (t)\varphi _p(t)\psi _+|_{\mathrm{}_{p+1}}A(a,b)P_p(t).$$
$`(5.31)`$
Remark 5.1. When applied to solutions $`(w,\phi )`$ of the system (2.11)-(2.12) obtained in Proposition 4.1, the time decay estimates of Proposition 5.3 will take the following typical form. Take $`|\rho ^{}|=t^{1\epsilon }`$ and $`h_0=t^{1\gamma +\epsilon }`$, which is adequate for Propostion 4.1. Then for $`(p+1)\gamma <1`$ and sufficiently small $`\epsilon `$
$$N_m(t)t^{1(m+1)(\gamma \epsilon )},Q_mt^{(m+1)(\gamma \epsilon )},P_p(t)t^{1(p+2)(\gamma \epsilon )}$$
and the condition $`P_p(1)<\mathrm{}`$ reduces to $`(p+2)(\gamma \epsilon )>1`$. For $`(p+1)(\gamma \epsilon )>1`$, the time decay saturates at $`N_p(t)1`$, $`Q_p(t)t^1`$ and $`P_p(t)t^\gamma `$, as explained in Section II.3.
Remark 5.2. We have kept the parameter $`k_0`$ in the statement of Proposition 5.3 because it plays a central role in the proof. For a given solution $`(w,\phi )`$, namely for given $`(k,\mathrm{})`$, we can optimize the results by maximizing $`k_0`$ as allowed by (5.24), namely by taking $`k_0=(k+\overline{\lambda }2)(\mathrm{}+\lambda )`$. The condition (5.25) then reduces to lower bounds on $`(k,\mathrm{})`$. Those bounds are stronger than (5.4) and therefore allow for the application of Proposition 5.1. Note also that the regularity obtained for the remainders $`(q_m,\psi _m)`$ for $`m1`$ is weaker than that of the estimating functions of the same level $`(w_m,\phi _m)`$ since $`(k_m,\mathrm{}_m)`$ in (5.26) contain only $`k_0k`$ whereas $`(k_m,\mathrm{}_m)`$ in (5.8) contain $`k`$ and $`w_+K_\rho ^k`$ in both cases.
Proof of Proposition 5.3. The proof is essentially the same as that of Proposition II.6.1 and proceeds by an induction on $`m`$, the starting point of which is the estimate for $`q_1`$. We assume the estimates (5.28) (5.29) to hold for $`(q_j,\psi _j)`$, $`0jm`$, and we derive them for $`(q_{m+1},\psi _{m+1})`$, with $`(q_m,\psi _m)`$ defined by (5.20) (5.21) and $`(q_0,\psi _0)=(w,\phi )`$.
We substitute the decompositions $`w=W_m+q_{m+1}`$ and $`\phi =\varphi _m+\psi _{m+1}`$ in the LHS of (2.11)-(2.12) and the decompositions $`w=W_{m1}+q_m`$ and $`\phi =\varphi _{m1}+\psi _m`$ in the RHS of the same, thereby obtaining
$$_tq_{m+1}=\left(2t^2\right)^1\{i\mathrm{\Delta }w+(2\phi +(\mathrm{\Delta }\phi ))q_m$$
$$+(2\psi _m+(\mathrm{\Delta }\psi _m))W_{m1}+\underset{\genfrac{}{}{0pt}{}{0i,jm1}{i+jm}}{}(2\phi _i+(\mathrm{\Delta }\phi _i))w_j\}$$
$`(5.32)`$
$$_t\psi _{m+1}=\left(2t^2\right)^1\left\{\left(\phi +\varphi _{m1}\right)\psi _m+\underset{\genfrac{}{}{0pt}{}{0i,jm1}{i+jm}}{}\phi _i\phi _j\right\}$$
$$+t^\gamma \left\{g_0(q_m,q_1)+g_0(q_m,W_{m1}w_0)+2g_0(q_{m+1},w_0)+\underset{\genfrac{}{}{0pt}{}{0i,jm1}{i+jm+1}}{}g_0(w_i,w_j)\right\}$$
$`(5.33)`$
for $`m1`$ and
$$\{\begin{array}{c}_tq_1=\left(2t^2\right)^1\left\{i\mathrm{\Delta }w+2\phi w+(\mathrm{\Delta }\phi )w\right\}\hfill \\ \\ _t\psi _1=\left(2t^2\right)^1|\phi |^2+t^\gamma g_0(q_1,w+w_+)\hfill \end{array}$$
$`(5.34)`$
for $`m=0`$. We let the exponents $`(k_m,\mathrm{}_m)`$ be undefined for the moment, except for the property of being decreasing in $`m`$ and of being not larger than the corresponding exponents of Proposition 5.2, which we denote momentarily by $`(\overline{k}_m,\overline{\mathrm{}}_m)`$ inside this proof, namely
$$k_m\overline{k}_m=km\overline{\lambda },\mathrm{}_m\overline{\mathrm{}}_m=k\lambda m\overline{\lambda }.$$
We estimate (5.32) by (3.23) (3.24), by (5.9) (5.10) (5.23) and by the induction assumption, and we estimate similarly (5.33) by (3.26) (3.27) (3.28) (3.29) and the same other ingredients, thereby obtaining
$$|_tq_{m+1}|_{k_{m+1}}A(a,b)t^2\left\{1+h_0Q_{m1}+N_m+\underset{mi+j2(m1)}{}N_iQ_{j1}\right\}$$
$`(5.35)`$
$$|_t\psi _{m+1}|_{\mathrm{}_{m+1}}A(a,b)\{t^2(h_0N_m+\underset{mi+j2(m1)}{}N_iN_j)$$
$$+t^\gamma (Q_0Q_{m1}+Q_m+\underset{m+1i+j2(m1)}{}Q_{i1}Q_{j1})\}$$
$`(5.36)`$
for $`m1`$, under the conditions
$$k_{m+1}(k_m1)\mathrm{}_m,\mathrm{}_{m+1}(\mathrm{}_m1)(k_{m+1}\lambda ),$$
$`(5.37)`$
$$\{\begin{array}{c}k_{m+1}+n/2<\left(\mathrm{}+k_m\right)\left(\mathrm{}_m+\overline{k}_{m1}\right)\hfill \\ \\ \mathrm{}_{m+1}+n/2<\left(\mathrm{}+\mathrm{}_m\right)\left(\mathrm{}_m+\overline{\mathrm{}}_{m1}\right)\hfill \\ \\ \mathrm{}_{m+1}+n/2+\lambda <\left(k_m+k_1\right)\left(k_m+\overline{k}_{m1}\right)\left(k_{m+1}+k\right),\hfill \end{array}$$
$`(5.38)`$
and
$$|_tq_1|_{k_1}Ct^2(a+h_0ab)$$
$`(5.39)`$
$$|_t\psi _1|_\mathrm{}_1Ct^2h_0^2b^2+Ct^\gamma A(a,b)Q_0$$
$`(5.40)`$
under the conditions
$$k_1(k2)\mathrm{},\mathrm{}_1(\mathrm{}1)(k_1\lambda ),$$
$`(5.41)`$
$$k_1+n/2<k+\mathrm{},\mathrm{}_1+n/2<2\mathrm{}(k_1+k\lambda )$$
$`(5.42)`$
for $`m=0`$.
As in the proof of Proposition 5.2, we saturate (5.37) and maximize $`\mathrm{}_m`$ with respect to $`k_m`$ by the choice (5.26). Then (5.41) reduces to the last two conditions of (5.24), while (5.38) is easily seen to reduce to $`k_0>n/2`$ and to the first condition of (5.24), and (5.42) follows from (5.24) and (5.25). The conditions $`k_m\overline{k}_m`$ and $`\mathrm{}_m\overline{\mathrm{}}_m`$ follow from the fact that $`kk_0`$, implied by (5.24).
Integrating (5.39) between $`t`$ and $`\mathrm{}`$ with initial condition (5.27) yields
$$|q_1(t)|_{k_1}C\left(a+h_0(1)\right)t^1+CabQ_0(t),$$
$`(5.43)`$
namely the estimate (5.28) for $`m=0`$ (since $`Q_0(t)t^1Q_0(1))`$ which is the starting point of the induction procedure. Integrating (5.35) between $`t`$ and $`\mathrm{}`$ with the initial condition $`q_{m+1}(\mathrm{})=0`$ for $`m1`$ and using (3.79) (3.75) yields (5.28). Similarly, integrating (5.40) and (5.36) between $`T`$ and $`t`$ and using (3.78) (3.76) and (3.80) (3.84) yields (5.29). Finally, if $`(p+2)\gamma >1`$ and $`P_p(1)<\mathrm{}`$, the RHS of (5.36) with $`m=p`$ is integrable at infinity in time, which proves the existence of the limit (5.30). Integrating (5.36) between $`t`$ and $`\mathrm{}`$ and using (3.78) (3.77) and (3.80) (3.82) yields (5.31).
$``$$``$
## 6 Cauchy problem at infinity and wave operators for the auxiliary system
In this section we derive the main technical result of this paper, which is in some sense the converse of Proposition 5.3, namely we prove that sufficiently regular asymptotic states $`(w_+,\psi _+)`$ generate solutions $`(w,\phi )`$ of the system (2.11)-(2.12) in the sense described in Section 2, thereby allowing for the definition of the local wave operator at infinity $`\mathrm{\Omega }_0`$ : $`(w_+,\psi _+)(w,\phi )`$. As a preliminary, and in order to allow for an easy proof of the gauge invariance of the construction, we first solve the linear transport equations (2.23) (2.24) with initial condition (2.25) and derive some asymptotic properties of their solutions.
In all this section, as in Proposition 4.1 and in contrast with Section 5, we are again solving the system (2.11)-(2.12) and we have to solve the simpler equations (2.23) (2.24) in the same framework. As a consequence, we again need the full machinery of Gevrey spaces with time dependent $`\rho `$ so as to be able to use (3.12), we need the integral norms in $`𝒳_\rho ^{k,\mathrm{}}`$ and the integration by parts (3.22) and (3.25). We therefore begin by choosing $`|\rho ^{}|`$ exactly as in Section 4. Now however in contrast with Proposition 4.1 where we kept $`t_0`$ finite, we want to take $`t_0=\mathrm{}`$, and therefore we must take $`\rho `$ to be increasing. Therefore, in all this section, we take
$$\rho (t)=\rho _{\mathrm{}}_t^{\mathrm{}}𝑑t_1|\rho ^{}(t_1)|,$$
$`(6.1)`$
taking $`\rho _{\mathrm{}}`$ sufficiently large for $`\rho (t)`$ to be nonnegative in the (asymptotic) region of interest, a sufficient condition for which being that $`\rho _{\mathrm{}}|\rho ^{}|;L^1([1,\mathrm{}))`$. Except for that condition, $`\rho _{\mathrm{}}`$ will be arbitrary (but fixed), and all subsequent estimates will be independent of $`\rho _{\mathrm{}}`$, for the same reasons as in the previous sections.
Solving the Cauchy problem either for the system (2.11)-(2.12) or for the transport equations (2.23) (2.24) with infinite initial time will be done as in II by first solving that system (or equation) with large but finite initial time $`t_0`$ and then letting $`t_0`$ tend to infinity. In order not to make this paper too cumbersome, we shall restrict our attention to solving that problem only for $`tt_0`$ when $`t_0`$ is finite. The solutions could easily be extended to $`[t_0,\mathrm{})`$ with a modified $`\rho `$ of the type (4.1), but that extension would be useless in the limit $`t_0\mathrm{}`$, and we shall refrain from performing it.
By analogy with the spaces $`𝒳_\rho ^{k,\mathrm{}}(I)`$ defined by (4.2) (4.3), we extend the definition (5.1) of the spaces $`𝒴_\rho ^{k,\mathrm{}}(I)`$ from the case of constant $`\rho `$ to the case of variable $`\rho `$ considered in this section by
$$𝒴_\rho ^{k,\mathrm{}}(I)=\{(w,\phi ):(F^1f\widehat{w},F^1f\widehat{\phi })𝒴_0^{k,\mathrm{}}(I)\}$$
$`(6.2)`$
with $`f`$ defined by (3.1), $`\rho `$ defined by (6.1) and $`𝒴_0^{k,\mathrm{}}(I)`$ defined by (5.1). Occasionally, we shall have to state that a single function $`w`$ or $`\phi `$ belongs to the $`w`$-subspace or to the $`\phi `$-subspace of some space $`𝒳_\rho ^{k,\mathrm{}}`$ or $`𝒴_\rho ^{k,\mathrm{}}`$. In order to avoid introducing additional notation, we shall then write $`(w,0)𝒳_\rho ^{k,0}(I)`$ or $`(0,\phi )𝒳_\rho ^{0,\mathrm{}}(I)`$, and similarly with $`𝒳`$ replaced by $`𝒴`$.
We shall have to consider norms of the type $`|w_{}|_k`$ or $`|\phi _{}|_{\mathrm{}}`$ for $`w_{}`$ or $`\phi _{}`$ that are differences of functions taken at different times, possibly leaving in doubt the value of $`t`$ appearing in $`\rho (t)`$ in the definition of the norm. In such cases it will be understood that the value of $`t`$ appearing in $`\rho (t)`$ should be the smaller of the times appearing in $`w_{}`$ or $`\phi _{}`$, thereby yielding the smaller of the corresponding values of $`\rho `$.
We begin with the study of the transport equation (2.23). As compared with the treatment of that problem given in II, however, a new difficulty arises. In order to compare $`V`$ with the asymptotic approximation $`W_p`$ to the anticipated solution of (2.11)-(2.12), it is no longer sufficient to take for $`V`$ the initial condition $`V(t_0)=w_+`$ when solving (2.23) with finite initial time $`t_0`$, and we have to use instead the better initial condition $`V(t_0)=W_p(t_0)`$. On the other hand, the results on the Cauchy problem for (2.23) (2.24) do not depend on detailed properties of $`\varphi _{p1}`$ and $`W_p`$. They are therefore stated in Propositions 6.1 and 6.2 in terms of general functions $`\varphi `$ and $`W`$, to be taken as $`\varphi _{p1}`$ and $`W_p`$ from Proposition 6.3 on.
In all this section, we use systematically the notation $`y_{(1)}`$, $`Y`$, $`z_{(1)}`$, $`Z`$ defined by (4.4)-(4.9).
We begin with the study of the transport equation (2.23) which we rewrite with general $`\varphi `$ as
$$_tV=\left(2t^2\right)^1\left(2\varphi +(\mathrm{\Delta }\varphi )\right)V.$$
$`(6.3)`$
Proposition 6.1. Let $`(\stackrel{~}{k},\stackrel{~}{\mathrm{}},\overline{k},k)`$ satisfy
$$\stackrel{~}{\mathrm{}}>n/2\nu ,\stackrel{~}{k}(\stackrel{~}{\mathrm{}}+\nu /2)\overline{k}k+1\nu ,k\nu /2.$$
$`(6.4)`$
Let $`1T<\mathrm{}`$. Let $`h_0`$ and $`h_1`$ be $`𝒞^1`$ positive functions defined in $`[T,\mathrm{})`$ with $`h_0`$ nondecreasing, $`h_1`$ nonincreasing and tending to zero at infinity, and $`h_1t^2\rho ^1h_0`$. Let $`w_+K_\rho _{_{\mathrm{}}}^{\stackrel{~}{k}}`$. Let $`(W,\varphi )`$ be such that $`(W,h_0^1\varphi )𝒴_\rho ^{\stackrel{~}{k},\stackrel{~}{\mathrm{}}}([T,\mathrm{}))`$ and that $`W(t)`$ tends to $`w_+`$ as $`t\mathrm{}`$, with an estimate
$$|W(t)w_+|_kc_1h_1(t)$$
$`(6.5)`$
for some constant $`c_1`$. Let
$$a=\underset{t}{Sup}|W(t)|_{\stackrel{~}{k}},b=\underset{t}{Sup}h_0^1(t)|\varphi (t)|_\stackrel{~}{\mathrm{}}.$$
$`(6.6)`$
Then
(1) There exist constants $`c`$ and $`C`$ such that if
$$bh_1(T)c,$$
$`(6.7)`$
there exists a unique solution $`V`$ of the equation (6.3) such that $`(V,0)𝒳_\rho ^{\overline{k},0}([T,\mathrm{}))`$ and such that the following estimates hold :
$$Y(V;[T,\mathrm{}),1,\overline{k})Ca,$$
$`(6.8)`$
$$Y(Vw_+;[T,\mathrm{}),h_1,k)Cab.$$
$`(6.9)`$
(2) $`V`$ is the limit as $`t_0\mathrm{}`$ of solutions $`V_{t_0}`$ of (6.3) such that $`V_{t_0}(t_0)=W(t_0)`$ and $`(V_{t_0},0)𝒳_\rho ^{\overline{k},0}([T,t_0])`$. The convergence is in the strong sense in $`𝒳_\rho ^{k^{},0}([T,T_1])`$ for $`k^{}<\overline{k}`$ and in the weak-$``$ sense in $`𝒳_\rho ^{\overline{k},0}([T,T_1])`$ for every $`T_1`$, $`T<T_1<\mathrm{}`$ and the following estimate holds for all $`t_0>T`$
$$Y(VV_{t_0};[T,t_0],1,k)C\left(ab+c_1\right)h_1(t_0).$$
$`(6.10)`$
(3) The solution $`V`$ is unique in $`L^{\mathrm{}}([T,\mathrm{}),L^2)`$ under the condition that $`V(t)w_+_2`$ tends to zero when $`t\mathrm{}`$.
Remark 6.1. From the uniqueness statement of Proposition 6.1, Part (2), it follows that for given $`\varphi `$ and $`w_+`$, $`V`$ is independent of $`W`$. Actually Parts (1) and (3) make no reference to $`W`$ and could be proved by taking $`W(t)w_+`$. $`W`$ appears only in the limiting process of Part (2), which however will be esssential to derive the more accurate estimates of Proposition 6.3.
Proof. We prove Parts (1) and (2) together.
We first solve (6.3) with initial data $`V_{t_0}(t_0)=W(t_0)`$ at finite $`t_0`$. This is a linear transport equation with $`𝒞^{\mathrm{}}`$ vector field $`\varphi `$ and $`𝒞^{\mathrm{}}`$ initial data and the existence and uniqueness of a solution, for instance with value in $`H^N`$, is a standard result. We concentrate on the Gevrey estimates and on the subsequent limit $`t_0\mathrm{}`$.
In the same way as in the proof of Lemma 3.5, from (3.22) (3.24) we obtain for $`tt_0`$
$$_t|V_{t_0}|_{\overline{k}}^22\rho ^{}|V_{t_0}|_{\overline{k}+\nu /2}^2Ct^2|V_{t_0}|_{\overline{k}+\nu /2}^2|\varphi |_\stackrel{~}{\mathrm{}}$$
$`(6.11)`$
under the conditions
$$\stackrel{~}{\mathrm{}}>n/2\nu ,\overline{k}\nu /2,\stackrel{~}{\mathrm{}}+1\overline{k}\nu /2$$
and
$$\overline{k}+\stackrel{~}{\mathrm{}}+\nu /2>\overline{k}\nu /2+n/2,\stackrel{~}{\mathrm{}}\overline{k}\nu /2$$
which follow from (6.4). Integrating (6.11) over time and using (6.6), we obtain as in the proof of Lemma 4.1
$$y^2y_1^2y_0^2+Cby_1^2h_1(t)$$
$`(6.12)`$
where $`y_{(1)}=y_{(1)}(V_{t_0};[t,t_0],1,\overline{k})`$ and $`y_0=|W(t_0)|_{\overline{k}}`$, which under the condition (6.7) yields
$$Y(V_{t_0};[T,t_0],1,\overline{k})Ca.$$
$`(6.13)`$
We next estimate the difference $`v_1(t)V_{t_0}(t)V_{t_0}(t_0)V_{t_0}(t)W(t_0)`$, which satisfies the equation
$$_tv_1=\left(2t^2\right)^1\left(2\varphi +(\mathrm{\Delta }\varphi )\right)V_{t_0}$$
$`(6.14)`$
with initial condition $`v_1(t_0)=0`$. Let $`\stackrel{~}{v}_1=h_1^1v_1`$. From (3.23) (3.24) with $`m=k\nu /2`$, $`k\overline{k}+\nu /2`$ and $`\mathrm{}=\stackrel{~}{\mathrm{}}`$, we obtain
$$_t|\stackrel{~}{v}_1|_k^22\rho ^{}|\stackrel{~}{v}_1|_{k+\nu /2}^2Ct^2h_1^1|\stackrel{~}{v}_1|_{k+\nu /2}|V_{t_0}|_{\overline{k}+\nu /2}|\varphi |_\stackrel{~}{\mathrm{}}$$
$`(6.15)`$
under the conditions
$$\overline{k}+\stackrel{~}{\mathrm{}}>k+n/2\nu ,\stackrel{~}{\mathrm{}}k\nu /20,\overline{k}+\nu /2k+1\nu /2,$$
which follow from (6.4).
Defining $`y_{(1)}=y_{(1)}(v_1;[T,t_0],h_1,k)`$, integrating over time and using (6.6) (6.13) and the Schwarz inequality we obtain in the same way as before
$$y^2y_1^2Caby_1\left\{\underset{t}{Sup}t^2\rho ^1h_0h_1^1\right\}=Caby_1$$
and therefore
$$Y(V_{t_0}W(t_0);[T,t_0],h_1,k)Cab.$$
$`(6.16)`$
We now take the limit $`t_0\mathrm{}`$, and for that purpose we estimate the difference $`V_{t_1}V_{t_0}`$ of two solutions corresponding to $`t_0`$ and $`t_1`$, with $`T<t_0t_1`$. Since the equation (6.3) is linear in $`V`$, the difference of two solutions is estimated in the same way as a single solution. In the same way as in the proof of (6.13), we obtain
$$Y(V_{t_1}V_{t_0};[T,t_0],1,k)C\left|V_{t_1}(t_0)V_{t_0}(t_0)\right|_k$$
$$C\left\{\left|V_{t_1}(t_0)W(t_1)\right|_k+\left|W(t_1)W(t_0)\right|_k\right\}.$$
$`(6.17)`$
We estimate the first norm in the last member of (6.17) by (6.16) with $`t_1`$ replacing $`t_0`$, and where we use the pointwise estimate taken at $`t=t_0`$, and we estimate the second norm by (6.5) used both for $`t=t_0`$ and $`t=t_1`$, thereby obtaining
$$Y(V_{t_1}V_{t_0};[T,t_0],1,k)C\left(ab+c_1\right)h_1(t_0).$$
$`(6.18)`$
From (6.18) it follows that when $`t_0\mathrm{}`$, $`V_{t_0}`$ converges to some $`V`$ such that $`(V,0)𝒳_\rho ^{k,0}([T,\mathrm{}))`$ strongly in $`𝒳_\rho ^{k,0}([T,T_1])`$ for all $`T_1`$, $`T<T_1<\mathrm{}`$. By standard compactness arguments and by (6.13) $`(V,0)𝒳_\rho ^{\overline{k},0}([T,\mathrm{}))`$, $`V`$ satisfies (6.8) and the convergence holds in the sense of Part (2) of the proposition. Furthermore, taking the limit $`t_1\mathrm{}`$ in (6.18) yields (6.10).
It remains only to prove (6.9). For that purpose, we take again $`T<t_0<t_1`$ and we estimate
$$Y(Vw_+;[T,t_0],h_1,k)Y(V_{t_1}W(t_1);[T,t_0],h_1,k)$$
$$+h_1(t_0)^1\left\{Y(VV_{t_1};[T,t_0],1,k)+Y(W(t_1)w_+;[T,t_0],1,k)\right\}$$
$$Cab+h_1(t_0)^1\left\{C\left(ab+c_1\right)h_1(t_1)+Cc_1h_1(t_1)\left(\rho (t_1)\rho (t_0)\right)^{1/2}\right\}$$
by (6.16) and (6.10) with $`t_1`$ replacing $`t_0`$, and by (6.5) and the fact that $`\rho (t)`$ is strictly increasing, so that
$$_T^{t_0}𝑑t\rho ^{}|W(t_1)w_+|_{\rho (t),k+\nu /2}^2C|W(t_1)w_+|_{\rho (t_1),k}^2|\rho (t_1)\rho (t_0)|^1\rho ^{}_1$$
where $`||_{\rho ,k}`$ denotes the norm in $`K_\rho ^k`$. Taking the limits $`t_1\mathrm{}`$ and $`t_0\mathrm{}`$ in that order yields (6.9), with the same constant as in (6.16).
Part (3) follows from an elementary energy estimate for the $`L^2`$ norm of the difference of two solutions, namely
$$V_1(t)V_2(t)_2V_1(t^{})V_2(t^{})_2\mathrm{exp}\left(Cb|h(t)h(t^{})|\right)$$
where $`h(t)=_t^{\mathrm{}}𝑑t_1t_1^2h_0(t_1)`$.
$``$$``$
We now turn to the transport equation (2.24), which we rewrite with general $`\varphi `$ as
$$_t\chi =t^2\varphi \chi .$$
$`(6.19)`$
Proposition 6.2. Let $`(\stackrel{~}{\mathrm{}},\overline{\mathrm{}},\mathrm{})`$ satisfy
$$\stackrel{~}{\mathrm{}}>n/2\nu ,\stackrel{~}{\mathrm{}}+\nu /2\overline{\mathrm{}}+1,\overline{\mathrm{}}\mathrm{}+1\nu ,\mathrm{}+1\nu /2.$$
$`(6.20)`$
Let $`1T<\mathrm{}`$ and let $`h_0`$, $`h_1`$ be as in Proposition 6.1. Let $`\psi _+Y_\rho _{_{\mathrm{}}}^\overline{\mathrm{}}`$ and let $`\beta =|\psi _+|_\overline{\mathrm{}}`$. Let $`\varphi `$ be such that $`(0,h_0^1\varphi )Y_\rho ^{0,\stackrel{~}{\mathrm{}}}([T,\mathrm{}))`$ with
$$b=\underset{t}{Sup}h_0(t)^1|\varphi (t)|_\stackrel{~}{\mathrm{}}.$$
$`(6.21)`$
Then
(1) There exist constants $`c`$ and $`C`$ such that if (6.7) holds, there exists a unique solution $`\chi `$ of the equation (6.19) such that $`(0,\chi )𝒳_\rho ^{0,\overline{\mathrm{}}}([T,\mathrm{}))`$ and such that the following estimates hold
$$Z(\chi ;[T,\mathrm{}),1,\overline{\mathrm{}})C\beta ,$$
$`(6.22)`$
$$Z(\chi \psi _+;[T,\mathrm{}),h_1,\mathrm{})Cb\beta .$$
$`(6.23)`$
(2) $`\chi `$ is the limit as $`t_0\mathrm{}`$ of solutions $`\chi _{t_0}`$ of (6.19) such that $`\chi _{t_0}(t_0)=\psi _+`$ and $`(0,\chi _{t_0})𝒳_\rho ^{0,\overline{\mathrm{}}}([T,t_0))`$. The convergence is in the strong sense in $`𝒳_\rho ^{0,\mathrm{}^{}}([T,T_1])`$ for $`\mathrm{}^{}\overline{\mathrm{}}`$ and in the weak-$``$ sense in $`𝒳_\rho ^{0,\overline{\mathrm{}}}([T,T_1])`$ for every $`T_1`$, $`T<T_1<\mathrm{}`$, and the following estimate holds
$$Z(\chi \chi _{t_0};[T,t_0],1,\mathrm{})Cb\beta h_1(t_0).$$
$`(6.24)`$
(3) The solution $`\chi `$ is unique in $`L^{\mathrm{}}(I,L^{\mathrm{}})`$ under the condition that $`\chi (t)\psi _+_{\mathrm{}}`$ tends to zero as $`t\mathrm{}`$.
(4) Let in addition $`(\stackrel{~}{k},\overline{k},k)`$ satisfy (6.4), let $`w_+K_\rho _{_{\mathrm{}}}^{\stackrel{~}{k}}`$ and let $`V`$ be defined by Proposition 6.1 for some $`W`$ (for instance $`W(t)w_+`$). Then for fixed $`\varphi `$, $`V\mathrm{exp}(i\chi )`$ is gauge invariant in the following sense. If $`(V,\chi )`$ and $`(V^{},\chi ^{})`$ are the solutions obtained from $`(w_+,\psi _+)`$ and $`(w_+^{},\psi _+^{})`$ with $`w_+\mathrm{exp}(i\psi _+)=w_+^{}\mathrm{exp}(i\psi _+^{})`$, then $`V(t)\mathrm{exp}(i\chi (t))=V^{}(t)\mathrm{exp}(i\chi ^{}(t))`$ for all $`tI`$.
Proof. Parts (1) and (2). The proof is very similar to that of Proposition 6.1 and we concentrate again on the Gevrey estimates and on the limit $`t_0\mathrm{}`$. Let $`\chi _{t_0}`$ be the solution of (6.19) with initial data $`\chi _{t_0}(t_0)=\psi _+`$. In addition to (6.19), it is convenient to use also the equation
$$_t\tau =t^2\left(S\tau +\tau S\right)$$
$`(6.25)`$
satisfies by $`\tau =\chi `$, with $`S=\varphi `$.
We estimate $`_t|\chi _{t_0}|_\overline{\mathrm{}}^2`$ in the same way as in the proof of Lemma 3.5. We estimate the contribution of $`S\tau _{t_0}`$ from (6.25) by (3.25) with $`(s,s^{},\mathrm{},\mathrm{}^{})`$ replaced by $`(S,\tau _{t_0},\stackrel{~}{\mathrm{}},\overline{\mathrm{}})`$ and the contribution of $`\tau _{t_0}S`$ by (3.26) with $`(s,s^{},\mathrm{},\mathrm{}^{},m)`$ replaced by $`(\tau _{t_0},S,\overline{\mathrm{}}+\nu /2,\stackrel{~}{\mathrm{}},\overline{\mathrm{}}\nu /2)`$, thereby obtaining
$$_t|\chi _{t_0}|_\overline{\mathrm{}}^22\rho ^{}|\chi _{t_0}|_{\overline{\mathrm{}}+\nu /2}^2Ct^2|\chi _{t_0}|_{\overline{\mathrm{}}+\nu /2}^2|\varphi |_\stackrel{~}{\mathrm{}}$$
$`(6.26)`$
under the conditions
$$\stackrel{~}{\mathrm{}}>n/2\nu ,\overline{\mathrm{}}+1\nu /2,\stackrel{~}{\mathrm{}}\overline{\mathrm{}}+1\nu /2$$
which follow from (6.20). Introducing $`z_{(1)}=z_{(1)}(\chi _{t_0};[T,t_0],1,\overline{\mathrm{}})`$ and integrating (6.26) over time, we obtain in the same way as before
$$z^2z_1^2z_0^2+Cbh_1(T)z_1^2$$
with $`z_0=\beta `$, which under the conditions (6.7) implies
$$Z(\chi _{t_0};[T,t_0],1,\overline{\mathrm{}})C\beta .$$
$`(6.27)`$
We next estimate the difference $`\chi _1(t)\chi _{t_0}(t)\chi _{t_0}(t_0)\chi _{t_0}(t)\psi _+`$ which satisfies the equations
$$_t\chi _1=t^2\varphi \chi _{t_0}$$
$$_t\tau _1=t^2\left(S\tau _{t_0}+\tau _{t_0}S\right).$$
Let $`\stackrel{~}{\chi }_1=h_1^1\chi _1`$. Using again (3.26) with $`(s,s^{})=(S,\tau _{t_0})`$ or $`(\tau _{t_0},S)`$, with $`m=\mathrm{}\nu /2`$ and with $`(\mathrm{},\mathrm{}^{})=(\stackrel{~}{\mathrm{}},\overline{\mathrm{}}+\nu /2)`$ or $`(\overline{\mathrm{}}+\nu /2,\stackrel{~}{\mathrm{}})`$, we obtain
$$_t|\stackrel{~}{\chi }_1|_{\mathrm{}}^22\rho ^{}|\stackrel{~}{\chi }_1|_{\mathrm{}+\nu /2}^2Ct^2h_1^1|\stackrel{~}{\chi }_1|_{\mathrm{}+\nu /2}|\chi _{t_0}|_{\overline{\mathrm{}}+\nu /2}|\varphi |_\stackrel{~}{\mathrm{}}$$
$`(6.28)`$
under the conditions
$$\stackrel{~}{\mathrm{}}+\overline{\mathrm{}}>\mathrm{}+n/2\nu ,\stackrel{~}{\mathrm{}}\mathrm{}+1\nu /2,\overline{\mathrm{}}\mathrm{}+1\nu ,\mathrm{}+1\nu /2$$
which also follow from (6.20). Defining now $`z_{(1)}=z_{(1)}(\chi _1;[T,t_0],h_1,\mathrm{})`$, integrating over time and using (6.21) (6.27), we obtain in the same way as before
$$z^2z_1^2Cb\beta z_1\underset{t}{Sup}\left\{t^2\rho ^1h_0h_1^1\right\}=Cb\beta z_1$$
and therefore
$$Z(\chi _{t_0}\psi _+;[T,t_0],h_1,\mathrm{})Cb\beta .$$
$`(6.29)`$
Starting from the basic estimates (6.27) and (6.29), the end of the proof is the same as that of Proposition 6.1 based on (6.13) (6.16), with the simplification that the initial condition at $`t_0`$ is given by a fixed $`\psi _+`$ instead of a time dependent $`W`$. The difference between two solutions $`\chi _{t_1}`$ and $`\chi _{t_0}`$ with $`T<t_0<t_1`$ is estimated with the help of the extension of (6.27) to that difference and of the pointwise part of (6.29) with $`t_0`$ replaced by $`t_1`$ taken at time $`t=t_0`$ as
$$Z(\chi _{t_1}\chi _{t_0};[T,t_0],1,\mathrm{})C\left|\chi _{t_1}(t_0)\psi _+\right|_{\mathrm{}}Cb\beta h_1(t_0)$$
$`(6.30)`$
from which the existence of $`\chi `$ with the properties and convergences stated in Part (2) follow. Taking the limit $`t_1\mathrm{}`$ in (6.30) yields (6.24), while taking the limit $`t_0\mathrm{}`$ in (6.27) (6.29) yields (6.22) (6.23) in the same way as in the proof of Proposition 6.1.
Part (3) follows from elementary estimates together with the estimate on $`\varphi `$ expressed by (6.21).
Part (4). It follows from (6.3) and (6.19) that $`V\mathrm{exp}(i\chi )`$ also satisfies (6.3), with gauge invariant initial condition $`V(\mathrm{})\mathrm{exp}(i\chi (\mathrm{}))=w_+\mathrm{exp}(i\psi _+)`$. The result then follows from the uniqueness statement of Proposition 6.1, part (3).
$``$$``$
Remark 6.2. Because of the linearity of the equations (6.3) (6.19), the solutions $`V`$ and $`\chi `$ constructed in Propositions 6.1 and 6.2 have obvious continuity properties with respect to $`w_+`$ and $`\psi _+`$ respectively. The continuity with respect to $`\varphi `$ is more delicate and will not be considered here.
We shall use the results of Propositions 6.1 and 6.2 in the special case where $`\varphi =\varphi _{p1}`$ and $`W=W_p`$ as defined by (5.22). In that case, $`V`$ satisfies an asymptotic estimate which is much more accurate than (6.9) and which shows that $`V`$ is a good approximation to $`W_p`$. In order to state that result, we need the results of Proposition 5.2. We need in particular the special function $`\overline{h}_0`$ defined by (5.7) and the associated functions $`\overline{N}_m`$ and $`\overline{Q}_m`$ associated with it according to (3.68) (3.69). The result can be stated as follows.
Proposition 6.3. Let $`p1`$ be an integer. Let $`(k_+,\overline{k},k)`$ satisfy
$$k_+>n/2,\stackrel{~}{\mathrm{}}>n/2\nu ,\stackrel{~}{k}k+1\nu /2,\stackrel{~}{k}\overline{k}k+1\nu ,k\nu /2$$
$`(6.31)`$
where $`\stackrel{~}{\mathrm{}}(\mathrm{}_{p1})=k_+\lambda (p1)\overline{\lambda }`$ and $`\stackrel{~}{k}(k_p)=k_+p\overline{\lambda }`$. Let $`w_+K_\rho _{_{\mathrm{}}}^{k_+}`$ and let $`a_+=|w_+|_{k_+}`$. Let $`\overline{h}_1`$ and $`h_2`$ be $`𝒞^1`$ positive nonincreasing functions defined in $`[1,\mathrm{})`$ and tending to zero at infinity, with $`\overline{h}_1t^2\rho ^1\overline{h}_0`$ and $`h_2t^2\rho ^1\overline{N}_p`$. Let $`\varphi =\varphi _{p1}`$ and $`W=W_p`$ be defined by (5.22), and let
$$b=\underset{t}{Sup}\overline{h}_0(t)^1|\varphi |_\stackrel{~}{\mathrm{}}\underset{t}{Sup}\overline{h}_0(t)^1|\varphi _{p1}(t)|_{\mathrm{}_{p1}}.$$
$`(6.32)`$
Let $`T,1<T<\mathrm{}`$ satisfy
$$b\overline{h}_1(T)c$$
$`(6.33)`$
for a suitable constant $`c`$ (see (6.7)) and let $`V`$ be the solution of the equation (6.3) constructed in Proposition 6.1, and such that $`(V,0)𝒳_\rho ^{\overline{k},0}([T,\mathrm{}))`$. Then $`V`$ satisfies the following estimate
$$Y(VW_p;[T,\mathrm{}),h_2,k)A(a_+).$$
$`(6.34)`$
Proof. One checks easily that the assumptions of Proposition 6.3 imply the relevant assumptions of Propositions 5.2 and 6.1. In particular (6.31) implies (6.4) and the exponents $`\stackrel{~}{k}`$ and $`\stackrel{~}{\mathrm{}}`$ are those given by Proposition 5.2. Let now $`V_{t_0}`$ be defined in Proposition 6.1, part (2), let $`r=V_{t_0}W_p`$ so that in particular $`r(t_0)=0`$, and $`r`$ satisfies the equation
$$_tr=\left(2t^2\right)^1\left\{\left(2\varphi _{p1}+(\mathrm{\Delta }\varphi _{p1})\right)r+\underset{\genfrac{}{}{0pt}{}{ip1,jp}{i+jp}}{}\left(2\phi _i+(\mathrm{\Delta }\phi _i)\right)w_j\right\}$$
$`(6.35)`$
obtained by taking the difference between (6.3) and the appropriate sum of (2.18). Let $`\stackrel{~}{r}=h_2^1r`$. We now estimate $`_t|\stackrel{~}{r}|_k^2`$. The contribution of the terms containing $`r`$ in the RHS are estimated exactly as in the proof of Proposition 6.1. The remaining terms are estimated by (3.23) (3.24) with $`m=k\nu /2`$, $`k\stackrel{~}{k}`$ and $`\mathrm{}\stackrel{~}{\mathrm{}}`$ under the conditions $`\stackrel{~}{\mathrm{}}>n/2\nu `$, $`0k\nu /2\stackrel{~}{k}1`$, which follow from (6.31). We obtain
$$_t|\stackrel{~}{r}|_k^22\rho ^{}|\stackrel{~}{r}|_{k+\nu /2}^2Ct^2|\stackrel{~}{r}|_{k+\nu /2}^2|\varphi |_\stackrel{~}{\mathrm{}}Ct^2h_2^1|\stackrel{~}{r}|_{k+\nu /2}|\phi _i|_\stackrel{~}{\mathrm{}}|w_j|_{\stackrel{~}{k}}$$
$$2\rho ^{}|\stackrel{~}{r}|_{k+\nu /2}^2Cbt^2\overline{h}_0(t)|\stackrel{~}{r}|_{k+\nu /2}^2A(a_+)t^2h_2^1\overline{N}_p|\stackrel{~}{r}|_{k+\nu /2}$$
$`(6.36)`$
by (6.32) (5.9) (5.10) (3.70) and (3.79). Defining $`y_{(1)}=y_{(1)}(r,[T,t_0],h_2,k)`$ and integrating (6.36) over time with the initial condition $`r(t_0)=0`$, we obtain in the same way as before
$$y^2y_1^2Cb\overline{h}_1(t)y_1^2+A(a_+)y_1$$
by using the properties of $`\overline{h}_1`$ and $`h_2`$, and therefore, under the condition (6.33)
$$Y(r;[T,t_0],h_2,k)A(a_+).$$
$`(6.37)`$
We now take $`T<t_0<t_1`$, and take the limit $`t_1\mathrm{}`$, $`t_0\mathrm{}`$ in that order in the estimate
$$Y(V_{t_1}W;[T,t_0],h_2,k)A(a_+)$$
which follows from (6.37). The estimate (6.34) then follows from the convergence (6.10).
$``$$``$
Remark 6.2. In order to appreciate the improvement of the asymptotic accuracy of Proposition 6.3, especially (6.34) over Proposition 6.1, especially (6.9), it is useful to consider again the special case $`\rho ^{}=t^{1\epsilon }`$. Proposition 6.1 with $`\varphi =\varphi _{p1}`$ is now applied with $`h_0=\overline{h}_0t^{1\gamma }`$ and therefore $`h_1=\overline{h}_1t^{\gamma +\epsilon }`$, so that the pointwise part of (6.9) states that
$$|V(t)w_+|_k<Cabt^{\gamma +\epsilon }.$$
On the other hand $`\overline{N}_pt^{1(p+1)\gamma }`$ for $`(p+1)\gamma <1`$ thereby allowing for $`h_2t^2\rho ^1\overline{N}_pt^{(p+1)\gamma +\epsilon }`$ in Proposition 6.3, so that the pointwise part of (6.34) states that
$$|V(t)W_p(t)|_k<A(a_+)t^{(p+1)\gamma +\epsilon }.$$
That improvement will play an essential role in the estimates of Proposition 6.4 below.
We now turn to the construction of solutions $`(w,\phi )`$ of the system (2.11)-(2.12) with given asymptotic states $`(w_+,\psi _+)`$. For that purpose, we first take a large positive $`t_0`$ and we construct a solution $`(w_{t_0},\phi _{t_0})`$ of (2.11)-(2.12) with initial data $`(V(t_0),\varphi _p(t_0)+\chi (t_0))`$ at $`t_0`$. The solution $`(w,\phi )`$ will be obtained therefrom by taking the limit $`t_0\mathrm{}`$, as explained in Section 2.
Proposition 6.4. Let $`(k,\mathrm{})`$ satisfy (3.48) and $`k1\nu /2`$, let $`p`$ be an integer such that $`(p+2)\gamma >1`$ and let $`k_+`$ and $`\mathrm{}_+`$ satisfy
$$k_+>n/2,k_+(k+2\nu )(\mathrm{}+\lambda +1)+p\overline{\lambda },\mathrm{}_+\mathrm{}+1,$$
$`(6.38)`$
where $`\lambda =\mu n+2`$ and $`\overline{\lambda }=\lambda 1`$. Let $`\overline{h}_0`$ be defined by (5.7) and let $`\overline{N}_m`$, $`\overline{Q}_m`$ be the associated estimating functions defined by (3.68) (3.69). Let $`\overline{h}_1`$, $`h_1`$, $`h_2`$ and $`h_3`$ be positive nonincreasing $`𝒞^1`$ functions defined in $`[1,\mathrm{})`$, tending to zero at infinity, and satisfying
$$\overline{h}_1t^2\rho ^1\overline{h}_0,h_2t^2\rho ^1\overline{N}_p,h_2\overline{Q}_p\mathrm{𝑖𝑓}p1,$$
$`(6.39)`$
$$h_3t^\gamma \rho ^1h_2,h_1t^2\rho ^1h_3h_2^1,h_3C\overline{h}_1.$$
$`(6.40)`$
Let $`w_+K_\rho _{_{\mathrm{}}}^{k_+}`$ and let $`(W_m,\varphi _m)`$, $`0mp`$, be defined by (5.22) so that $`(W_m,\overline{h}_0^1\varphi _m)𝒴_\rho _{_{\mathrm{}}}^{k_m,\mathrm{}_m}([1,\mathrm{}))`$ by Proposition 5.2. Let $`V`$ be defined by Proposition 6.1 with $`(W,\varphi )=(W_p,\varphi _{p1})`$ and $`(\stackrel{~}{k},\stackrel{~}{\mathrm{}})=(k_p,\mathrm{}_{p1})`$. Let $`\psi _+Y_\rho _{_{\mathrm{}}}^\mathrm{}_+`$ and let $`\chi `$ be defined by Proposition 6.2 with the same $`(W,\varphi )`$. Let
$$a_+=|w_+|_{\rho _{_{\mathrm{}}},k_+},b_+=|\psi _+|_{\rho _{_{\mathrm{}}},\mathrm{}_+}.$$
Then there exists $`T`$, $`1T<\mathrm{}`$, depending only on $`(\gamma ,p,a_+,b_+)`$ such that for all $`t_0T`$, the system (2.11)-(2.12) with initial data $`w(t_0)=V(t_0)`$, $`\phi (t_0)=\varphi _p(t_0)+\chi (t_0)`$ has a unique solution $`(w_{t_0},\phi _{t_0})𝒳_\rho ^{k,\mathrm{}}([T,t_0])`$. One can define $`T`$ by a condition of the type
$$A(a_+,b_+)\left(h_1(T)\overline{h}_1(T)h_2(T)\right)=1$$
$`(6.41)`$
and the solution satifies the estimates
$$Y(w_{t_0}V;[T,t_0],h_2,k)Y(w_{t_0}W_p;[T,t_0],h_2,k)A(a_+,b_+)$$
$`(6.42)`$
$$Z(\phi _{t_0}\varphi _p\chi ;[T,t_0],h_3,\mathrm{})Z(\phi _{t_0}\varphi _p\psi _+;[T,t_0],h_3,\mathrm{})A(a_+,b_+)$$
$`(6.43)`$
$$Y(w_{t_0};[T,t_0],1,k)A(a_+,b_+),$$
$`(6.44)`$
$$Z(\phi _{t_0};[T,t_0],\overline{h}_0,\mathrm{})A(a_+,b_+).$$
$`(6.45)`$
Remark 6.3. As mentioned previously, in the same way as in Propositions 6.1 and 6.2, we could easily (but we shall not) extend the solution $`(w_{t_0},\phi _{t_0})`$ to the interval $`[T,\mathrm{})`$.
Remark 6.4. In order to understand the time decay estimates implied by (6.42) (6.43), it is useful to consider again the special case $`\rho ^{}=t^{1\epsilon }`$. Saturating as far as possible the inequalities in (6.39) (6.40), we obtain
$$\overline{h}_1t^{\gamma +\epsilon },h_2t^{(p+1)\gamma +\epsilon },h_3t^{1(p+2)\gamma +2\epsilon },h_1=t^{\gamma +2\epsilon }$$
for $`(p+1)\gamma <1`$, and the assumptions are satisfied for $`\epsilon `$ sufficiently small, namely $`2\epsilon <(p+2)\gamma 1`$. In particular the condition that $`h_3`$ be decreasing in $`t`$ essentially imposes the condition $`(p+2)\gamma >1`$.
Note also that in the condition (6.41) $`h_2`$ yields only a marginal restriction as soon as $`p1`$, while for $`p=0`$ it is natural to take $`h_2=\overline{h}_1`$ (see (6.39) with $`\overline{N}_0=\overline{h}_0)`$. Finally, the condition $`h_3C\overline{h}_1`$ will in general be automatically satisfied for any reasonable choice of the estimating functions $`\overline{h}_1`$ and $`h_3`$.
Proof. The proof follows the same pattern as that of Proposition 4.1, involving a parabolic regularization, possibly a regularization of the initial data, the local resolution of the regularized system by a fixed point method, the derivation of a priori estimates uniform in the regularization, and a limiting procedure. The only difference lies in the a priori estimates of the solutions in $`𝒳_\rho ^{k,\mathrm{}}([T,t_0])`$. Those estimates are much more elaborate than previously and will ensure in particular that $`T`$ can be taken independent of $`t_0`$, contrary to what happened in Lemma 4.1 (see especially (4.13)). We concentrate on the proof of those estimates, omitting the parabolic regularization terms for brevity. Their contribution will be briefly discussed at the end of that proof.
Let $`(w_{t_0},\phi _{t_0})𝒳_\rho ^{k,\mathrm{}}([T,t_0])`$ be a solution of the system (2.11)-(2.12) with initial condition $`(w_{t_0}(t_0),\phi _{t_0}(t_0))=(V(t_0),\varphi _p(t_0)+\chi (t_0))`$ where $`V`$ and $`\chi `$ are defined by Propositions 6.1 and 6.2. Instead of estimating $`(w_{t_0},\phi _{t_0})`$ directly, we estimate the differences $`q=w_{t_0}V`$ and $`\psi =\phi _{t_0}\varphi _p\chi `$. For convenience we also introduce the gradients $`\sigma =\psi `$, $`s_{t_0}=\phi _{t_0}`$, $`\tau =\chi `$, as well as $`s_m=\phi _m`$, $`S_m=\varphi _m`$ for $`0mp`$. Comparing the equations (2.11)-(2.12) and (6.3) (6.19) with $`\varphi =\varphi _{p1}`$, we obtain
$$_tq=\left(2t^2\right)^1\left\{i\mathrm{\Delta }w_{t_0}+2s_{t_0}q+2(\sigma +s_p+\tau )V+(\sigma )w_{t_0}+((S_p+\tau ))q+((s_p+\tau ))V\right\}$$
$`(6.46)`$
$$_t\psi =\left(2t^2\right)^1\left\{|\sigma |^2+2\sigma (S_p+\tau )+|\tau |^2+2\tau s_p+\underset{\genfrac{}{}{0pt}{}{i,jp}{i+jp}}{}s_is_j\right\}$$
$$+t^\gamma \left\{g_0(q,q)+2g_0(q,V)+q_0(VW_p,V+W_p)+\underset{\genfrac{}{}{0pt}{}{i,jp}{i+j>p}}{}g_0(w_i,w_j)\right\}.$$
$`(6.47)`$
It is convenient to write also the equation for $`\sigma =\psi `$, namely
$$_t\sigma =t^2\left\{s_{t_0}\sigma +\sigma (S_p+\tau )+(\tau +s_p)\tau +\tau s_p+\underset{\genfrac{}{}{0pt}{}{i,jp}{i+jp}}{}s_is_j\right\}$$
$$+t^\gamma \left\{g_0(q,q)+2g_0(q,V)+q_0(VW_p,V+W_p)+\underset{\genfrac{}{}{0pt}{}{i,jp}{i+j>p}}{}g_0(w_i,w_j)\right\}.$$
$`(6.48)`$
We define $`\stackrel{~}{q}=h_2^1q`$ and $`\stackrel{~}{\psi }=h_3^1\psi `$ and we estimate $`_t|\stackrel{~}{q}|_k^2`$ and $`_t|\stackrel{~}{\psi }|_{\mathrm{}}^2`$ by exactly the same method as in Lemmas 3.5, 3.6, 3.7, based on Lemma 3.4, using in particular (3.22) (3.23) (3.24) for $`\stackrel{~}{q}`$ and (3.25) (3.26) (3.28) for $`\stackrel{~}{\psi }`$, and omitting the terms with $`h_2^{}`$ and $`h_3^{}`$. We obtain for $`tt_0`$
$$_t|\stackrel{~}{q}|_k^22\rho ^{}|\stackrel{~}{q}|_{k+\nu /2}^2Ct^2|\stackrel{~}{q}|_{k+\nu /2}\{h_2^1|V|_{k+2\nu /2}$$
$$+|\stackrel{~}{q}|_{k+\nu /2}|\psi +\varphi _p+\chi |_{\mathrm{}}+h_2^1|V|_{k+1\nu /2}|\psi +\phi _p+\chi |_{\mathrm{}}$$
$$+(|\stackrel{~}{q}|_k+h_2^1|V|_k)|\psi |_{\mathrm{}+\nu /2}+|\stackrel{~}{q}|_k|\varphi _p+\chi |_{\mathrm{}+\nu /2}+h_2^1|V|_k|\phi _p+\chi |_{\mathrm{}+\nu /2}\}$$
$`(6.49)`$
$$_t|\stackrel{~}{\psi }|_{\mathrm{}}^22\rho ^{}|\stackrel{~}{\psi }|_{\mathrm{}+\nu /2}^2Ct^2|\stackrel{~}{\psi }|_{\mathrm{}+\nu /2}\{|\stackrel{~}{\psi }|_{\mathrm{}+\nu /2}|\psi +\varphi _p+\chi |_{\mathrm{}}$$
$$+|\stackrel{~}{\psi }|_{\mathrm{}}|\varphi _p+\chi |_{\mathrm{}+1\nu /2}+h_3^1|\phi _p+\chi |_{\mathrm{}}|\chi |_{\mathrm{}+1\nu /2}+h_3^1|\chi |_{\mathrm{}}|\phi _p|_{\mathrm{}+1\nu /2}$$
$$+h_3^1\underset{\genfrac{}{}{0pt}{}{i,jp}{i+jp}}{}|\phi _i|_{\mathrm{}}|\phi _j|_{\mathrm{}+1\nu /2}\}Ct^\gamma h_3^1|\stackrel{~}{\psi }|_{\mathrm{}+\nu /2}\{h_2|\stackrel{~}{q}|_{k+\nu /2}$$
$$\left(|q|_k+|V|_k\right)+h_2|\stackrel{~}{q}|_k|V|_{k+\nu /2}+|VW_p|_k|V+W_p|_{k+\nu /2}$$
$$+|VW_p|_{k+\nu /2}|V+W_p|_k+\underset{\genfrac{}{}{0pt}{}{i,jp}{i+j>p}}{}|w_i|_k|w_j|_{k+\nu /2}\}.$$
$`(6.50)`$
In order to continue the estimates, we need some information on $`V`$, $`\chi `$ and on the $`(\phi _m,w_m)`$. Applying Proposition 6.1 with $`\varphi =\varphi _{p1}`$ and Proposition 6.3, both with $`\overline{k}=k+2\nu `$, we rewrite (6.8) and (6.34) as
$$Y(V;[T,\mathrm{}),1,k+2\nu )a$$
$`(6.51)`$
$$Y(VW_p;[T,\mathrm{}),h_2,k)a$$
$`(6.52)`$
for some $`a`$ depending on $`a_+`$, under the conditions
$$k_+>n/2,\mathrm{}_{p1}>n/2\nu ,k_pk+2\nu $$
$`(6.53)`$
where (cf (5.8))
$$k_m=k_+m\overline{\lambda },\mathrm{}_m=k_+\lambda m\overline{\lambda }.$$
Similarly, applying Proposition 6.2 with $`\varphi =\varphi _{p1}`$, $`\overline{\mathrm{}}=\mathrm{}+1`$, $`h_0=\overline{h}_0`$ and therefore $`h_1=\overline{h}_1`$, we rewrite (6.22) (6.23) as
$$Z(\chi ;[T,\mathrm{}),1,\mathrm{}+1)b$$
$`(6.54)`$
$$Z(\chi \psi _+;[T,\mathrm{}),\overline{h}_1,\mathrm{})b$$
$`(6.55)`$
for some $`b`$ depending on $`a_+`$ and $`b_+`$, under the conditions
$$\mathrm{}_{p1}>n/2\nu ,\mathrm{}_{p1}\mathrm{}+2\nu /2,\mathrm{}_+\mathrm{}+1.$$
$`(6.56)`$
Finally, from Proposition 5.2, we obtain
$$|w_j(t)|_{k+\nu /2}a\overline{Q}_{j1}(t)\mathrm{for}1jp,$$
$`(6.57)`$
$$|W_p(t)|_{k+\nu /2}a,$$
$`(6.58)`$
for some $`a`$ depending on $`a_+`$, under the conditions
$$k_+>n/2,k_pk+\nu /2,$$
$`(6.59)`$
and
$$|\phi _j(t)|_{\mathrm{}+1\nu /2}b\overline{N}_j(t)\mathrm{for}0jp,$$
$`(6.60)`$
$$|\varphi _p(t)|_{\mathrm{}+1\nu /2}b\overline{h}_0(t)$$
$`(6.61)`$
for some $`b`$ depending on $`a_+`$, under the conditions
$$k_+>n/2,\mathrm{}_p\mathrm{}+1\nu /2.$$
$`(6.62)`$
In the estimates (6.51) (6.52) (6.54) (6.55) (6.57) (6.58) (6.60) (6.61), we use two common letters $`a`$ and $`b`$ to refer to estimates of amplitudes and phases respectively. The constant $`a`$ depends only on $`a_+`$, while $`b`$ depends on $`a_+`$ and on $`b_+`$. The conditions (6.53) (6.56) (6.59) (6.62) are implied by (6.38), which is the statement of
$$k_+>n/2,k_pk+2\nu ,\mathrm{}_p\mathrm{}+1,\mathrm{}_+\mathrm{}+1.$$
With the estimates (6.51)-(6.61) available, we continue to estimate $`(q,\psi )`$ by integrating (6.49) (6.50) over time in the interval $`[t,t_0]`$. We define $`y_{(1)}=y_{(1)}(q;[t,t_0],h_2,k)`$ and $`z_{(1)}=z_{(1)}(\psi ;[t,t_0],h_3,\mathrm{})`$. We proceed exactly as in the proofs of Lemmas 4.1, 4.2 and 4.3, using the Schwarz inequality for the time integrals whenever necessary. We use furthermore the fact that (6.49) (resp. (6.50)) contains $`|\stackrel{~}{q}|_{k+\nu /2}`$ (resp. $`|\stackrel{~}{\psi }|_{\mathrm{}+\nu /2})`$ as a factor in its RHS, thereby yielding a factor $`y_1`$ (resp. $`z_1`$) after integration, and the elementary fact that
$$y^2y_1^2Ay_1yy_1A$$
and its analogue for $`(z,z_1)`$. We then obtain the following estimates, where Sup means that the Supremum of the function of time that follows is taken in the interval $`[t,t_0]`$, and where an overall constant $`C`$ is omitted for brevity.
$$yy_1a\mathrm{Sup}\left(t^2\rho ^1h_2^1\right)+ab\mathrm{Sup}\left(t^2\rho ^1\overline{N}_ph_2^1\right)$$
$$+b(y+y_1)\mathrm{Sup}\left(t^2\rho ^1\overline{h}_0\right)+a(z+z_1)\mathrm{Sup}\left(t^2\rho ^1h_2^1h_3\right)$$
$$+(yz_1+y_1z)\mathrm{Sup}\left(t^2\rho ^1h_3\right),$$
$`(6.63)`$
$$zz_1b^2\mathrm{Sup}\left(t^2\rho ^1\overline{h}_0h_3^1\overline{N}_p\right)+b(z+z_1)\mathrm{Sup}\left(t^2\rho ^1\overline{h}_0\right)$$
$$+zz_1\mathrm{Sup}\left(t^2\rho ^1h_3\right)+a(y+y_1)\mathrm{Sup}\left(t^\gamma \rho ^1h_3^1h_2\right)$$
$$+yy_1\mathrm{Sup}\left(t^\gamma \rho ^1h_3^1h_2^2\right)+pa^2\mathrm{Sup}\left(t^\gamma \rho ^1h_3^1(h_2+\overline{Q}_p)\right)$$
$`(6.64)`$
where the factor $`p`$ in the last term simply means that that term is absent for $`p=0`$. The various Sup in time are estimated in an obvious way with the help of the conditions (6.39) (6.40) which are taylored for that purpose. The only non-obvious term is the coefficient of $`b^2`$ in (6.64), which is estimated by (6.39) (6.40) as
$$\left(t^2\rho ^1\overline{N}_ph_2^1\right)\left(t^\gamma \rho ^1h_2h_3^1\right)\left(t^\gamma \rho ^{}\overline{h}_0\right)\rho ^{}_1$$
since
$$\overline{h}_0=_1^t𝑑t_1t_1^\gamma \rho ^{}(t_1)^1\rho ^{}(t_1)t^\gamma \rho ^{}(t)^1_1^t𝑑t_1\rho ^{}(t_1)\rho ^{}_1t^\gamma \rho ^1$$
$`(6.65)`$
by the monotony of $`t^\gamma \rho ^1`$. Absorbing the factor $`\rho ^{}_1`$ in the (again omitted) overall constant and defining as previously $`Y=yy_1`$ and $`Z=zz_1`$, we end up with
$$Ya+ab+bY\overline{h}_1+aZh_1+YZh_1h_2$$
$`(6.66)`$
$$Zb^2+pa^2+bZ\overline{h}_1+Z^2h_1h_2+aY+Y^2h_2$$
$`(6.67)`$
where the functions $`h_1`$, $`\overline{h}_1`$ and $`h_2`$ are taken at time $`t`$ where they take their Supremum in $`[t,t_0]`$, since they are assumed to be decreasing.
In order to conclude the estimates, we impose the conditions
$$4b\overline{h}_1(T)1,4(1+b)h_2(T)1,Zh_1(t)(1+b)$$
$`(6.68)`$
which together imply $`4Zh_1(t)h_2(t)1`$ for all $`tT`$. We then obtain
$$Y2a(1+b)+2aZh_14a(1+b)$$
$`(6.69)`$
$$Z2(b^2+pa^2)+2aY+2Y(Yh_2)2(b^2+pa^2)+4aY2(b^2+pa^2)+16a^2(1+b).$$
$`(6.70)`$
The condition $`Zh_11+b`$ then reduces to
$$\left(2(b^2+pa^2)(1+b)^1+16a^2\right)h_11$$
which is implied by
$$\left(2b+(16+2p)a^2\right)h_1(T)1.$$
$`(6.71)`$
The condition (6.71) together with the first two conditions of (6.68) then take the form (6.41), while the estimates (6.69) (6.70) yield the estimate of the first terms in the LHS of (6.42) (6.43). The estimates of the second terms follow from those of the first ones and from (6.52) and (6.55), together with the condition $`h_3C\overline{h}_1`$ from (6.40).
Finally the estimates (6.44) (6.45) follow immediately from (6.42) (6.43), from (6.51) or (6.58) and from (6.61) (6.54).
We now discuss briefly the contribution of the parabolic regularization terms in the previous proof. We regularize the system (2.11)-(2.12) in the same way as in the proof of Proposition 4.1 (see (4.51), where however the sign of the regularizing terms should be changed since we are now solving the equations for decreasing $`t`$ starting from $`t_0`$). Instead of (6.49) (6.50), we then obtain
$$_t|\stackrel{~}{q}|_k^22\theta (|\stackrel{~}{q}|_k^2+h_2^1\mathrm{Re}<\stackrel{~}{q},V>_k)+\text{previous terms},$$
$$_t|\stackrel{~}{\psi }|_{\mathrm{}}^22\theta (|\stackrel{~}{\psi }|_{\mathrm{}}^2+h_3^1<\stackrel{~}{\psi },(\varphi _p+\chi )>_{\mathrm{}})+\text{previous terms}$$
where $`<,>_k`$ and $`<,>_{\mathrm{}}`$ denote the scalar products in $`K_\rho ^k`$ and $`Y_\rho ^{\mathrm{}}`$. The scalar products are controlled since we have assumed that $`(V,\varphi _p+\chi )𝒳_{\rho ,loc}^{k+1,\mathrm{}+1}`$ by imposing $`\overline{k}=k+2\nu `$, $`\overline{\mathrm{}}=\mathrm{}+1`$ and $`\mathrm{}_p\mathrm{}+1`$. They produce additional terms in the final estimates which are uniformly bounded in $`\theta `$ in a neighborhood of zero.
With the a priori estimates (6.44) (6.45) replacing Lemma 4.1, the proof of Proposition 6.4 proceeds in the same way as that of Proposition 4.1, as mentioned above. In particular, the required regularity estimates and difference estimates are provided by Parts 3 of Lemmas 4.2 and 4.3, which are taylored for that purpose. The assumption on $`(w,\phi )`$ made in those parts follow from (6.44) (6.45) and from the relation
$$\overline{h}_0\rho ^{}_1h_0$$
which follows from (6.65), while the regularity assumptions on the initial data required in Lemma 4.2 are ensured by the previous regularity of $`V`$ and $`\varphi _p+\chi `$.
$``$$``$
We can now take the limit $`t_0\mathrm{}`$ of the solution $`(w_{t_0},\phi _{t_0})`$ constructed in Proposition 6.4, for fixed $`(w_+,\psi _+)`$.
Proposition 6.5. Let the assumptions of Proposition 6.4 be satisfied. Then
(1) There exists $`T`$, $`1T<\mathrm{}`$, depending only on $`(\gamma ,p,a_+,b_+)`$ and there exists a unique solution $`(w,\phi )`$ of the system (2.11)-(2.12) in the interval $`[T,\mathrm{})`$ such that $`(w,\overline{h}_0^1\phi )𝒳_\rho ^{k,\mathrm{}}([T,\mathrm{}))`$ and such that the following estimates hold
$$Y(wV;[T,\mathrm{}),h_2,k)Y(wW_p;[T,\mathrm{}),h_2,k)A(a_+,b_+),$$
$`(6.72)`$
$$Z(\phi \varphi _p\chi ;[T,\mathrm{}),h_3,\mathrm{})Z(\phi \varphi _p\psi _+;[T,\mathrm{}),h_3,\mathrm{})A(a_+,b_+),$$
$`(6.73)`$
$$Y(w;[T,\mathrm{}),1,k)A(a_+,b_+),$$
$`(6.74)`$
$$Z(\phi ;[T,\mathrm{}),\overline{h}_0,\mathrm{})A(a_+,b_+).$$
$`(6.75)`$
One can define $`T`$ by a condition of the type (6.41).
(2) Let $`(w_{t_0},\phi _{t_0})𝒳_\rho ^{k,\mathrm{}}([T,t_0])`$ be the solution of the system (2.11)-(2.12) constructed in Proposition 6.4. Then $`(w_{t_0},\phi _{t_0})`$ converges to $`(w,\phi )`$ in norm in $`𝒳_\rho ^{k^{},\mathrm{}^{}}([T,T_1])`$ for $`k^{}<k`$, $`\mathrm{}^{}<\mathrm{}`$ and in the weak-$``$ sense in $`𝒳_\rho ^{k,\mathrm{}}([T,T_1])`$ for any $`T_1`$, $`T<T_1<\mathrm{}`$, and in the weak-$``$ sense in $`K_\rho ^kY_\rho ^{\mathrm{}}`$ pointwise in $`t`$ for all $`tT`$.
(3) The map $`(w_+,\psi _+)(w,\phi )`$ defined in Part (1) is continuous from the norm topology of $`(w_+,\psi _+)`$ in $`K_\rho _{_{\mathrm{}}}^{k_+}Y_\rho _{_{\mathrm{}}}^\mathrm{}_+`$ to the norm topology of $`(w,\phi )`$ in $`𝒳_\rho ^{k^{},\mathrm{}^{}}([T,T_1])`$ for $`k^{}<k`$, $`\mathrm{}^{}<\mathrm{}`$ and to the weak-$``$ topology in $`𝒳_\rho ^{k,\mathrm{}}([T,T_1])`$ for any $`T_1`$, $`T<T_1<\mathrm{}`$, and to the weak-$``$ topology in $`K_\rho ^kY_\rho ^{\mathrm{}}`$ pointwise in $`t`$ for all $`tT`$.
Remark 6.5. For simplicity, we have not stated the strongest continuity properties that would follow by tracking more accurately the exponents in the proof of Part (3). Actually the required topology on $`(w_+,\psi _+)`$ could be weakened to the norm topology of $`K_\rho _{_{\mathrm{}}}^k^{}Y_\rho _{_{\mathrm{}}}^{\mathrm{}^{}}`$ on the bounded sets of $`K_\rho _{_{\mathrm{}}}^{k_+}Y_\rho _{_{\mathrm{}}}^\mathrm{}_+`$ for suitable $`(k^{},\mathrm{}^{})`$ smaller than $`(k_+,\mathrm{}_+)`$.
Proof. Parts (1) and (2) will follow from the convergence of $`(w_{t_0},\phi _{t_0})`$ when $`t_0\mathrm{}`$ in the topologies stated in Part (2). Let $`Tt_0t_1`$, let $`(w_{t_0},\phi _{t_0})`$ and $`(w_{t_1},\phi _{t_1})`$ be the corresponding solutions of the system (2.11)-(2.12) obtained in Proposition 6.4, and let $`(w_{},\phi _{})=(w_{t_0}w_{t_1},\phi _{t_0}\phi _{t_1})`$. From (6.42) (6.43) and their analogues for $`(w_{t_1},\phi _{t_1})`$, it follows that
$$\{\begin{array}{c}Y(w_{};[T,t_0],h_2,k)A(a_+,b_+)\hfill \\ \\ Z(\phi _{};[T,t_0],h_3,\mathrm{})A(a_+,b_+)\hfill \end{array}$$
$`(6.76)`$
so that in particular
$$\{\begin{array}{c}\left|w_{}(t_0)\right|_kA(a_+,b_+)h_2(t_0)\hfill \\ \\ \left|\phi _{}(t_0)\right|_{\mathrm{}}A(a_+,b_+)h_3(t_0).\hfill \end{array}$$
$`(6.77)`$
We now apply Lemma 4.3, part (3) to $`(w_{},\phi _{})`$. For that purpose we take $`h_0=t^\gamma \rho ^1`$, so that by (6.40)
$$h_1t^2\rho ^1h_0,h_2h_0h_3$$
$`(6.78)`$
and in particular $`(h_0,h_1)`$ satisfy the assumptions of Lemma 4.3. The assumptions (4.35) (restricted to the relevant interval $`[T,t_0])`$ and (4.38) follow from (6.44) (6.45) and from (6.65), while the condition (4.26) can be included in (6.41). We now apply (4.41) with $`k^{}=k1+\nu `$ $`\mathrm{}^{}=\mathrm{}1+\nu `$, together with (6.77) (6.78), thereby obtaining
$$Y(w_{};[T,t_0],h_0^1,k1+\nu )A(a_+,b_+)\left\{h_0(t_0)\right|w_{}(t_0)|_{k1+\nu }+h_1(T)|\phi _{}(t_0)|_{\mathrm{}1+\nu }\}$$
$$A(a_+,b_+)h_3(t_0)$$
$`(6.79)`$
$$Z(\phi _{};[T,t_0],1,\mathrm{}1+\nu )A(a_+,b_+)\left\{|\phi _{}(t_0)|_{\mathrm{}1+\nu }+h_0(t_0)|w_{}(t_0)|_{k1+\nu }\right\}$$
$$A(a_+,b_+)h_3(t_0).$$
$`(6.80)`$
From (6.79) (6.80) and from the fact that $`h_3`$ tends to zero at infinity, it follows that there exists $`(w,\phi )𝒳_{\rho ,loc}^{k1+\nu ,\mathrm{}1+\nu }([T,\mathrm{}))`$ such that $`(w_{t_0},\phi _{t_0})`$ converges to $`(w,\phi )`$ in norm in $`𝒳_\rho ^{k1+\nu ,\mathrm{}1+\nu }([T,T_1])`$ for all $`T_1`$, $`T<T_1<\mathrm{}`$. From that convergence, from (6.42)-(6.45) and from standard compactness, continuity and interpolation arguments, it follows that $`(w,\overline{h}_0^1\phi )𝒳_\rho ^{k,\mathrm{}}([T,\mathrm{}))`$, that $`(w,\phi )`$ satisfies the estimates (6.72)-(6.75) and that $`(w_{t_0},\phi _{t_0})`$ converges to $`(w,\phi )`$ in the other topologies stated in Part (2). Furthermore $`(w,\phi )`$ satisfies the system (2.11)-(2.12), and uniqueness of $`(w,\phi )`$ under the conditions (6.72) (6.73) follows from Proposition 4.3 and from the fact that $`h_3(t)`$ tends to zero at infinity.
Part (3). Let $`(w_+,\psi _+)`$ and $`(w_+^{},\psi _+^{})`$ belong to a fixed bounded set of $`K_\rho _{_{\mathrm{}}}^{k_+}Y_\rho _{_{\mathrm{}}}^\mathrm{}_+`$, so that (5.17) (5.18) and (6.72)-(6.75) hold with fixed $`A`$. Let $`(W_p,\varphi _p)`$ and $`(W_p^{},\varphi _p^{})`$ be the associated functions defined by (5.22) and let $`(w,\phi )`$ and $`(w^{},\phi ^{})`$ be the associated solutions of the system (2.11)-(2.12) defined in Part (1). We assume that $`(w_+^{},\psi _+^{})`$ is close to $`(w_+,\psi _+)`$ in the sense that
$$|w_+w_+^{}|_{k_+}\epsilon ,|\psi _+\psi _+^{}|_{\mathrm{}}\epsilon _0.$$
$`(6.81)`$
Let $`w_{}=ww^{}`$, $`\phi _{}=\phi \phi ^{}`$ and let $`t_0>T`$ be defined by $`h_2(t_0)=\epsilon `$, which can be done for any (sufficiently small) $`\epsilon >0`$ and which implies that $`t_0\mathrm{}`$ when $`\epsilon 0`$. It follows from (6.72) and (5.17) that
$$|w_{}(t_0)|_kAh_2(t_0)+|W_p(t_0)W_p^{}(t_0)|_kA\left(h_2(t_0)+\epsilon \right)Ah_2(t_0)$$
$`(6.82)`$
and from (6.73) and (5.18) that
$$\begin{array}{cc}|\phi _{}(t_0)|_{\mathrm{}}\hfill & Ah_3(t_0)+|\varphi _p(t_0)\varphi _p^{}(t_0)|_{\mathrm{}}+|\psi _+\psi _+^{}|_{\mathrm{}}\hfill \\ & \\ & A\left(h_3(t_0)+\overline{h}_0(t_0)\epsilon \right)+|\psi _+\psi _+^{}|_{\mathrm{}}Ah_3(t_0)+\epsilon _0.\hfill \end{array}$$
$`(6.83)`$
We now apply Lemma 4.3, part (3) with $`k^{}=k1+\nu `$, $`\mathrm{}^{}=\mathrm{}1+\nu `$, thereby obtaining
$$\begin{array}{cc}Y(w_{};[T,t_0],h_0^1,k1+\nu )\hfill & A\left(h_0(t_0)h_2(t_0)+h_1(T)(h_3(t_0)+\epsilon _0)\right)\hfill \\ & \\ & A\left(h_3(t_0)+\epsilon _0\right)\hfill \end{array}$$
$`(6.84)`$
$$\begin{array}{cc}Z(\phi _{};[T,t_0],1,\mathrm{}1+\nu )\hfill & A\left\{h_3(t_0)+\epsilon _0+h_0(t_0)h_2(t_0)\right\}\hfill \\ & \\ & A\left(h_3(t_0)+\epsilon _0\right).\hfill \end{array}$$
$`(6.85)`$
When $`\epsilon `$ tends to zero, $`h_3(t_0)`$ tends to zero, and therefore (6.84) (6.85) imply norm continuity in $`𝒳_\rho ^{k1+\nu ,\mathrm{}1+\nu }([T,T_1])`$ for all $`T_1`$, $`T<T_1<\mathrm{}`$. The other continuities follow therefrom, from the estimates (6.74) (6.75) and from standard continuity, interpolation and compactness arguments.
$``$$``$
## 7 Asymptotics and wave operators for $`u`$
In this section we complete the construction of the wave operators for the equation (1.1) and we derive asymptotic properties of solutions in their range. The construction relies in an essential way on those of Section 6, especially Proposition 6.5, and will require a discussion of the gauge invariance of those constructions. This section follows closely Section II.7.
We first define the wave operator for the auxiliary system (2.11)-(2.12).
Definition 7.1. We define the wave operator $`\mathrm{\Omega }_0`$ as the map
$$\mathrm{\Omega }_0:(w_+,\psi _+)(w,\phi )$$
$`(7.1)`$
from $`K_\rho _{_{\mathrm{}}}^{k_+}Y_\rho _{_{\mathrm{}}}^\mathrm{}_+`$ to the space of $`(w,\phi )`$ such that $`(w,\overline{h}_0^1\phi )𝒳_\rho ^{k,\mathrm{}}([T,\mathrm{}))`$ for some $`T`$, $`1T<\mathrm{}`$, where $`\rho `$, $`k_+`$, $`\mathrm{}_+`$, $`k`$, $`\mathrm{}`$ satisfy the assumptions of Proposition 6.5 and $`(w,\phi )`$ is the solution of the system (2.11)-(2.12) obtained in Part (1) of that proposition.
In order to study the gauge invariance of $`\mathrm{\Omega }_0`$, we need some information on the Cauchy problem at finite times for the equation (1.1). We define the operator $`JJ(t)=x+it`$, which satisfies the commutation relation
$$iMD=JMD,$$
$`(7.2)`$
where $`M`$ and $`D`$ are defined by (2.4) (2.5). For any interval $`I[1,\mathrm{})`$, any nonnegative integer $`k`$ and any nonnegative $`𝒞^1`$ function $`\rho `$ defined in $`I`$, we define the space
$$\begin{array}{cc}X_\rho ^k(I)\hfill & =\{u:D^{}M^{}u𝒞(I,K_\rho ^k)\}\hfill \\ & \\ & =\{u:<J(t)>^kf(J(t))u𝒞(I,L^2)\}.\hfill \end{array}$$
$`(7.3)`$
where $`<\lambda >=(1+\lambda ^2)^{1/2}`$ for any real number or self-adjoint operator $`\lambda `$ and where the second equality follows from (7.2) and from Remark 3.1. (The space $`X_0^k(I)`$ was denoted $`𝒳^k(I)`$ in II). We recall the following result (see Proposition I.7.1).
Proposition 7.1. Let $`k`$ be a positive integer and let $`0<\mu <2k`$. Then the Cauchy problem for the equation (1.1) with initial data $`u(t_0)=u_0`$ such that $`<J(t_0)>^k`$ $`u_0L^2`$ at some initial time $`t_01`$ is locally well posed in $`𝒳_0^k()`$, namely
(1) There exists $`T>0`$ such that (1.1) has a unique solution with initial data $`u(t_0)=u_0`$ in $`𝒳_0^k([1(t_0T),t_0+T])`$.
(2) For any interval $`I`$, $`t_0I[1,\mathrm{})`$, (1.1) with initial data $`u(t_0)=u_0`$ has at most one solution in $`𝒳_0^k(I)`$.
(3) The solution of Part (1) depends continuously on $`u_0`$ in the norms considered there.
We come back from the system (2.11)-(2.12) to the equation (1.1) by reconstructing $`u`$ from $`(w,\phi )`$ by (2.7) and accordingly we define the map
$$\mathrm{\Lambda }:(w,\phi )u=MD\mathrm{exp}(i\phi )w.$$
$`(7.4)`$
It follows immediately from Lemma 3.3 that the map $`\mathrm{\Lambda }`$ satisfies the following property.
Lemma 7.1. Let $`\mathrm{}+2>n/2`$ and $`0k\mathrm{}+2`$. Then for any interval $`I[1,\mathrm{})`$ and any nonnegative $`𝒞^1`$ function $`\rho `$ defined in $`I`$, the map $`\mathrm{\Lambda }`$ is bounded and continuous from $`𝒴_{\rho ,loc}^{k,\mathrm{}}(I)`$ (defined by (5.1) (6.2)) to $`X_\rho ^k(I)`$, with norm estimates in compact intervals independent of $`\rho `$.
We now give the following definition
Definition 7.2. Two solutions $`(w,\phi )`$ and $`(w^{},\phi ^{})`$ of the system (2.11)-(2.12) in $`𝒴_{\rho ,loc}^{k,\mathrm{}}(I)`$ for some $`k`$, $`\mathrm{}`$, $`\rho `$ and some interval $`I[1,\mathrm{})`$ are said to be gauge equivalent if $`\mathrm{\Lambda }(w,\phi )=\mathrm{\Lambda }(w^{},\phi ^{})`$, or equivalently if
$$\mathrm{exp}(i\phi (t))w(t)=\mathrm{exp}(i\phi ^{}(t))w^{}(t)$$
$`(7.5)`$
for all $`tI`$.
The following sufficient condition for gauge equivalence follows immediately from Lemma 7.1 and from the uniqueness statement of Proposition 7.1, part (2).
Lemma 7.2. Let $`\mathrm{}+2>n/2`$ and $`0k\mathrm{}+2`$. Let $`I[1,\mathrm{})`$ be an interval, let $`\rho `$ be a strictly positive $`𝒞^1`$ function defined in $`I`$, and let $`(w,\phi )`$ and $`(w^{},\phi ^{})`$ be two solutions of the system (2.11)-(2.12) in $`𝒴_{\rho ,loc}^{k,\mathrm{}}(I)`$. In order that $`(w,\phi )`$ and $`(w^{},\phi ^{})`$ be gauge equivalent, it is sufficient that (7.5) holds for one $`tI`$.
We now turn to the study of the gauge equivalence of asymptotic states (if any) for solutions of the system (2.11)-(2.12) such as those obtained in Proposition 4.1. For that purpose, we need $`\rho `$ to be decreasing, namely to be defined by (4.1) for some $`t_01`$.
Proposition 7.2. Let $`k0`$ and $`\mathrm{}>n/2\nu `$. Let $`\rho `$ be defined by (4.1) for some $`t_01`$ and let $`h_0`$ and $`h_1`$ be as in Proposition 4.1. Let $`(w,\phi )`$ and $`(w^{},\phi ^{})`$ be two gauge equivalent solutions of the system (2.11)-(2.12) such that $`(w,h_0^1\phi )`$, $`(w^{},h_0^1\phi ^{})𝒳_\rho ^{k,\mathrm{}}([T,\mathrm{}))`$ for some $`T`$, $`t_0T<\mathrm{}`$. Let
$$b=Z(\phi ;[T,\mathrm{}),h_0,\mathrm{})Z(\phi ^{};[T,\mathrm{}),h_0,\mathrm{}).$$
$`(7.6)`$
Then
(1) There exists $`\omega Y_{\rho (\mathrm{})}^{\mathrm{}1+\nu }`$ such that $`\phi ^{}(t)\phi (t)`$ converges to $`\omega `$ strongly in $`Y_{\rho (\mathrm{})}^{\mathrm{}^{}}`$ for $`\mathrm{}^{}<\mathrm{}1+\nu `$ and weakly in $`Y_{\rho (\mathrm{})}^{\mathrm{}1+\nu }`$. The following estimate holds :
$$\left|\phi ^{}(t)\phi (t)\omega \right|_{\mathrm{}2+\nu }Cb\left|\phi ^{}(T_1)\phi (T_1)\right|_{\mathrm{}1+\nu }h_1(t)$$
$`(7.7)`$
for $`tT_1`$, for some $`T_1`$ sufficiently large, namely satisfying
$$bh_1(T_1)c$$
$`(7.8)`$
for some constant $`c`$.
(2) Assume in addition that $`1\nu /2k\mathrm{}+1`$ and let $`w_+`$ and $`w_+^{}`$ be the limits of $`w(t)`$ and $`w^{}(t)`$ as $`t\mathrm{}`$ obtained in Proposition 4.2. Then $`w_+^{}=w_+\mathrm{exp}(i\omega )`$.
(3) Let $`p0`$ be an integer. Assume in addition that $`w_+`$, $`w_+^{}K_{\rho (\mathrm{})}^{k_+}`$ where $`k_+`$ satisfies
$$k_+>n/2,k_+(p+1)\overline{\lambda }1$$
$`(7.9)`$
and let $`\varphi _p`$, $`\varphi _p^{}`$ be associated with $`(w_+,w_+^{})`$ according to (5.22) and Proposition 5.2. Assume that the following limits exist
$$\underset{t\mathrm{}}{lim}(\phi (t)\varphi _p(t))=\psi _+,\underset{t\mathrm{}}{lim}(\phi ^{}(t)\varphi _p^{}(t))=\psi _+^{}$$
$`(7.10)`$
as strong limits in $`L^{\mathrm{}}`$. Then $`\psi _+^{}=\psi _++\omega `$.
Proof. Part (1). Let $`\phi _\pm (t)=\phi ^{}(t)\pm \phi (t)`$. By the same estimates as in the proof of Lemma 3.7, we obtain
$$_t|\phi _{}|_{\mathrm{}^{}}^22\rho ^{}|\phi _{}|_{\mathrm{}^{}+\nu /2}^2Ct^2|\phi _{}|_{\mathrm{}^{}+\nu /2}\left\{|\phi _{}|_{\mathrm{}^{}+\nu /2}|\phi _+|_{\mathrm{}}+|\phi _{}|_{\mathrm{}^{}}|\phi _+|_{\mathrm{}+\nu /2}\right\}$$
$`(7.11)`$
for $`\nu /2\mathrm{}^{}+1\mathrm{}+\nu `$. Defining $`z_{(1)}=z_{(1)}(\phi _{};[T_1,t],1,\mathrm{}^{})`$ with $`TT_1<t`$, integrating over time and using (7.6), we obtain
$$\begin{array}{cc}z^2z_1^2\hfill & z_0^2+C(\underset{[T_1,t]}{Sup}t^2|\rho ^{}|^1h_0)bz_1(z+z_1)\hfill \\ & \\ & z_0^2+Cbh_1(T_1)z_1(z+z_1)\hfill \end{array}$$
where $`z_0=|\phi _{}(T_1)|_{\mathrm{}^{}}`$, which under the condition (7.8) with suitable $`c`$, yields
$$Z(\phi _{};[T_1,\mathrm{}),1,\mathrm{}^{})C|\phi _{}(T_1)|_{\mathrm{}^{}}.$$
$`(7.12)`$
From
$$_t\phi _{}=\left(2t^2\right)^1\left(\phi _{}\phi _+\right)$$
we next estimate directly
$$\begin{array}{cc}|_t\phi _{}|_\mathrm{}^{}1\hfill & Ct^2|\phi _{}|_{\mathrm{}^{}}|\phi _+|_{\mathrm{}}\hfill \\ & \\ & Ct^2h_0b|\phi _{}(T_1)|_{\mathrm{}^{}}\hfill \end{array}$$
$`(7.13)`$
by (7.6) and (7.12). The last member of (7.13) is integrable in time since
$$_t^{\mathrm{}}𝑑t_1t_1^2h_0(t_1)\rho ^{}_1h_1(t),$$
$`(7.14)`$
which for $`\mathrm{}^{}=\mathrm{}1+\nu `$ proves the existence of the limit $`\omega Y_{\rho (\mathrm{})}^{\mathrm{}2+\nu }`$ for $`\phi _{}`$ together with the estimate (7.7). The fact that actually $`\omega Y_{\rho (\mathrm{})}^{\mathrm{}1+\nu }`$ and the other convergences stated in Part (1) follow therefrom and from (7.12) by standard compactness and interpolation arguments.
Parts (2) and (3). The proof is identical with that of the corresponding statements in Proposition II.7.2 and will be omitted.
$``$$``$
Remark 7.1. The assumptions made on $`(k,\mathrm{})`$ in Parts (1) and (2) of Proposition 7.2 are implied by the assumptions (3.48) and $`k1\nu /2`$ of Proposition 4.1, so that Parts (1) and (2) apply directly to the solutions of the system (2.11)-(2.12) constructed in that proposition. In Part (3), the assumption (7.9) is that required in Proposition 5.2 to ensure the existence and appropriate estimates of $`\varphi _p`$, $`\varphi _p^{}`$. That assumption, together with the existence of the limits (7.10), is ensured under the assumptions of Proposition 5.3, namely if $`(k,\mathrm{})`$ satisfy in addition (5.24) for some $`k_0`$ satisfying (5.25) and if $`(p+2)\gamma >1`$ and $`P_p(1)<\mathrm{}`$.
Proposition 7.2 prompts us to make the following definition.
Definition 7.3. Two pairs of asymptotic states $`(w_+,\psi _+)`$ and $`(w_+^{},\psi _+^{})`$ are gauge equivalent if $`w_+\mathrm{exp}(i\psi _+)=w_+^{}\mathrm{exp}(i\psi _+^{})`$.
With this definition, Proposition 7.2 implies that two gauge equivalent solutions of the system (2.11)-(2.12) in $`(\mathrm{\Omega }_0)`$ are images of two gauge equivalent pairs of asymptotic states. One should however not overlook the following fact. The wave operator $`\mathrm{\Omega }_0`$ is defined through Proposition 6.5 which uses an increasing $`\rho `$ defined by (6.1) whereas Proposition 7.2 uses a decreasing $`\rho `$ (essentially in the proof of (7.12)). In order to apply Proposition 7.2 to the solutions constructed in Proposition 6.5 with increasing $`\rho `$, one has therefore to take some large $`t_0`$, to define
$$\stackrel{~}{\rho }(t)=\rho (t_0)_{t_0}^t\rho ^{}(t_1)𝑑t_1$$
and to apply Proposition 7.2 with that new $`\stackrel{~}{\rho }`$, thereby ending up with information on $`(\omega ,w_+,w_+^{},`$ $`\psi _+,\psi _+^{})`$ in spaces $`K`$ or $`Y`$ associated with $`\stackrel{~}{\rho }(\mathrm{})=\rho (\mathrm{})2_{t_0}^{\mathrm{}}\rho ^{}(t)𝑑t<\rho (\mathrm{})`$. That fact of course does not impair the algebraic relations expressing gauge invariance.
We now turn to the converse result, namely to the fact that gauge equivalent asymptotic states generate gauge equivalent solutions through $`\mathrm{\Omega }_0`$.
Proposition 7.3. Let $`(k,\mathrm{})`$ and $`(k_+,\mathrm{}_+)`$ satisfy the assumptions of Proposition 6.5, namely (3.48), $`k1\nu /2`$ and (6.38). Let $`\rho `$ and the estimating functions of time also satisfy the assumptions of Proposition 6.5. Let $`(w_+,\psi _+)`$, $`(w_+^{},\psi _+^{})K_\rho _{_{\mathrm{}}}^{k_+}Y_\rho _{_{\mathrm{}}}^\mathrm{}_+`$ be gauge equivalent and let $`(w,\phi )`$, $`(w^{},\phi ^{})`$ be their images under $`\mathrm{\Omega }_0`$. Then $`(w,\phi )`$ and $`(w^{},\phi ^{})`$ are gauge equivalent.
Proof. The proof is identical with that of Proposition II.7.3. We reproduce it for completeness.
Let $`t_0`$ be sufficiently large and let $`(w_{t_0},\phi _{t_0})`$ and $`(w_{t_0}^{},\phi _{t_0}^{})`$ be the solutions of the system (2.11)-(2.12) constructed by Proposition 6.4. From the initial conditions
$$w_{t_0}(t_0)=V(t_0),w_{t_0}^{}(t_0)=V^{}(t_0),$$
$$\phi _{t_0}(t_0)=\varphi _p(t_0)+\chi (t_0),\phi _{t_0}^{}(t_0)=\varphi _p^{}(t_0)+\chi ^{}(t_0),$$
from the fact that $`\varphi _p=\varphi _p^{}`$ by Proposition 5.2, part (2) and that $`V\mathrm{exp}(i\chi )=V^{}\mathrm{exp}(i\chi ^{})`$ by Proposition 6.2, part (4), it follows that
$$w_{t_0}(t_0)\mathrm{exp}(i\phi _{t_0}(t_0))=w_{t_0}^{}(t_0)\mathrm{exp}(i\phi _{t_0}^{}(t_0))$$
and therefore by Lemma 7.2, $`(w_{t_0},\phi _{t_0})`$ and $`(w_{t_0}^{},\phi _{t_0}^{})`$ are gauge equivalent, namely
$$w_{t_0}(t)\mathrm{exp}(i\phi _{t_0}(t))=w_{t_0}^{}(t)\mathrm{exp}(i\phi _{t_0}^{}(t))$$
$`(7.15)`$
for all $`t`$ for which both solutions are defined.
We now take the limit $`t_0\mathrm{}`$ for fixed $`t`$ in (7.15). By Proposition 6.5, part (2), for fixed $`t`$, $`(w_{t_0},\phi _{t_0})`$ and $`(w_{t_0}^{},\phi _{t_0}^{})`$ converge respectively to $`(w,\phi )`$ and $`(w^{},\phi ^{})`$ in $`K_\rho ^k^{}Y_\rho ^{\mathrm{}^{}}`$. By Lemma 3.3, one can take the limit $`t_0\mathrm{}`$ in (7.15), thereby obtaining (7.5), so that $`(w,\phi )`$ and $`(w^{},\phi ^{})`$ are gauge equivalent.
$``$$``$
We can now define the (local) wave operators (at infinity) for $`u`$.
Definition 7.4. The wave operator $`\mathrm{\Omega }`$ is defined as the map
$$\mathrm{\Omega }:u_+u=(\mathrm{\Lambda }\mathrm{\Omega }_0)(Fu_+,0)$$
$`(7.16)`$
where $`\mathrm{\Omega }_0`$ and $`\mathrm{\Lambda }`$ are defined by Definition 7.1 and by (7.4).
We collect in the following proposition the information on $`\mathrm{\Omega }`$ that follows from the previous study, in particular from Propositions 6.5 and 7.3.
Proposition 7.4. Let $`p0`$ be an integer with $`(p+2)\gamma >1`$. Let $`\rho `$ (defined by (6.1)) and the estimating functions of time satisfy the assumptions of Proposition 6.5 (see especially (6.39) (6.40) and Remark 6.4). Let
$$\lambda =\mu n+22\nu ,$$
$`(7.17)`$
$$k1\nu /2,k>1\nu +\mu /2.$$
$`(7.18)`$
Let $`k_+`$ satisfy (6.38) for some $`\mathrm{}`$ satisfying $`\mathrm{}>n/2\nu `$, $`\mathrm{}k\nu `$. Then
(1) The wave operator $`\mathrm{\Omega }`$ maps $`FK_\rho _{_{\mathrm{}}}^{k_+}`$ to $`X_\rho ^k([T,\mathrm{}))`$ for some $`T`$, $`1T<\mathrm{}`$. (Actually $`T`$ depends on $`u_+`$).
(2) $`\mathrm{\Omega }`$ is injective.
Proof. Part (1) follows from the definitions, from Lemma 7.1 and from Proposition 6.5 with $`\psi _+=0`$. The only point to be checked is the fact that (7.17) (7.18) imply the existence of $`\mathrm{}`$ such that $`(k,\mathrm{})`$ satisfies (3.48). Now the $`\mathrm{}`$ dependent part of (3.48) reduces to
$$k\nu \mathrm{}k+\nu \lambda $$
$`(7.19)`$
$$n/2\nu <\mathrm{}<2k\lambda +\nu n/2$$
$`(7.20)`$
and the compatibility of (7.19) (7.20) for $`\mathrm{}`$ is easily seen to reduce to (7.17) and to the second inequality in (7.18).
Part (2) follows from Proposition 7.2 and from the fact that a gauge equivalence class of asymptotic states contains at most one element with $`\psi _+=0`$.
$``$$``$
Remark 7.2. One may wonder whether the restriction to asymptotic states with $`\psi _+=0`$ in (7.16) restricts the range of $`\mathrm{\Omega }`$ as compared with that of $`\mathrm{\Lambda }\mathrm{\Omega }_0`$. From Proposition 7.3, it follows that
$$(\mathrm{\Lambda }\mathrm{\Omega }_0)(w_+,\psi _+)=(\mathrm{\Lambda }\mathrm{\Omega }_0)(w_+\mathrm{exp}(i\psi _+),0)$$
in so far as $`w_+\mathrm{exp}(i\psi _+)`$ has the regularity needed to apply Proposition 6.5. This is the case if $`\mathrm{}_+k_+2`$ by Lemma 3.3. That condition however is significantly stronger than the condition on $`\mathrm{}_+`$ contained in (6.38), especially for large $`p`$, i.e. for small $`\gamma `$. Therefore, there is actually a restriction of the range for regularity reasons, in spite of the (algebraic) gauge invariance of the construction expressed by Proposition 7.3.
We next collect the information and in particular the asymptotic estimates obtained for the solutions of the equation (1.1) in $`(\mathrm{\Omega })`$.
Proposition 7.5. Let $`0<\mu n2+2\nu n`$. Let $`0<\gamma 1`$ and let $`p0`$ be an integer with $`(p+2)\gamma >1`$. Let $`\rho `$ (defined by 6.1) and the estimating functions of time satisfy the assumptions of Proposition 6.5 (especially (6.39) (6.40)). Let $`(k,\mathrm{},k_+)`$ satisfy (7.18) (3.48) (6.38). Let $`u_+FK_\rho _{_{\mathrm{}}}^{k_+}`$ and $`a_+=Fu_+;K_\rho _{_{\mathrm{}}}^{k_+}`$. Then
(1) There exists $`T`$, $`1T<\mathrm{}`$, and there exists a unique solution $`uX_\rho ^k([T,\mathrm{}))`$ of the equation (1.1) which can be represented as
$$u=MD\mathrm{exp}(i\phi )w$$
where $`(w,\phi )`$ is a solution of the system (2.11)-(2.12) such that $`(w,\overline{h}_0^1\phi )𝒳_\rho ^{k,\mathrm{}}([T,\mathrm{}))`$ and such that
$$|w(t)Fu_+|_{k1+\nu }h_0(t)0$$
$`(7.21)`$
$$|\phi (t)\varphi _p(t)|_{\mathrm{}1+\nu }0$$
$`(7.22)`$
when $`t\mathrm{}`$. The time $`T`$ can be defined by a condition of the type (6.41) with $`b_+=0`$.
(2) The solution is obtained as $`u=\mathrm{\Omega }u_+`$, following Definition 7.4.
(3) The map $`\mathrm{\Omega }`$ is continuous from $`FK_\rho _{_{\mathrm{}}}^{k_+}`$ to the norm topology of $`X_\rho ^k^{}(I)`$ for $`k^{}<k`$ and to the weak-$``$ topology of $`X_\rho ^k(I)`$ for any compact interval $`I[T,\mathrm{})`$, and to the weak topology of $`MDK_\rho ^k`$ pointwise in $`t`$.
(4) The solution $`u`$ satisfies the following estimate for $`tT`$ :
$$<J(t)>^kf(J(t))\left(\mathrm{exp}[i\varphi _p(t,x/t)]u(t)M(t)D(t)Fu_+\right)_2A(a_+)h_3(t)$$
$`(7.23)`$
for some estimating function $`A(a_+)`$.
(5) Let $`2r<\mathrm{}`$. The solution $`u`$ satisfies the following estimate :
$$u(t)\mathrm{exp}[i\varphi _p(t,x/t)]M(t)D(t)Fu_+_rA(a_+)\rho (t)^\beta t^{\delta (r)}h_3(t)$$
$`(7.24)`$
where $`\delta (r)=n/2n/r`$, $`\beta =\nu ^1(\delta (r)k)0`$ if $`r<\mathrm{}`$ or $`k>n/2`$, $`\beta =\nu ^1(n/2k+\epsilon )`$ if $`r=\mathrm{}`$ and $`kn/2`$.
Proof. Parts (1) (2) (3) follow from Propositions 6.5 and 4.3, from Definition 7.4 and from Proposition 7.4.
Part (4). From the commutation relation (7.2), from Remark 3.1 and from Lemma 3.3, it follows that the LHS of (7.23) is estimated by
$$_2C|\mathrm{exp}(i(\varphi _p\phi ))wFu_+|_kC\{|wFu_+|_k+|\mathrm{exp}(i(\varphi _p\phi ))1|_{\mathrm{}}|w|_k\}$$
$$C\left\{\left|wW_p\right|_k+\left|W_pw_+\right|_k+\mathrm{exp}\left(C\left|\varphi _p\phi \right|_{\mathrm{}}\right)\left|\varphi _p\phi \right|_{\mathrm{}}|w|_k\right\}$$
and the result follows from the estimates (6.72) (5.9) (6.73) (6.74).
Part (5) follows from Part (4) and from the inequality
$$\begin{array}{cc}v_r\hfill & =t^{\delta (r)}D^{}M^{}v_rCt^{\delta (r)}<>^mD^{}M^{}v_2\hfill \\ & \\ & =Ct^{\delta (r)}<J(t)>^mv_2\hfill \end{array}$$
which follows from (7.2) and Sobolev inequalities with $`m=\delta (r)`$ if $`r<\mathrm{}`$, $`m=n/2+\epsilon `$ if $`r=\mathrm{}`$, together with the fact that
$$|\xi |^{mk}f(\xi )^1C\rho ^\beta .$$
$``$$``$
Remark 7.3. In (7.23) and (7.24), one could replace $`MDFu_+`$ by $`U(t)u_+`$ since the difference is small in the relevant norms. One could also replace $`Fu_+`$ by $`W_p`$, but this would not produce any improvement in the estimates, since the main contribution comes from the phase. The estimate (7.24) is a rather weak one, since we have omitted the function $`f`$ which generates the Gevrey regularity. It is only one example of a large number of similar estimates exhibiting the typical $`t^{\delta (r)}`$ decay associated with $`L^r`$ norms.
The final step of the standard construction of the wave operators for the equation (1.1) would consist in extending the solutions $`u`$ to arbitrary finite times, and defining the maps $`\mathrm{\Omega }_1:u_+u(1)`$ where $`u=\mathrm{\Omega }u_+`$. In order not to waste the Gevrey regularity of the local solutions at infinity, this would require a treatment of the global Cauchy problem at finite times for arbitrarily large data in the same Gevrey framework. This is a rather different problem and we shall refrain from considering it here.
Acknowledgements. Part of this work was done while one of the authors (G.V.) was visiting the Institut des Hautes Etudes Scientifiques (IHES), Bures-sur-Yvette, France and the Laboratoire de Physique Théorique (LPT), Université de Paris-Sud, France. He is very grateful to Professor Jean-Pierre Bourguignon, Director of the IHES, and to Professor Dominique Schiff, Director of the LPT, for the kind hospitality extended to him.
## Appendix A
In this appendix we derive a number of properties of the function $`\stackrel{~}{f}`$ defined by (3.2), which make it a possible substitute for the function $`f_0`$ defined by (3.1) in the definition of the spaces $`𝒳_\rho ^{k,\mathrm{}}`$. In all this appendix, we assume $`0<\nu 1`$ and we take $`\rho =1`$. The parameter $`\rho `$ can be reintroduced easily by scaling. Accordingly we define
$$\stackrel{~}{f}(\xi )=\underset{j0}{}(j!)^{1/\nu }|\xi |^j,$$
$`(\mathrm{A}.1)`$
$$F(\xi )=\underset{j0}{}(j+1)^1(j!)^{1/\nu }|\xi |^{j+1}.$$
$`(\mathrm{A}.2)`$
We also use (A.1) and (A.2) to define $`\stackrel{~}{f}`$ and $`F`$ when applied to $`\xi IR^+`$. With that convention, we have $`\stackrel{~}{f}(\xi )=\stackrel{~}{f}(|\xi |)`$ and $`F(\xi )=F(|\xi |)`$ for all $`\xi IR^n`$. Furthermore $`\stackrel{~}{f}=dF/d|\xi |`$.
We first derive some preliminary estimates which allow in particular for a comparison of $`\stackrel{~}{f}`$ and $`f_0`$.
Lemma A.1. The following estimates hold for all $`\xi IR^n`$ :
$$\underset{j1}{}j(j!)^{1/\nu }|\xi |^j|\xi |^\nu \underset{j0}{}(j!)^{1/\nu }|\xi |^j\underset{j0}{}(j+1)(j!)^{1/\nu }|\xi |^j,$$
$`(\mathrm{A}.3)`$
$$\underset{j0}{}(j+1)^1(j!)^{1/\nu }|\xi |^{j+1}|\xi |^{1\nu }\underset{j0}{}(j!)^{1/\nu }|\xi |^j,$$
$`(\mathrm{A}.4)`$
$$\stackrel{~}{f}^1(d\stackrel{~}{f}/d|\xi |)|\xi |^{\nu 1}F^1\stackrel{~}{f}=F^1(dF/d|\xi |),$$
$`(\mathrm{A}.5)`$
$$|\xi |^\nu \stackrel{~}{f}d(|\xi |\stackrel{~}{f})/d|\xi |,$$
$`(\mathrm{A}.6)`$
$$F(a)(|\xi |a)^{\nu 1}\mathrm{exp}\left(\nu ^1(|\xi |^\nu a^\nu )\right)\stackrel{~}{f}(\xi )\mathrm{exp}(\nu ^1|\xi |^\nu )$$
$`(\mathrm{A}.7)`$
for all $`a>0`$,
$$\stackrel{~}{f}(\xi )=(2\pi )^{(\nu 1)/2\nu }\nu ^{1/2}|\xi |^{(\nu 1)/2}\mathrm{exp}\left(\nu ^1|\xi |^\nu \right)(1+o(1))\mathrm{𝑤ℎ𝑒𝑛}|\xi |\mathrm{}.$$
$`(\mathrm{A}.8)`$
Proof. (A.3a) and (A.4) follow from the Hölder inequality on $`\mathrm{Z}\mathrm{}^+`$ with the measure $`(j!)^{1/\nu }|\xi |^j`$ and the exponents $`1/\nu `$ and $`1/(1\nu )`$, applied respectively to the pairs of functions $`(j,1)`$ and $`((j+1)^1,1)`$.
(A.3b) follows similarly from the Hölder inequality with the measure $`(j+1)(j!)^{1/\nu }|\xi |^j`$ and the exponents $`\nu ^1(1+\nu )`$ and $`1+\nu `$, applied to the functions $`(j+1)^1`$ and 1.
(A.5) is a rewriting of (A.3a) and (A.4), while (A.6) is a rewriting of (A.3b).
(A.7) follows from (A.5). In fact (A.5) states that the functions $`\stackrel{~}{f}(\xi )\mathrm{exp}(|\xi |^\nu /\nu )`$ and $`F(\xi )\mathrm{exp}(|\xi |^\nu /\nu )`$ are respectively decreasing (and therefore less than one) and increasing in $`|\xi |`$. The first fact yields (A.7b) while both of them together with $`\stackrel{~}{f}|\xi |^{\nu 1}F`$ yield (A.7a). Note also that (A.7b) follows directly from the definition (A.1) and from the fact that
$$\stackrel{~}{f}(\xi )=(j!)^1|\xi |^{\nu j};\mathrm{}^{1/\nu }^{1/\nu }(j!)^1|\xi |^{\nu j};\mathrm{}^1^{1/\nu }=\mathrm{exp}(\nu ^1|\xi |^\nu )$$
by the standard embedding of $`\mathrm{}^p`$ spaces.
(A.8) is proved in (see (8.07) p. 308) in the special case $`\nu 1/4`$, but the proof extends easily to the whole range $`0<\nu 1`$.
$``$$``$
The estimates (A.7) and the asymptotic property (A.8) compare $`\stackrel{~}{f}`$ with $`f_0`$ by stating that essentially $`\stackrel{~}{f}f_0^{1/\nu }`$. On the other hand (A.6) is the analogue of the fact that $`df_0/d|\xi |=\nu |\xi |^{\nu 1}f_0`$ and allows for the construction of function space norms satisfying an inequality which can replace (3.12) in the subsequent estimates.
We next show that $`\stackrel{~}{f}`$ satisfies estimates similar to those of Lemma 3.1.
Lemma A.2. Let $`(\xi ,\eta )IR^n`$. Then $`\stackrel{~}{f}`$ satisfies the estimates
$$\stackrel{~}{f}(\xi )\stackrel{~}{f}(\xi \eta )\stackrel{~}{f}(\eta )\mathrm{𝑓𝑜𝑟}\mathrm{𝑎𝑙𝑙}(\xi ,\eta ),$$
$`(\mathrm{A}.9)`$
$$\stackrel{~}{f}(\xi )\stackrel{~}{f}(\xi \eta )\mathrm{exp}(|\eta |^\nu )\mathrm{𝑓𝑜𝑟}|\xi ||\eta ||\xi \eta |,$$
$`(\mathrm{A}.10)`$
$$(\stackrel{~}{f}(\eta )\stackrel{~}{f}(\xi ))|\eta |^{1\nu }|\xi \eta |\stackrel{~}{f}(\eta )\mathrm{𝑓𝑜𝑟}|\xi ||\eta |,$$
$`(\mathrm{A}.11)`$
$$|\stackrel{~}{f}(\xi )\stackrel{~}{f}(\eta )||\eta |^{1\nu }|\xi \eta |^{1\nu }\stackrel{~}{f}(\xi \eta )\stackrel{~}{f}(\eta )\mathrm{𝑓𝑜𝑟}\mathrm{𝑎𝑙𝑙}(\xi ,\eta ),$$
$`(\mathrm{A}.12)`$
$$|\stackrel{~}{f}(\xi )\stackrel{~}{f}(\eta )||\eta |^{1\nu }|\xi \eta |^{1\nu }(\mathrm{exp}(|\xi \eta |^\nu )1)\stackrel{~}{f}(\eta )\mathrm{𝑓𝑜𝑟}|\xi ||\xi \eta ||\eta |,$$
$`(\mathrm{A}.13)`$
$$|\stackrel{~}{f}(\xi )\stackrel{~}{f}(\eta )||\eta |^{1\nu }C|\xi \eta |^{1\nu }(\stackrel{~}{f}(\xi \eta )1)\mathrm{exp}(|\eta |^\nu )\mathrm{𝑓𝑜𝑟}|\xi ||\eta ||\xi \eta |.$$
$`(\mathrm{A}.14)`$
In (A.14) one can take $`C=1`$, except in the region $`|\xi ||\xi \eta ||\eta |`$ where $`C=2^{1\nu }`$.
The function $`\stackrel{~}{f}(\xi )\stackrel{~}{f}(1)`$ satisfies the same estimates as $`\stackrel{~}{f}(\xi )`$.
Proof. (A.9) is trivial except if $`|\xi ||\eta ||\xi \eta |`$. In all cases, we estimate
$$\begin{array}{cc}\stackrel{~}{f}(\xi )\hfill & \underset{j,k0}{}(j!)^1(k!)^1((j+k)!)^{11/\nu }|\xi \eta |^j|\eta |^k\hfill \\ & \\ & \underset{j,k0}{}(j!)^{1/\nu }(k!)^{1/\nu }|\xi \eta |^j|\eta |^k=\stackrel{~}{f}(\xi \eta )\stackrel{~}{f}(\eta )\hfill \end{array}$$
$`(\mathrm{A}.15)`$
since $`(j+k)!j!k!`$.
(A.10) is trivial in the allowed region except if $`|\eta ||\xi \eta ||\xi |`$. In that case, we rewrite the first inequality in (A.15) as
$$\begin{array}{cc}\stackrel{~}{f}(\xi )\hfill & \underset{k0}{}\left\{\underset{j0}{}\left((j!)^{1/\nu }|\xi \eta |^j\right)^\nu \left(((j+k)!)^{1/\nu }|\xi \eta |^{j+k}\right)^{1\nu }\right\}(k!)^1|\eta |^k|\xi \eta |^{(\nu 1)k}\hfill \\ & \\ & \stackrel{~}{f}(\xi \eta )\mathrm{exp}\left(|\eta ||\xi \eta |^{\nu 1}\right)\stackrel{~}{f}(\xi \eta )\mathrm{exp}(|\eta |^\nu )\hfill \end{array}$$
$`(\mathrm{A}.16)`$
by the Hölder inequality applied to the sum over $`j`$ for fixed $`k`$ and the fact that $`|\eta ||\xi \eta |`$.
(A.11). We estimate
$$\stackrel{~}{f}(\eta )\stackrel{~}{f}(\xi )=\underset{j1}{}(j!)^{1/\nu }(|\eta |^j|\xi |^j)|\xi \eta |\underset{j1}{}(j!)^{1/\nu }j|\eta |^{j1}|\xi \eta ||\eta |^{\nu 1}\stackrel{~}{f}(\eta )$$
by (A.3a).
(A.12) follows from (A.11) and from $`|\xi \eta |^\nu \stackrel{~}{f}(\xi \eta )`$ if $`|\xi ||\eta |`$. If $`|\xi ||\eta |`$, we estimate
$$\stackrel{~}{f}(\xi )\stackrel{~}{f}(\eta )\underset{j1,k0}{}(j!)^1(k!)^1((j+k)!)^{11/\nu }|\xi \eta |^j|\eta |^k$$
$$=\underset{j,k0}{}((j+1)!)^1(k!)^1((j+k+1)!)^{11/\nu }|\xi \eta |^{j+1}|\eta |^k$$
$$\underset{j,k0}{}(j+1)^1(j!)^{1/\nu }|\xi \eta |^{j+1}(k+1)((k+1)!)^{1/\nu }|\eta |^k$$
since $`(j+k+1)!j!(k+1)!`$,
$$\mathrm{}|\xi \eta |^{1\nu }\stackrel{~}{f}(\xi \eta )|\eta |^{\nu 1}\stackrel{~}{f}(\eta )$$
by (A.3a) and (A.4).
(A.13) follows from (A.11) and $`|\xi \eta |^\nu (\mathrm{exp}(|\xi \eta |^\nu )1)`$ if $`|\xi ||\eta |`$. If $`|\xi \eta ||\eta ||\xi |`$, we estimate in the same way as in (A.16)
$$(\stackrel{~}{f}(\xi )\stackrel{~}{f}(\eta ))|\eta |^{1\nu }|\eta |^{1\nu }\underset{j1}{}\left\{\underset{k0}{}((k!)^{1/\nu }|\eta |^k)^\nu \left(((j+k)!)^{1/\nu }|\eta |^{j+k}\right)^{1\nu }\right\}$$
$$\times (j!)^1|\xi \eta |^j|\eta |^{(\nu 1)j}$$
$$\stackrel{~}{f}(\eta )\left\{|\eta |^{1\nu }(\mathrm{exp}(|\xi \eta ||\eta |^{\nu 1})1)\right\}$$
$`(\mathrm{A}.17)`$
by the Hölder inequality applied to the sum over $`k`$ for fixed $`j`$. Now for fixed $`|\xi \eta |`$, the last bracket in (A.17) is a decreasing function of $`|\eta |`$ and for $`|\eta ||\xi \eta |`$ is therefore less than its value for $`|\eta |=|\xi \eta |`$, which yields (A.13) in that case.
(A.14). If $`|\xi ||\eta ||\xi \eta |`$, (A.14) with $`C=1`$ follows from $`|\eta ||\xi \eta |`$ and from
$$|\stackrel{~}{f}(\xi )\stackrel{~}{f}(\eta )|\stackrel{~}{f}(\xi \eta )1.$$
If $`|\eta ||\xi \eta ||\xi |`$, (A.14) with $`C=1`$ follows from $`|\eta ||\xi \eta |`$ and from a minor variant of (A.16) with the sum over $`j`$ restricted to $`j1`$, so that
$$\stackrel{~}{f}(\xi )\stackrel{~}{f}(\eta )(\stackrel{~}{f}(\xi \eta )1)\mathrm{exp}(|\eta |^\nu ).$$
$`(\mathrm{A}.18)`$
If $`|\xi ||\xi \eta ||\eta |`$, (A.14) with $`C=2^{1\nu }`$ follows from $`|\eta |2|\xi \eta |`$ and from (A.18) with $`\xi `$ and $`\eta `$ interchanged.
The last statement of Lemma A.2 is obvious as regards (A.9) (A.10) and follows from the fact that
$$|\stackrel{~}{f}(\xi )a\stackrel{~}{f}(\eta )a||\stackrel{~}{f}(\xi )\stackrel{~}{f}(\eta )|$$
for all $`\xi `$, $`\eta `$ and $`a>0`$ as regards (A.11)-(A.14), in the same way as in Lemma 3.1.
$``$$``$
It follows from Lemma A.2 that $`\stackrel{~}{f}(\xi )`$ and $`\stackrel{~}{f}(\xi )\stackrel{~}{f}(1)`$ satisfy the basic estimates (A.9) and (A.12) which are used throughout this paper, thereby making those functions into suitable substitutes for $`f_0`$ and $`f`$ in the definition of the spaces. Note also that by (A.7) (A.8), $`\mathrm{exp}(|\eta |^\nu )\stackrel{~}{f}(\eta )^\nu `$ (up to a small power of $`\eta `$), which makes (A.10) (A.13) (A.14) into close analogues of (3.4) (3.6) (3.7).
We finally use the function $`\stackrel{~}{f}`$ to relate the definition of spaces such as $`K_\rho ^k`$ or $`Y_\rho ^{\mathrm{}}`$ to more standard definition of Gevrey spaces. Since this part is meant to be only illustrative, we restrict our attention to space dimension $`n=1`$. The Gevrey class $`G_{1/\nu }`$ can be defined as the vector space of $`𝒞^{\mathrm{}}`$ functions $`u`$ such that there exists a constant $`C`$ such that
$$c_jC^j(j!)^{1/\nu }^ju\mathrm{}_j^p(L_x^q)$$
$`(\mathrm{A}.19)`$
where $`1p`$, $`q\mathrm{}`$ and $`\{c_j\}`$ is a sequence of positive numbers with at most polynomial increase or decrease at infinity. Of course if one allows for all possible $`C>0`$, the parameters $`p`$, $`q`$, and $`\{c_j\}`$ are irrelevant, and one can take $`p=q=\mathrm{}`$ and $`c_j1`$, which yields the standard definition. Here however we fix $`C`$, in fact $`C=\rho ^{1/\nu }`$, so as to obtain a Banach space, and those parameters become important. Since moreover we want a Hilbert space in order to apply the energy method in a convenient way, we take $`p=q=2`$, thereby obtaining the Gevrey Hilbert space $`X`$ with norm
$$u;X=b_j\rho ^{j/\nu }^ju;\mathrm{}_j^2(L_x^2)$$
$`(\mathrm{A}.20)`$
where $`b_j=c_j(j!)^{1/\nu }`$. Let now
$$f(\xi )=f(|\xi |)=\underset{j0}{}a_j|\xi |^j=f_+(\xi )+f_{}(\xi )$$
where $`\{a_j\}`$ is a sequence of positive numbers and where $`f_+`$ and $`f_{}`$ denote the sums over even and odd $`j`$ respectively. We claim that for a suitable relation between $`\{a_j\}`$ and $`\{b_j\}`$, the norm in $`X`$ is equivalent to the norm
$$u_{}=f(\rho ^{1/\nu }\xi )\widehat{u}_2.$$
$`(\mathrm{A}.21)`$
In fact
$$f_+(\xi )^2+f_{}(\xi )^2f(\xi )^22\left(f_+(\xi )^2+f_{}(\xi )^2\right)$$
so that if we define $`\{b_j\}`$ by
$$f_+(\xi )^2+f_{}(\xi )^2=\underset{j0}{}b_j^2|\xi |^{2j}$$
$`(A.22)`$
then
$$u;Xu_{}\sqrt{2}u;X$$
$`(A.23)`$
which proves the equivalence. The relation (A.22) can be rewritten as
$$b_k^2=\underset{0j2k}{}a_ja_{2kj}.$$
In the special case of $`\stackrel{~}{f}`$, where $`a_j=(j!)^{1/\nu }`$, it is obvious that
$$a_jb_j\sqrt{2j+1}a_j$$
and one can show that actually when $`j\mathrm{}`$
$$b_j=a_j(\pi \nu j)^{1/4}(1+o(1))$$
$`(A.24)`$
which together with (A.8) makes it possible to define equivalent norms of the form (A.20) for the spaces $`K_\rho ^k`$ and $`Y_\rho ^{\mathrm{}}`$.
## Appendix B
In this appendix we prove a Lemma which exemplifies the fact that for $`\rho >0`$ and $`\nu <1`$, the lower condition $`\mathrm{}+2>n/2`$ on $`\mathrm{}`$ needed in Lemma 3.3 in order to make $`Y_\rho ^{\mathrm{}}`$ into an algebra can be relaxed by using the estimate (3.4) instead of (3.3) (but not uniformly in $`\nu `$ and $`\rho `$). Following Remark 3.1, we consider the Hilbert space $`K`$ with norm
$$u;K=\overline{f}\widehat{u}_2$$
where $`\overline{f}`$ is either $`ff_1`$ or $`f_0f_1`$ with $`f_0`$, $`f`$ defined by (3.1) and $`f_1`$ defined by (3.10) for some $`k_<`$, $`k_>IR^+`$.
Lemma B.1. Let $`K`$ be as above, with $`0<\nu <1`$, $`\rho >0`$, $`k_>0`$ and $`0k_<<n/2`$. Then $`K`$ is an algebra, namely there exists a constant $`C`$ such that for all $`u_1`$, $`u_2K`$
$$u_1u_2;KCu_1;Ku_2;K.$$
$`(\mathrm{B}.1)`$
One can take
$$C^2=𝑑\eta \overline{f}(\eta )^2\left(1+2^{2k}f_0(\eta )^{2\nu }\right)$$
$`(\mathrm{B}.2)`$
where $`k=k_<k_>`$ and the integral converges under the assumptions made on $`\nu `$, $`\rho `$ and $`k_<`$.
Proof. By the Schwarz inequality, we estimate
$$\overline{f}\widehat{u_1u_2}_2^2=𝑑\xi \overline{f}(\xi )^2\left|𝑑\eta \widehat{u}_1(\eta )\widehat{u}_2(\xi \eta )\right|^2𝑑\xi \overline{f}(\xi )^2\left\{𝑑\eta \overline{f}(\eta )^2\overline{f}(\xi \eta )^2\right\}$$
$$\times \{d\eta \overline{f}(\eta )^2\overline{f}(\xi \eta )^2|\widehat{u}_1(\eta )|^2|\widehat{u}_2(\xi \eta )|^2\}C^2\overline{f}\widehat{u}_1_2^2\overline{f}\widehat{u}_2_2^2$$
where
$`C^2`$ $`=`$ $`\underset{\xi }{Sup}\overline{f}(\xi )^2{\displaystyle 𝑑\eta \overline{f}(\eta )^2\overline{f}(\xi \eta )^2}`$
$`=`$ $`2\underset{\xi }{Sup}\overline{f}(\xi )^2{\displaystyle _{|\eta ||\xi \eta |}}\overline{f}(\eta )^2\overline{f}(\xi \eta )^2.`$
Now for $`|\xi ||\xi \eta |`$, $`\overline{f}(\xi )\overline{f}(\xi \eta )`$ while for $`|\eta ||\xi \eta ||\xi |`$ one has $`|\xi |2|\xi \eta |`$ so that
$$f_1(\xi )2^{k_<k_>}f_1(\xi \eta )=2^kf_1(\xi \eta ),$$
and
$$f_{(0)}(\xi )f_{(0)}(\xi \eta )f_0(\eta )^\nu $$
by (3.4). Therefore
$$C^22\underset{\xi }{Sup}\left\{_{|\eta ||\xi ||\xi \eta |}𝑑\eta \overline{f}(\eta )^2+_{|\eta ||\xi \eta ||\xi |}𝑑\eta \overline{f}(\eta )^22^{2k}f_0(\eta )^{2\nu }\right\}.$$
The Supremum over $`\xi `$ is easily seen to be the limit $`|\xi |\mathrm{}`$, namely
$`C^2`$ $``$ $`2\left\{{\displaystyle _{\xi \eta 0}}𝑑\eta \overline{f}(\eta )^2+{\displaystyle _{\xi \eta 0}}𝑑\eta \overline{f}(\eta )^22^{2k}f_0(\eta )^{2\nu }\right\}`$
$`=`$ $`{\displaystyle 𝑑\eta \overline{f}(\eta )^2\left(1+2^{2k}f_0(\eta )^{2\nu }\right)}`$
and the integral converges for small $`\eta `$ by the condition $`k_<<n/2`$ and for large $`\eta `$ by the conditions $`\rho >0`$ and $`\nu <1`$.
$``$$``$
The same result with essentially the same proof holds if one replaces $`f_0`$ by $`\stackrel{~}{f}`$ in the definition of $`K`$. One then has to replace (3.4) by (A.10) in the proof. The same result also holds for arbitrary $`k_>IR`$ (still with $`0k_<<n/2)`$, but the proof is more cumbersome for $`k_><0`$.
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# Spin-orbit splitting in non-relativistic and relativistic self-consistent models
## I Introduction
One often asked question about the relativistic approach to nuclear structure is: what are the main physical effects that a relativistic approach can handle better than a non-relativistic approach? Among various arguments one feature seems to emerge quite obviously, namely the fact that, in a description based on the Dirac equation the nucleon spin degree of freedom is incorporated very naturally. Therefore, nuclear properties related to the spin and particularly spin-orbit effects in nuclei should come out from a relativistic description without any special adjustments whereas a non-relativistic approach would need ad hoc parameters. This was indeed so in the Dirac phenomenology analysis of spin observables measured in medium energy nucleon-nucleus scattering. The interplay of two deep phenomenological potentials, the attractive scalar field and the repulsive vector field gives the right magnitude for the central and spin-orbit optical potentials.
However, if one tries to calculate microscopically the mean field experienced by a nucleon bound to a nucleus and predict spin-orbit splittings, it is not so clear how successful will be the various relativistic models. In this work, we would like to analyze the spin-orbit splittings calculated in different self-consistent models: non-relativistic Hartree-Fock with Skyrme-type forces, relativistic mean field theory (RMFT) and relativistic Hartree-Fock (RHF). Of course, a mean field approach like Hartree-Fock is designed for reproducing global ground state properties like total binding energies or densities but it is not supposed to describe correctly single-particle spectra because core polarization effects are known to modify importantly single-particle energies. However, the core polarization effects should affect essentially in the same way the two members of a spin-orbit doublet if they are both below, or both above the Fermi level. Therefore, it is legitimate to study the question of spin-orbit splittings in Hartree-Fock or mean field frameworks for nuclei where both spin-orbit partners are occupied.
We carry out this study by looking at the evolution of the spin-orbit splittings of a $`nlj_>nlj_<`$ proton pair of states along an isotopic chain. The predictions of different models are compared with experimental values when they are available. The difficulties of such comparisons lie in the fact that the experimental information is scarce and it suffers sometimes of large uncertainties. This points to the necessity of having more data on spin-orbit splittings if one aims at more quantitative conclusions about the validity of different models concerning their spin-orbit properties. Within these limitations, our analysis will be done for the Ca, O and Sn isotopic chains.
The outline of the paper is as follows. In Section II and in the Appendix we derive the spin-orbit potentials corresponding to the different models: Skyrme-Hartree-Fock (SHF), RMFT and RHF. In Section III we discuss the results for the <sup>40</sup>Ca-<sup>48</sup>Ca case as well as those concerning the O and Sn isotopes. Concluding remarks are made in Section IV.
## II Formalism
Let us derive the main expressions for the spin-orbit potentials in non-relativistic and relativistic approaches. In either case the general method is to rewrite the Hartree or Hartree-Fock equations for the single-particle states in a Schrödinger-like form and to identify a central potential and a spin-orbit potential. These potentials are local in the cases considered here, but they are generally state-dependent.
### A Non-relativistic approach
The SHF model is very simple because the Hartree-Fock equations for the single-particle wave functions $`\varphi _i`$ and energies $`\epsilon _i`$ take the form of a Schrödinger equation with an effective mass $`m^{}(𝐫)`$ and a local, state-independent potential. All non-locality effects are described by this effective mass and consequently, the SHF equations contain also first derivatives of the wave function. It is natural to introduce an asymptotically equivalent wave function:
$`\stackrel{~}{\varphi }_i(𝐫)`$ $`=`$ $`\left({\displaystyle \frac{m^{}(𝐫)}{m}}\right)^{\frac{1}{2}}\varphi _i(𝐫),`$ (1)
which now satisfies a Schrödinger equation with a constant mass $`m`$ and a state-dependent local potential:
$`V(\epsilon _i,𝐫)`$ $`=`$ $`V_0(\epsilon _i,𝐫)+V_{LS}(𝐫)𝐥.𝐬.`$ (2)
The spin-orbit potential is (for spherical symmetry):
$`V_{LS}(q,r)`$ $`=`$ $`{\displaystyle \frac{m_q^{}(r)}{m}}[{\displaystyle \frac{w}{r}}[\rho ^{}(r)+\rho _q{}_{}{}^{}(r)]{\displaystyle \frac{1}{4r}}[(t_1x_1+t_2x_2)J(r)+(t_2t_1)J_q(r)]],`$ (3)
where $`w`$, $`t_i`$, $`x_i`$ are Skyrme force parameters, $`\rho _q`$ and $`J_q`$ are nucleon densities and spin densities ($`q`$= $`n`$ or $`p`$) with $`\rho =\rho _n+\rho _p`$ and $`J=J_n+J_p`$. The spin densities are practically zero in spin-saturated nuclei but they can give some contributions in spin-unsaturated ones. The above expression shows that $`V_{LS}`$ is surface peaked, and that proton spin-orbit splittings $`\mathrm{\Delta }_{LS}`$ should have little isotopic dependence in the SHF model as we shall see in the next sections. One should, however, keep in mind that another contribution to $`\mathrm{\Delta }_{LS}`$ comes from the energy-dependent central potential. One can easily check that the energy dependence of $`V_0(\epsilon _i,𝐫)`$ is entirely contained in a term $`\left(1m^{}(𝐫)/m\right)\epsilon _i`$.
Recently, some attempts have been made in order to introduce more freedom in the spin-orbit term of the Skyrme parametrization with the aim of improving isotope shift predictions in the Pb region. The authors of Ref. modify directly the Skyrme energy functional so that the first term of Eq.(3) becomes $`w\rho {}_{}{}^{}(r)/r+w{}_{}{}^{}\rho _{q}^{}{}_{}{}^{}(r)/r`$ and the spin density terms are dropped. It is found that a better description of data is reached when $`w^{}`$ is nearly $`w`$ (parameter set SkI4), i.e., when the proton spin-orbit potential depends almost entirely on the derivative of the neutron density.
### B Relativistic approach
Let us now calculate the spin-orbit potential in the general RHF case. The corresponding expressions for the RMFT case will be easily obtained by dropping all contributions corresponding to exchange (Fock) terms. Starting from the RHF equation for the 4-component spinors we shall obtain a Schrödinger-like equation for the upper component which allows us to identify the central and spin-orbit potentials. We use the notations introduced in Ref..
We start from Eq.(C3) of Ref. for the spinor $`(G_i,F_i)`$:
$`G_{i}^{}{}_{}{}^{}`$ $`=`$ $`({\displaystyle \frac{\kappa _i}{r}}\mathrm{\Sigma }_{T,i}P_i)G_i+(m+E_i+\mathrm{\Sigma }_{S,i}\mathrm{\Sigma }_{0,i}Q_i)F_i`$ (4)
$`F_{i}^{}{}_{}{}^{}`$ $`=`$ $`(mE_i+\mathrm{\Sigma }_{S,i}+\mathrm{\Sigma }_{0,i}+R_i)G_i+({\displaystyle \frac{\kappa _i}{r}}+\mathrm{\Sigma }_{T,i}+S_i)F_i,`$ (5)
where the $`\mathrm{\Sigma }`$’s are the direct (Hartree) self-energies whereas $`P,Q,R,S`$ are related to the exchange (Fock) contributions. We shall drop the state indices $`i`$ from now on. The first step is to eliminate the lower component $`F`$ and obtain a second order differential equation for the upper component $`G`$:
$`\mathrm{\Lambda }G{}_{}{}^{\prime \prime }+\alpha _1G{}_{}{}^{}+\alpha _0G`$ $`=`$ $`0,`$ (6)
where
$`\mathrm{\Lambda }`$ $`=`$ $`(m+E+\mathrm{\Sigma }_S\mathrm{\Sigma }_0Q)^1,`$ (7)
$`\alpha _1`$ $`=`$ $`\mathrm{\Lambda }{}_{}{}^{}+\mathrm{\Lambda }(PS),`$ (8)
$`\alpha _0`$ $`=`$ $`\mathrm{\Lambda }{}_{}{}^{}({\displaystyle \frac{\kappa }{r}}+\mathrm{\Sigma }_T+P)+\mathrm{\Lambda }({\displaystyle \frac{\kappa }{r^2}}+\mathrm{\Sigma }{}_{}{}^{}{}_{T}{}^{}+P{}_{}{}^{})`$ (10)
$`(ME+\mathrm{\Sigma }_S+\mathrm{\Sigma }_0+R)\mathrm{\Lambda }({\displaystyle \frac{\kappa }{r}}+\mathrm{\Sigma }_T+P)({\displaystyle \frac{\kappa }{r}}+\mathrm{\Sigma }_T+S).`$
Next, we look for a modified function $`\stackrel{~}{G}`$:
$`G`$ $`=`$ $`\lambda \stackrel{~}{G},`$ (11)
such that $`G`$ and $`\stackrel{~}{G}`$ are identical asymptotically while the differential equation satisfied by $`\stackrel{~}{G}`$ contains no first derivative. This condition determines $`\lambda `$ to be:
$`\lambda `$ $`=`$ $`C\mathrm{\Lambda }^{\frac{1}{2}}e^{\frac{1}{2}{\scriptscriptstyle (PS)𝑑r}}.`$ (12)
where the constant C is determined from the condition $`\stackrel{~}{G}^2𝑑r=1`$. In the limit of RMFT the exchange quantities $`P`$ and $`S`$ are zero and one recovers the familiar result $`\lambda =\mathrm{\Lambda }^{\frac{1}{2}}`$.
We can now write the differential equation satisfied by $`\stackrel{~}{G}`$ in the form:
$`{\displaystyle \frac{\mathrm{}^2}{2m}}\left(\stackrel{~}{G}{}_{}{}^{\prime \prime }+[{\displaystyle \frac{\lambda ^{\prime \prime }}{\lambda }}{\displaystyle \frac{1}{2}}({\displaystyle \frac{\alpha _1}{\mathrm{\Lambda }}})^2+{\displaystyle \frac{\alpha _0}{\mathrm{\Lambda }}}]\stackrel{~}{G}\right)=0.`$ (13)
The potential in state $`i`$ can be identified from the coefficient of $`\stackrel{~}{G}`$. It is possible to separate out the spin-orbit part $`V_{LS}`$ of this potential by observing that the quantities $`\mathrm{\Lambda },P,Q,R,S`$ depend on $`\kappa (2j+1)(lj)`$ (see, e.g., Ref.) which in turn can be expressed in terms of $`𝐥.𝐬`$. The complete expressions for $`V_{LS}`$ are given in the Appendix.
## III EVOLUTION OF SPLITTINGS ALONG ISOTOPIC CHAINS
We now turn to an analysis of spin-orbit splittings in some finite nuclei. We concentrate on the quantities:
$`\mathrm{\Delta }(nlj_<nlj_>)\epsilon _{nlj_<}\epsilon _{nlj_>}.`$ (14)
### A Ca isotopes
In Ca isotopes, the experimental situation concerning the $`1d_{5/2}`$ and $`1d_{3/2}`$ proton states is not very accurately established. In <sup>40</sup>Ca, the generally accepted value of $`1d_{3/2}`$ single-particle energy is $`\epsilon _{1d3/2}=8.3`$ MeV , , However, the values of $`\epsilon _{1d5/2}`$ are much more dispersed: -15.1 MeV , -15.5 MeV , -16 MeV , -14.3 MeV . Note that the most recent values of Ref. have been obtained as the centroid energies of the single-particle spectroscopic strength distributions:
$`\epsilon _{nlj}{\displaystyle \frac{_iE_iS_i}{_iS_i}}`$ (15)
and therefore, they take into account the strong fragmentation of the single-particle states. The fragment energies $`E_i`$ and spectroscopic factors $`S_i`$ were taken from Refs. ,,,, and . The corresponding experimental spin-orbit splittings of <sup>40</sup>Ca are plotted on Fig. 1. In <sup>48</sup>Ca, the experimental proton energies $`\epsilon _{1d3/2}`$ are respectively -16.2 MeV and -15.5 MeV whereas the $`\epsilon _{1d5/2}`$ energies are -21.5 MeV and -20.5 MeV . This gives the spin-orbit energies in <sup>48</sup>Ca shown in Fig. 1.
At this point we must mention that some theoretical papers, e.g., Ref. quote a somewhat smaller value for the experimental $`\mathrm{\Delta }_{S.O.}(1d3/21d5/2)`$ in <sup>48</sup>Ca based on an earlier paper but the experimental source remains uncertain.
In Fig. 1 we display the results calculated with the following models: SHF with a standard SIII force and with the SkI4 parametrization of Ref., RMFT with the commonly used NL-SH and NL3 parametrizations, RHF with ($`\sigma `$, $`\omega `$, $`\rho `$, $`\pi `$) mesons (model e) of Ref.). The two parametrizations NL-SH and NL3 have been chosen as the most representative and successful in describing nuclear ground states in a wide range of nuclei.
All calculated results for $`\mathrm{\Delta }(1d_{3/2}1d_{5/2})`$ show a linear A-dependence between $`A`$=40 to $`A`$=48. However, whereas the values of SHF and RMFT models decrease only slightly those of the RHF model exhibit a large decrease.
The relatively small variations of the SIII results can be understood by examining the $`V_{LS}`$ potential of Eq.(3). When going from <sup>40</sup>Ca to <sup>48</sup>Ca one acquires extra neutron densities $`\rho _{7/2}(r)`$ and $`J_{7/2}(r)`$ due to the filling of the $`1f_{7/2}`$ orbital while the core neutron and proton densities do not change much. For the SIII parametrization, $`x_1`$ = $`x_2`$ = 0 so that $`J_{7/2}(r)`$ gives no contribution to $`V_{LS}(p,r)`$ whereas $`\rho {}_{}{}^{}{}_{7/2}{}^{}(r)`$ yields a contribution with a node at the surface, i.e., a small effect on $`\mathrm{\Delta }(1d_{3/2}1d_{5/2})`$. The same analysis holds for the SkI4 results. For the RMFT case the interpretation of results is also simple. The spin-orbit potential reduces to the first term of Eq.(A6), i.e., $`2\mathrm{\Lambda }{}_{}{}^{}/\mathrm{\Lambda }r`$. With the definition of $`\mathrm{\Lambda }`$ given by Eq.(10) it is seen that the variation of spin-orbit potential is obtained by folding the derivative of the scalar and vector $`1f_{7/2}`$ densities with the (short-ranged) $`\sigma `$ and $`\omega `$ form factors and therefore, it has a node near the surface and the variation of spin-orbit splitting must be small. We should note, however, that the central potential that one can identify from Eq.(13) has some state dependence due to the effective mass $`m^{}=m+\mathrm{\Sigma }_S`$ and this effect can also contribute to the energy splitting. Furthermore, the kinetic energy and the central potential contribute also to this splitting through the j-dependence of the radial part in the wave function.
The RHF results are more difficult to interpret. From Eq.(13) we identify the kinetic energy $`T`$, central potential $`V_0`$ and spin-orbit potential $`V_{LS}`$ such that:
$`T_i+V_0_i+V_{LS}_i`$ $`=`$ $`\epsilon _i(1+{\displaystyle \frac{\epsilon _i}{2m}}).`$ (16)
The expressions for $`V_{LS}`$ given in the Appendix are complicated. The non-locality of the original RHF mean field induces a strong state dependence in the local equivalent potentials and therefore, variations with $`j_>`$ and $`j_<`$ will be found in all three terms of Eq.(16). This is illustrated in Table I where the values of the differences:
$`\mathrm{\Delta }T`$ $``$ $`T_{1d3/2}T_{1d5/2},`$ (17)
$`\mathrm{\Delta }V_0`$ $``$ $`V_0_{1d3/2}V_0_{1d5/2},`$ (18)
$`\mathrm{\Delta }V_{LS}`$ $``$ $`V_{LS}_{1d3/2}V_{LS}_{1d5/2},`$ (19)
are shown. In <sup>40</sup>Ca the calculated spin-orbit splitting is practically equal to $`\mathrm{\Delta }V_{LS}`$, with a strong cancellation of the two other terms. In <sup>48</sup>Ca $`\mathrm{\Delta }V_{LS}`$ is only one third of $`\epsilon _{1d3/2}\epsilon _{1d5/2}`$ , the rest coming from an increase of $`\mathrm{\Delta }V_0`$ and $`\mathrm{\Delta }T`$.
It is possible to analyze further the respective role of the isovector $`\pi `$ and $`\rho `$ mesons in the decrease of $`\mathrm{\Delta }(1d_{3/2}1d_{5/2})`$ from <sup>40</sup>Ca to <sup>48</sup>Ca. The key effect comes from the $`\pi `$ meson and the strong non-locality related to its light mass. If one artificially increases $`m_\pi `$ up to 600 MeV (and renormalize accordingly the $`f_\pi `$ coupling constant to keep reasonable single-particle levels) one tends to a more local situation and the strong decrease of $`\mathrm{\Delta }(1d_{3/2}1d_{5/2})`$ vanishes. The role of the $`\rho `$ meson can be studied by repeating the calculations without its contributions to the Fock terms. It is found that the decrease of $`\mathrm{\Delta }(1d_{3/2}1d_{5/2})`$ still remains.
### B Oxygen isotopes
The O chain is particularly interesting because of the prospects of reaching experimentally new, unstable isotopes with a large neutron excess. We examine here the evolution of the $`1p_{1/2}1p_{3/2}`$ splitting of proton levels. In Fig. 2 we summarize the results obtained with SHF (SIII and SkI4 parametrizations), RMFT (NL-SH and NL3 parametrizations) and RHF (model $`e)`$ of Ref.). We note that the results of NL-SH and NL3 are rather similar for spin-orbit splittings in O isotopes. This behaviour remains also in the Sn isotopes (see next subsection). These two parametrizations differ mostly by their predicted compression modulus K and therefore, this indicates that K has little effect on spin-orbit properties in the RMFT. One can distinguish three intervals ending at $`A`$=22, 24 and 28 which correspond to the filling of $`1d_{5/2}`$, $`2s_{1/2}`$ and $`1d_{3/2}`$ subshells, respectively. The trend in the $`1d_{5/2}`$ subshell resembles that of the Ca isotopes, namely a relatively small variation of $`\mathrm{\Delta }(1p_{1/2}1p_{3/2})`$ for the SHF and RMFT models and a large decrease for the RHF model. A similar discussion as for the Ca isotopes can be done here. It would be interesting to measure experimentally the variations of spin-orbit splittings in the chain of Oxygen isotopes. For the isotopes heavier than $`A`$=22 RMFT shows a maximum at $`A`$=24 followed by a decrease whereas for RHF the behaviour is opposite, with an increase of the spin-orbit splitting after $`A`$=22.
### C Sn isotopes
In the chain of Sn isotopes proton spin-orbit splittings are experimentally known in several nuclei. In Fig. 3 are shown the calculated and measured values of $`\mathrm{\Delta }(2p_{1/2}2p_{3/2})`$ for protons, and in Fig. 4 are displayed the calculated results for $`\mathrm{\Delta }(1f_{5/2}1f_{7/2})`$. From Fig. 3 it can be seen that none of the models is able to reproduce the very small empirical values of $`\mathrm{\Delta }(2p_{1/2}2p_{3/2})`$. In the range A=112-124 the SIII results show relatively small variations whereas the 3 other models give large fluctuations. Of course, it would be more satisfactory to calculate these nuclei using a Hartree-Fock-BCS description. It is not clear, however, that the proton spin-orbit splittings would be significantly affected so as to bring them into agreement with the data. Indeed, in Ref. a HF-BCS was performed for Sn isotopes using the effective interaction SkI4 and a density-dependent pairing force. The values of $`\mathrm{\Delta }(2p_{1/2}2p_{3/2})`$ thus obtained vary from 1.6 MeV in <sup>112</sup>Sn to 2.0 MeV in <sup>124</sup>Sn. As for the $`1f_{5/2}1f_{7/2}`$ splitting, Fig. 4 shows a wide range of predictions. For instance, in A=100 RHF is as low as 2.2 MeV while the other models are in the 5-6 MeV range . It would be very interesting to have data on this $`1f_{5/2}1f_{7/2}`$ case in order to shed light on this issue.
## IV CONCLUSION
We have examined in this work some predictions of spin-orbit splittings in the framework of non-relativistic and relativistic mean field approaches. We selected, as a representative case of the non-relativistic self-consistent approach, the Skyrme-Hartree-Fock model because it is widely used and its analytic simplicity lends itself to an easy interpretation of the calculated splittings. Numerical results were obtained with two versions of the effective Skyrme interaction, the standard parametrization SIII and the more recently proposed SkI4 which contains more degrees of freedom in its spin-orbit part. On the relativistic side, the RMFT is also a successful model for describing nuclear ground states and we have examined its spin-orbit predictions with the often used non-linear versions NL-SH and NL3. Here also the analytic form of the one-body spin-orbit potential is simple enough to allow some insight into the numerical results. In addition, we have included in our study the RHF model of Ref. in order to examine the role of the pion in the spin-orbit splittings.
One conclusion which can be drawn is that there is no clear advantage of the RMFT over the non-relativistic Skyrme-Hartree-Fock approach. It is true that in RMFT, the spin-orbit properties come out together with the central part of the mean field whereas in the non-relativistic approach one has the freedom of one or more additional spin-orbit parameters. Nevertheless, it cannot be concluded that the RMFT spin-orbit splittings describe the data particularly well. In Refs. , an interesting formal connection is made between RMFT and Skyrme-HF, but we can see here that, at a quantitative level, not only their spin-orbit predictions do differ, but even between SIII and SkI4 there are sizable differences in $`\mathrm{\Delta }_{S.O.}`$
All models fail to reproduce the very small $`2p_{1/2}2p_{3/2}`$ proton splittings in Sn isotopes and it does not seem that pairing correlations can improve this prediction. As for the RHF approach, the light pion mass produces strongly non-local Fock fields, i.e., a strong state dependence of the mean fields. Consequently, the splitting between spin-orbit partners is due not only to $`V_{LS}𝐥.𝐬_{lj}`$ but also to a large extent to the state dependence of the central potential. Thus, the evolution of the predicted splittings along an isotopic chain does not resemble that of less non-local models like RMFT or Skyrme. It appears that the isotopic dependence of $`\mathrm{\Delta }`$ disagrees with the data in the case of RHF.
The conclusion made earlier that RHF describes well the experimental $`\mathrm{\Delta }_{S.O.}`$ both in <sup>40</sup>Ca and <sup>48</sup>Ca was too premature in view of the uncertainties in the experimental values.
Finally, the question of spin-orbit predictions in the framework of self-consistent theories is still open. Existing parametrizations of effective interactions, relativistic as well as non-relativistic, need further improvements. The RMF predictions seem to follow a common trend whereas the Skyrme-HF results may differ somewhat depending on the parametrizations, as it can be seen with SIII and SkI4 in Sn isotopes. Improving the parametrizations necessitates better comparisons with experiment. In this respect, it would be very helpful to have more data on the evolution of spin-orbit splittings along isotopic chains. Without better data, it is not possible at the moment to state that the spin-orbit problem is understood either by RMFT or by the non-relastivistic approach.
###### Acknowledgements.
We would like to thank M. Bender and R. Wyss for their comments and A.M. Oros for the data on Ca isotopes. This work was supported in part by contract DGICYT PB97-0360 (Spain).
## A
In this Appendix we give the expressions for $`V_{LS}`$ in the general case of RHF. The special case of RMFT is easily obtained by setting all functions $`P,Q,R,S`$ to zero. The notations are those of Ref. .
Eq.(13) is a Schrödinger-type equation for a state of angular momentum ($`l,j`$). We split the coefficient of $`\stackrel{~}{G}`$ into a centrifugal term, a central potential and a spin-orbit potential:
$`{\displaystyle \frac{\lambda ^{\prime \prime }}{\lambda }}{\displaystyle \frac{1}{2}}({\displaystyle \frac{\alpha _1}{\mathrm{\Lambda }}})^2+{\displaystyle \frac{\alpha _0}{\mathrm{\Lambda }}}`$ $``$ $`{\displaystyle \frac{l(l+1)}{r^2}}{\displaystyle \frac{2m}{\mathrm{}^2}}[V_0(r)+V_{LS}(r)𝐥.𝐬_{lj}\epsilon (1+\epsilon /2m)].`$ (A1)
The functions $`P,Q,R,S`$ can be expressed as second order polynomials in $`\kappa `$ , namely, $`P=P_0+P_1k+P_2k^2`$ and similarly for $`Q,R,S`$. Here $`\kappa =(2j+1)(lj)=\frac{1}{2}\mathrm{\Omega }(2j+1)`$ where $`\mathrm{\Omega }=+1`$ if $`j=l\frac{1}{2}`$ and $`\mathrm{\Omega }=1`$ if $`j=l+\frac{1}{2}`$. It is easy to relate powers of $`\kappa `$ to $`𝐥.𝐬_{lj}`$. For instance, we have:
$`2𝐥.𝐬_{lj}`$ $`=`$ $`\mathrm{}^2(1+\kappa ).`$ (A2)
Thus, all contributions to $`V_{LS}(r)`$ can be identified by expressing the $`\kappa `$-dependence into a $`𝐥.𝐬_{lj}`$-dependence. One then obtains:
$`V_{LS}(r)`$ $`=`$ $`V_{LS}^{(I)}+V_{LS}^{(II)}+V_{LS}^{(III)},`$ (A3)
where
$`V_{LS}^{(I)}`$ $`=`$ $`{\displaystyle \frac{1}{m}}\{{\displaystyle \frac{\mathrm{\Lambda }^{}}{\mathrm{\Lambda }}}({\displaystyle \frac{1}{r}}+P_1)+{\displaystyle \frac{1}{r^2}}P{}_{}{}^{}{}_{1}{}^{}+{\displaystyle \frac{R_1}{\mathrm{\Lambda }}}+({\displaystyle \frac{1}{r}}+P_1)(\mathrm{\Sigma }_T^D+S_0)+({\displaystyle \frac{1}{r}}+S_1)(\mathrm{\Sigma }_T^D+P_0)`$ (A6)
$`+{\displaystyle \frac{\mathrm{\Lambda }^{}}{\mathrm{\Lambda }}}P_2+P{}_{}{}^{}{}_{2}{}^{}{\displaystyle \frac{R_2}{\mathrm{\Lambda }}}(\mathrm{\Sigma }_T^D+P_0)S_2({\displaystyle \frac{1}{r}}+P_1)({\displaystyle \frac{1}{r}}+S_1)P_2(\mathrm{\Sigma }_T^D+S_0)`$
$`+[l(l+1)+1][({\displaystyle \frac{1}{r}}+P_1)S_2+({\displaystyle \frac{1}{r}}+S_1)P_2][2l(l+1)+1]P_2S_2\},`$
$`V_{LS}^{(II)}`$ $`=`$ $`{\displaystyle \frac{1}{m}}\{{\displaystyle \frac{\mathrm{\Lambda }^{}}{\mathrm{\Lambda }}}(P_1S_1)(P_0S_0)(P_1S_1)`$ (A9)
$`+{\displaystyle \frac{\mathrm{\Lambda }^{}}{\mathrm{\Lambda }}}(P_2S_2)+{\displaystyle \frac{1}{2}}(P_1S_1)^2+(P_0S_0)(P_2S_2)`$
$`[l(l+1)+1](P_1S_1)(P_2S_2)+[2l(l+1)+1]{\displaystyle \frac{1}{2}}(P_2S_2)^2\},`$
$`V_{LS}^{(III)}`$ $`=`$ $`{\displaystyle \frac{1}{m}}\{{\displaystyle \frac{\mathrm{\Lambda }^{}}{2\mathrm{\Lambda }}}(P_1S_1)+{\displaystyle \frac{1}{2}}(P{}_{}{}^{}{}_{1}{}^{}S{}_{}{}^{}{}_{1}{}^{}){\displaystyle \frac{1}{2}}(P_0S_0)(P_1S_1)`$ (A12)
$`+{\displaystyle \frac{\mathrm{\Lambda }^{}}{2\mathrm{\Lambda }}}(P_2S_2){\displaystyle \frac{1}{2}}(P{}_{}{}^{}{}_{2}{}^{}S{}_{}{}^{}{}_{2}{}^{})+{\displaystyle \frac{1}{4}}(P_1S_1)^2+{\displaystyle \frac{1}{2}}(P_0S_0)(P_2S_2)`$
$`[l(l+1)+1]{\displaystyle \frac{1}{2}}(P_1S_1)(P_2S_2)+[2l(l+1)+1]{\displaystyle \frac{1}{4}}(P_2S_2)^2\}.`$
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# Color Superconductivity in High Density Effective Theory
## 1 Introduction
At high density, quarks in dense matter interact weakly with each other and form a Fermi sea, due to asymptotic freedom. When the energy is much less than the quark chemical potential ($`E\mu `$), only the quarks near the Fermi surface are relevant. The dynamics of quarks near the Fermi surface is effectively one-dimensional, since excitations along the Fermi surface do not cost any energy. The momentum perpendicular to the Fermi momentum just labels the degeneracy, similarly to the perpendicular momentum of charged particle under external magnetic field. This dimensional reduction due to the presence of Fermi surface makes possible for quarks to form a Cooper pair for any arbitrary weak attraction, since the critical coupling for the condensation in (1+1) dimensions is zero, known as the Cooper theorem in condensed matter.
While, in the BCS theory, such attractive force for electron Cooper pair is provided by phonons, for dense quark matter, where phonons are absent, the gluon exchange interaction provides the attraction, as one-gluon exchange interaction is attractive in the color anti-triplet channel. One therefore expects that color anti-triplet Cooper pairs will form and quark matter is color superconducting, which is indeed shown more than 20 years ago .
Recent development in color superconductivity, started from 1998, was spurred by recent two seminal works. The first one is by Alford, Rajagopal, and Wilczek , who convincingly argued that for three massless flavors, the ground state of quark matter is a so-called color-flavor locking (CFL) phase, in which the Cooper pair takes the following form, neglecting the small sextet component,
$$\psi _{L\alpha }^a(\stackrel{}{p})\psi _{L\beta }^b(\stackrel{}{p})=\psi _{R\alpha }^a(\stackrel{}{p})\psi _{R\beta }^b(\stackrel{}{p})=ϵ^{abI}ϵ_{\alpha \beta I}K(p_F),$$
(1)
where $`a,b(=1,2,3)`$ denote the color indices and $`\alpha ,\beta (=1,2,3)`$ denote the flavor indices.
The interesting feature of the CFL phase is that chiral symmetry is broken and the excitations in CFL phase have integral multiplet of electron charge. Though the usual quark-antiquark condensate is absent at high density, at least at the leading order, the chiral symmetry is spontaneously broken in the CFL phase. The flavor indices of Cooper pairs are locked to their color indices so that the unbroken symmetry that leaves the Cooper pair condensate invariant is the simultaneous rotation in the flavor and color space, breaking both color and chiral symmetry down to their diagonal subgroup,
$$SU(3)_c\times SU(3)_L\times SU(3)_RSU(3)_{c+L+R}.$$
(2)
The second work is done by Son , who showed that the Cooper pair gap in high density quark matter is very different from the usual BCS gap, due to the long range (color) magnetic interaction among quarks. By the renormalization group (RG) analysis, aided by the analysis of the Eliashberg equation, he found the Cooper pair gap depends on the coupling as,
$$\mathrm{\Delta }\frac{\mu }{g_s^5}\mathrm{exp}\left(\frac{3\pi ^2}{\sqrt{2}g_s}\right),$$
(3)
which was confirmed by more careful analysis .
In this talk, I will derive the above results in terms of a high density effective theory derived in . I will also calculate the critical temperature and the critical density and mention the mass of low-lying excitations in the CFL phase.
## 2 High density effective theory
QCD at high density has two distinct scales; one is an extrinsic scale, $`\mu `$, the quark chemical potential, and the other is the intrinsic scale, $`\mathrm{\Lambda }_{\mathrm{QCD}}`$. If the density is high enough, two scales are well separated, $`\mu \mathrm{\Lambda }_{\mathrm{QCD}}`$. To study a low-energy physics below a scale $`\mathrm{\Lambda }`$, an effective theory approach, where heavy modes ($`\omega >\mathrm{\Lambda }`$) are separated from light modes ($`\omega <\mathrm{\Lambda }`$) systematically, has been quite powerful.
Since we are interested in a cold dense matter where the relevant excitations are quasi-quarks near the Fermi surface, it will be useful to construct an effective theory that deals only with those relevant degrees of freedom . A dense matter with a fixed quark number is described by the QCD Lagrangian density with a chemical potential $`\mu `$,
$$_{\mathrm{QCD}}=\overline{\psi }i\overline{)}D\psi \frac{1}{4}F_{\mu \nu }^aF^{a\mu \nu }+\mu \overline{\psi }\gamma _0\psi ,$$
(4)
where the covariant derivative $`D_\mu =_\mu +ig_sA_\mu ^aT^a`$ and we neglect the mass of quarks for simplicity.
At energy just below $`\mu `$, we may decompose the quark momentum as
$$p^\mu =\mu v^\mu +l^\mu ,\left|l^\mu \right|<\mu ,$$
(5)
where $`\stackrel{}{v}_F`$ is a Fermi velocity and $`v^\mu =(0,\stackrel{}{v}_F)`$. We expand the quark propagator in powers of $`1/\mu `$:
$`S_F(p)={\displaystyle \frac{i}{(1+iϵ)p^0\gamma ^0\stackrel{}{p}\stackrel{}{\gamma }+\mu \gamma ^0}}=P_+{\displaystyle \frac{i\gamma ^0}{lV+iϵl^0}}+P_{}{\displaystyle \frac{i\gamma ^0}{2\mu }}+\mathrm{},`$ (6)
where $`V^\mu =(1,\stackrel{}{v}_F)`$ and the ellipsis denote higher order terms in $`1/\mu `$ expansion. In the second line of Eq. (6) we have introduced projection operators
$$P_\pm =\frac{1\pm \stackrel{}{\alpha }\stackrel{}{v}_F}{2},$$
(7)
where $`\stackrel{}{\alpha }=\gamma ^0\stackrel{}{\gamma }`$. The projection operators $`P_+`$ and $`P_{}`$ project out the states near the Fermi surface and the states in the Dirac sea, respectively. We see that the propagating modes are the states near the Fermi surface.
Using the techniques developed in heavy quark effective theory , we Fourier-decompose the quark field as
$$\psi (x)=\underset{\stackrel{}{v}_F}{}e^{i\mu \stackrel{}{v}_F\stackrel{}{x}}\psi (\stackrel{}{v}_F,x)$$
(8)
where
$$\psi (\stackrel{}{v}_F,x)=_{\left|l^\mu \right|<\mu }\frac{d^4l}{(2\pi )^4}\psi (\stackrel{}{v}_F,l)e^{ilx}\psi _+(\stackrel{}{v}_F,x)+\psi _{}(\stackrel{}{v}_F,x)$$
(9)
with $`\psi _\pm (\stackrel{}{v}_F,x)=P_\pm \psi (\stackrel{}{v}_F,x)`$. The low energy effective Lagrangian that consists of the light degrees of freedom (gluons and $`\psi _+`$) is obtained by matching all one-light-particle irreducible amplitudes in QCD with the vertex functions in the effective theory. As shown in , in the effective theory the quark propagator becomes
$$S_F(\stackrel{}{v}_F;l)=\frac{1+\stackrel{}{\alpha }\stackrel{}{v}_F}{2}\frac{i\gamma ^0}{lV+iϵl^0},$$
(10)
and in addition to the quark-gluon minimal coupling $`i\gamma ^0V^\mu g_s`$ there is marginal four-Fermi interaction for quarks with opposite Fermi velocities,
$`_{4\mathrm{f}}^1`$ $`=`$ $`{\displaystyle \frac{g_{us;tv}^S}{2\mu ^2}}[\psi _{L}^{}{}_{t}{}^{}(\stackrel{}{v}_F,x)\psi _{L}^{}{}_{s}{}^{}(\stackrel{}{v}_F,x)\psi _{L}^{}{}_{v}{}^{}(\stackrel{}{v}_F,x)\psi _{L}^{}{}_{u}{}^{}(\stackrel{}{v}_F,x)+(LR)]`$ (11)
$`+`$ $`{\displaystyle \frac{g_{us;tv}^P}{2\mu ^2}}[\psi _{L}^{}{}_{t}{}^{}(\stackrel{}{v}_F,x)\psi _{L}^{}{}_{s}{}^{}(\stackrel{}{v}_F,x)\psi _{R}^{}{}_{v}{}^{}(\stackrel{}{v}_F,x)\psi _{R}^{}{}_{u}{}^{}(\stackrel{}{v}_F,x)+(LR)].`$
To summarize, the high density effective theory has several interesting features: (1) In the leading order, only $`\gamma ^0`$ enters in the Dirac matrices. (2) Anti-quarks are systematically decoupled. (3) There appear marginal four-quark operators naturally. (4) It offers a systematic high-density expansion.
## 3 Cooper pair gap
To describe the Cooper-pair gap equation, we introduce a 8-component field, following the Nambu-Gorkov formalism, $`\mathrm{\Psi }(\stackrel{}{v}_F,x)(\psi (\stackrel{}{v}_F,x),\psi _c(\stackrel{}{v}_F,x))^T`$, where we reverted the notation $`\psi `$ for $`\psi _+`$ and introduced the charge conjugate field $`\psi _c(\stackrel{}{v}_F,x)=C\overline{\psi }^T(\stackrel{}{v}_F,x)`$. The charge conjugation matrix, $`C`$, satisfies $`C^1\gamma _\mu C=\gamma _\mu ^T`$. The inverse propagator for the Nambu-Gorkov field is
$$S^1(\stackrel{}{v}_F,l)=i\gamma _0\left(\begin{array}{cc}lV& \mathrm{\Delta }^{}(l_{})\\ \mathrm{\Delta }(l_{})& l\overline{V}\end{array}\right),$$
(12)
where $`\overline{V}^\mu =(1,\stackrel{}{v}_F)`$ and $`\mathrm{\Delta }`$ is the Cooper-pair gap.
The effective action for the fermion two-point function $`S`$ is given as
$$\mathrm{\Gamma }=\mathrm{Tr}\mathrm{ln}S^1+\mathrm{Tr}\left(S^1S_0^1\right)S+\left(2\mathrm{P}\mathrm{I}\mathrm{diagrams}\right),$$
(13)
where the 2PI diagrams are two-particle irreducible vacuum diagrams. For the gluon propagator, we use an in-medium propagator, which is in the hard dense loop (HDL) approximation given as
$$iD_{\mu \nu }(k)=\frac{P_{\mu \nu }^T}{k^2G}+\frac{P_{\mu \nu }^L}{k^2F}\xi \frac{k_\mu k_\nu }{k^4},$$
(14)
where $`\xi `$ is the gauge parameter and the projectors are defined by
$`P_{ij}^T`$ $`=`$ $`\delta _{ij}{\displaystyle \frac{k_ik_j}{|\stackrel{}{k}|^2}},P_{00}^T=0=P_{0i}^T`$ (15)
$`P_{\mu \nu }^L`$ $`=`$ $`g_{\mu \nu }+{\displaystyle \frac{k_\mu k_\nu }{k^2}}P_{\mu \nu }^T.`$ (16)
The medium effect is incorporated in $`F`$ and $`G`$, which becomes in the weak coupling limit ($`\left|k_0\right|\left|\stackrel{}{k}\right|`$)
$`F(k_0,\stackrel{}{k})M^2,G(k_0,\stackrel{}{k}){\displaystyle \frac{\pi }{4}}M^2{\displaystyle \frac{k_0}{|\stackrel{}{k}|}},`$ (17)
where $`M=\sqrt{N_f/2}g_s\mu /\pi `$, the Debye mass. The gap equations, obtained by extremizing the effective action, $`0=\delta \mathrm{\Gamma }/\delta S`$, are given in Euclidean space as
$`\mathrm{\Delta }(p_{})`$ $`=`$ $`{\displaystyle }{\displaystyle \frac{d^4q}{(2\pi )^4}}[{\displaystyle \frac{2}{3}}g_s^2\{{\displaystyle \frac{VP^T\overline{V}}{(pq)_{}^2+\stackrel{}{q}_{}^2+\frac{\pi }{4}M^2|p_0q_0|/|\stackrel{}{p}\stackrel{}{q}|}}`$ (18)
$``$ $`{\displaystyle \frac{1}{(pq)_{}^2+\stackrel{}{q}_{}^{}{}_{}{}^{2}+M^2}}\xi {\displaystyle \frac{(pq)_{}^2}{(pq)^4}}\}+{\displaystyle \frac{g_{\overline{3}}}{\mu ^2}}]{\displaystyle \frac{\mathrm{\Delta }(q_{})}{q_{}^2+\mathrm{\Delta }^2(q_{})}},`$
where $`g_{\overline{3}}`$ is a four-Fermi coupling. Since the gluon coupling is vectorial, the gluon exchange interaction in the gap equation does not distinguish the handedness of quarks and thus it will generate same condensates regardless of handedness; $`\left|\psi _L\psi _L\right|=\left|\psi _R\psi _R\right|=\left|\psi _L\psi _R\right|`$, suppressing other quantum numbers. But, the four-Fermi interaction in the effective Lagrangian, Eq. (11), lifts the degeneracy, since the gap in $`LL`$ or $`RR`$ channel will be bigger than the one in $`LR`$ channel due to the difference in the four-Fermi couplings, $`g^S>g^P`$. The $`LL`$ or $`RR`$ condensate is energetically more preferred than the $`LR`$ condensate. We also note that since in the effective theory the gluons are blind not only to flavors but also to the Dirac indices of quarks, the diquark Cooper-pair can be written as color anti-triplet.
Since quarks are anti-commuting, the only possible way to form diquark (S-wave) condensate is either in spin-singlet or in spin-triplet:
$`\psi _{L}^{}{}_{i\alpha }{}^{a}(\stackrel{}{v}_F,x)\psi _{L}^{}{}_{j\beta }{}^{b}(\stackrel{}{v}_F,x)`$ $`=`$ $`\psi _{R}^{}{}_{i\alpha }{}^{a}(\stackrel{}{v}_F,x)\psi _{R}^{}{}_{j\beta }{}^{b}(\stackrel{}{v}_F,x)`$ (19)
$`=`$ $`ϵ_{ij}ϵ^{abc}K_{\left[\alpha \beta \right]c}(p_F)+\delta _{ij}ϵ^{abc}K_{\left\{\alpha \beta \right\}c}(p_F),`$ (20)
where $`a,b,c=1,2,3`$ are color indices, $`\alpha ,\beta ,\gamma =u,d,s,\mathrm{},N_f`$ flavor indices, and $`i,j=1,2`$ spinor indices. Indices in the bracket and in the curled bracket are anti-symmetrized and symmetrized, respectively. But, the spin-one component of the gap, $`K_{\{\alpha \beta \}c}`$, vanishes algebraically, since $`\psi (\stackrel{}{v}_F,x)=1/2\left(1+\stackrel{}{\alpha }\stackrel{}{v}_F\right)\psi (\stackrel{}{v}_F,x)`$ and $`(1+\stackrel{}{\alpha }\stackrel{}{v}_F)_{il}(1\stackrel{}{\alpha }\stackrel{}{v}_F)_{lj}=0`$.
When $`N_f=3`$, the spin-zero component of the condensate becomes (flavor) anti-triplet,
$$K_{\left[\alpha \beta \right]c}(p_F)=ϵ_{\alpha \beta \gamma }K_c^\gamma (p_F).$$
(21)
Using the global color and flavor symmetry, one can always diagonalize the spin-zero condensate as $`K_c^\gamma =\delta _c^\gamma K_\gamma `$. To determine the parameters, $`K_u`$, $`K_d`$, and $`K_s`$, we need to minimize the vacuum energy for the condensate. The vacuum energy is given as in the leading HDL approximation
$`V(\mathrm{\Delta }){\displaystyle \frac{\mu ^2}{4\pi }}{\displaystyle \underset{i=1}{\overset{9}{}}}{\displaystyle \frac{d^2l_{}}{(2\pi )^2}\left[\mathrm{ln}\left(\frac{l_{}^2}{l_{}^2+\mathrm{\Delta }_i^2(l_{})}\right)+\frac{1}{2}\frac{\mathrm{\Delta }_i^2(l_{})}{l_{}^2+\mathrm{\Delta }_i^2(l_{})}\right]},`$ (22)
where $`\mathrm{\Delta }_i`$’s are the eigenvalues of color anti-symmetric and flavor anti-symmetric $`9\times 9`$ gap, $`\mathrm{\Delta }_{\alpha \beta }^{ab}`$.
Approximating $`\mathrm{\Delta }_i`$ to be constant, one can easily perform the momentum integration in (22) to get
$`V(\mathrm{\Delta })0.43{\displaystyle \frac{\mu ^2}{4\pi ^2}}{\displaystyle \underset{i}{}}\left|\mathrm{\Delta }_i(0)\right|^2.`$ (23)
Were $`\mathrm{\Delta }_i`$ independent of each other, the global minimum should occur at $`\mathrm{\Delta }_i(0)=\mathrm{const}.`$ for all $`i=1,\mathrm{},9`$. But, due to the global color and flavor symmetry, only three of them are independent. Similarly to the condensate, the gap can be also diagonalized by the color and flavor symmetry as
$$\mathrm{\Delta }_{ab}^{\alpha \beta }=ϵ_{\alpha \beta \gamma }ϵ^{abc}\mathrm{\Delta }_\gamma \delta _c^\gamma .$$
(24)
Without loss of generality, we can take $`\left|\mathrm{\Delta }_u\right|\left|\mathrm{\Delta }_d\right|\left|\mathrm{\Delta }_s\right|`$. Let $`\mathrm{\Delta }_d/\mathrm{\Delta }_u=x`$ and $`\mathrm{\Delta }_s/\mathrm{\Delta }_u=y`$. Then, the vacuum energy becomes
$$V(\mathrm{\Delta })0.43\frac{\mu ^2}{4\pi ^2}\left|\mathrm{\Delta }_u\right|^2f(x,y),$$
(25)
where $`f(x,y)`$ is a complicate function of $`1x,y1`$ that has a maximum at $`x=1=y`$, $`f(x,y)13.4`$. Therefore, the vacuum energy has a global minimum when $`\mathrm{\Delta }_u=\mathrm{\Delta }_d=\mathrm{\Delta }_s`$, or in terms of the eigenvalues of the gap
$$\mathrm{\Delta }_i=\mathrm{\Delta }_u(1,1,1,1,1,1,1,1,2).$$
(26)
Now, we analyze the SD gap equation Eq. (18) to see if it admits a nontrivial solution. Since the color-flavor locking gap is preferred if it exists, we may write the gap as
$$\mathrm{\Delta }_{\alpha \beta }^{ab}=ϵ^{abI}ϵ_{\alpha \beta I}\mathrm{\Delta }.$$
(27)
We first note that the main contribution to the integration comes from the loop momenta in the region $`q_{}^2\mathrm{\Delta }^2`$ and $`|\stackrel{}{q}_{}|M^{2/3}\mathrm{\Delta }^{1/3}`$. Therefore, we find that the leading contribution is by the first term due to the Landau-damped magnetic gluons. For this momentum range, we can take $`|\stackrel{}{p}\stackrel{}{q}||\stackrel{}{q}_{}|`$ and
$`VP^T\overline{V}=v_F^iv_F^j\left(\delta _{ij}\widehat{k}_i\widehat{k}_j\right)=1+O\left({\displaystyle \frac{\mathrm{\Delta }^{4/3}}{M^{4/3}}}\right).`$ (28)
We also note that the term due to the four-Fermi operator is negligible, since $`g_{\overline{3}}g_s^4`$ at the matching scale $`\mu `$.
Neglecting $`(pq)_{}^2`$ in the denominator to integrate over $`\stackrel{}{q}_{}`$, we get
$`\mathrm{\Delta }(p_{})={\displaystyle \frac{g_s^2}{9\pi }}{\displaystyle \frac{d^2q_{}}{(2\pi )^2}\frac{\mathrm{\Delta }(q_{})}{q_{}^2+\mathrm{\Delta }^2}\left[\mathrm{ln}\left(\frac{\mu ^3}{\frac{\pi }{4}M^2|p_0q_0|}\right)+\frac{3}{2}\mathrm{ln}\left(\frac{\mu ^2}{M^2}\right)+\frac{3}{2}\xi \right]}.`$ (29)
We see that in this approximation $`\mathrm{\Delta }(p_{})\mathrm{\Delta }(p_0)`$. Then, we can integrate over $`\stackrel{}{v}_F\stackrel{}{q}`$ to get
$`\mathrm{\Delta }(p_0)`$ $`=`$ $`{\displaystyle \frac{g_s^2}{36\pi ^2}}{\displaystyle _\mu ^\mu }𝑑q_0{\displaystyle \frac{\mathrm{\Delta }(q_0)}{\sqrt{q_0^2+\mathrm{\Delta }^2}}}\mathrm{ln}\left({\displaystyle \frac{\overline{\mathrm{\Lambda }}}{|p_0q_0|}}\right)`$ (30)
where $`\overline{\mathrm{\Lambda }}=4\mu /\pi (\mu /M)^5e^{3/2\xi }`$. If we take $`\mathrm{\Delta }\mathrm{\Delta }(0)`$ for a rough estimate of the gap,
$$1=\frac{g_s^2}{36\pi ^2}\left[\mathrm{ln}\left(\frac{\overline{\mathrm{\Lambda }}}{\mathrm{\Delta }}\right)\right]^2\mathrm{or}\mathrm{\Delta }\overline{\mathrm{\Lambda }}\mathrm{exp}\left(\frac{6\pi }{g_s}\right).$$
(31)
To take into account the energy dependence of the gap, we convert the Schwinger-Dyson equation (30) into a differential equation, approximating the kernel as
$$\mathrm{ln}|p_0q_0|\mathrm{ln}[\mathrm{max}.(|p_0|,|q_0|)],$$
(32)
to get
$$p\mathrm{\Delta }^{\prime \prime }(p)+\mathrm{\Delta }^{}(p)+\frac{2\alpha _s}{9\pi }\frac{\mathrm{\Delta }(p)}{\sqrt{p^2+\mathrm{\Delta }^2}}=0,$$
(33)
with boundary conditions $`p\mathrm{\Delta }^{}=0`$ at $`p=\mathrm{\Delta }`$ and $`\mathrm{\Delta }=0`$ at $`p=\overline{\mathrm{\Delta }}`$, where $`pp_0`$. When $`p\mathrm{\Delta }(p)`$, the equation becomes
$$p\mathrm{\Delta }^{\prime \prime }+\mathrm{\Delta }^{}+\frac{r^2}{4}\frac{\mathrm{\Delta }(p)}{|\mathrm{\Delta }|}=0,$$
(34)
where $`r^2=2g_s^2/(9\pi ^2)`$ and $`|\mathrm{\Delta }|`$ is the gap at $`p=0`$. We find $`\mathrm{\Delta }(p)=|\mathrm{\Delta }|J_0\left(r\sqrt{p/|\mathrm{\Delta }|}\right)`$ for $`p\mathrm{\Delta }|`$. When $`p\mathrm{\Delta }`$, the differential equation (33) becomes
$$p\mathrm{\Delta }^{\prime \prime }+\mathrm{\Delta }^{}+\frac{r^2}{4}\frac{\mathrm{\Delta }}{p}=0,$$
(35)
whose solution is $`\mathrm{\Delta }(p)=B\mathrm{sin}\left[(r/2)\mathrm{ln}\overline{\mathrm{\Lambda }}/p\right]`$. By matching two solutions at the boundary $`p=|\mathrm{\Delta }|`$ we get
$$B|\mathrm{\Delta }|\mathrm{and}|\mathrm{\Delta }|=\overline{\mathrm{\Lambda }}e^{\pi /r}.$$
(36)
The gap is therefore given as at the leading order in the weak coupling expansion
$$\left|\mathrm{\Delta }\right|=c\frac{\mu }{g_s^5}\mathrm{exp}\left(\frac{3\pi ^2}{\sqrt{2}g_s}\right),$$
(37)
where $`c=2^7\pi ^4N_f^{5/2}e^{3\xi /2+1}`$. This agrees with the RG analysis done by Son (see also ) and also with the Schwinger-Dyson approach in full QCD . The $`1/g_s`$ behavior of the exponent of the gap at high density is due to the double logarithmic divergence in the gap equation (18), similarly to the case of chiral symmetry breaking under external magnetic fields .
## 4 Critical density and temperature
In this section we calculate the critical density and temperature. First, we add the $`1/\mu `$ corrections to the gap equation Eq. (18) to see how the formation of Cooper pair changes when the density decreases. The leading $`1/\mu `$ corrections to the quark-gluon interactions are
$$_1=\frac{1}{2\mu }\underset{\stackrel{}{v}_F}{}\psi ^{}(\stackrel{}{v}_F,x)\left(\gamma _{}D\right)^2\psi (\stackrel{}{v}_F,x)=\underset{\stackrel{}{v}_F}{}\left[\psi ^{}\frac{D_{}^2}{2\mu }\psi +g_s\psi ^{}\frac{\sigma _{\mu \nu }F^{\mu \nu }}{4\mu }\psi \right].$$
(38)
In the leading order in the HDL approximation, the loop correction to the vertex is neglected and the quark-gluon vertex is shifted by the $`1/\mu `$ correction as
$$ig_sv_F^iig_sv_F^iig_s\frac{l_{}^i}{\mu },$$
(39)
where $`l_i`$ is the momentum carried away from quarks by gluons. We note that since the $`1/\mu `$ correction to the quark-gluon vertex does not depend on the Fermi velocity of the quark, it generates a repulsion for quark pairs. For a constant gap approximation, $`\mathrm{\Delta }(p_{})\mathrm{\Delta }`$, the gap equation becomes in the leading order, as $`p0`$,
$`1={\displaystyle \frac{g_s^2}{9\pi }}{\displaystyle \frac{\mathrm{d}^2l_{}}{(2\pi )^2}\left[\mathrm{ln}\left(\frac{\overline{\mathrm{\Lambda }}}{|l_0|}\right)\frac{3}{2}\right]\frac{1}{l_{}^2+\mathrm{\Delta }^2}}={\displaystyle \frac{g_s^2}{36\pi ^2}}\mathrm{ln}\left({\displaystyle \frac{\overline{\mathrm{\Lambda }}}{\mathrm{\Delta }}}\right)\left[\mathrm{ln}\left({\displaystyle \frac{\overline{\mathrm{\Lambda }}}{\mathrm{\Delta }}}\right)3\right].`$ (40)
When $`\mu <\mu _ce^3\mathrm{\Delta }`$, the gap due to the long-range color magnetic interaction disappears. Since the phase transition for color superconducting phase is believed to be of first order , we may assume that the gap has the same dependence on the chemical potential $`\mu `$ as the leading order. Then, the critical density for the color superconducting phase transition is given by
$$\mu _c=e^3\mu _c\mathrm{exp}\left[\frac{3\pi ^2}{\sqrt{2}g_s(\mu _c)}\right].$$
(41)
Therefore, if the strong interaction coupling is too strong at the scale of the chemical potential, the gap does not form. To form the Cooper pair gap, the strong coupling at the scale of the chemical potential has to be smaller than $`g_s(\mu _c)=\pi ^2/\sqrt{2}`$. By using the one-loop $`\beta `$ function for three light flavors, $`\beta (g_s)=9/(16\pi ^2)g_s^3`$, and the experimental value for the strong coupling constant, $`\alpha _s(1.73\mathrm{GeV})=0.32_{0.053}^{+0.031}(\mathrm{exp})\pm 0.016(\mathrm{theo})`$ , we get $`0.13\mathrm{GeV}<\mu _c<0.31\mathrm{GeV}`$, which is about the same order as the one estimated by the instanton induced four-Fermi interaction or by general effective four-Fermi interactions . But, this should be taken as an order of magnitude, since for such a small chemical potential the higher order terms in $`1/\mu `$ expansion, which we have neglected, are as important as the leading term.
We now consider the temperature effect on the gap, which is quite important to understand the heavy ion collision or the final stage of the evolution of giant stars. The super dense and hot quark matter will go through a phase transition as it cools down by emitting weakly interacting particles like neutrinos.
At finite temperature, $`T`$, the gap equation (18) becomes, following the imaginary-time formalism developed by Matsubara,
$`\mathrm{\Delta }(\omega _n^{})={\displaystyle \frac{g_s^2}{9\pi }}T{\displaystyle \underset{n=\mathrm{}}{\overset{+\mathrm{}}{}}}{\displaystyle \frac{\mathrm{d}q}{2\pi }\frac{\mathrm{\Delta }(\omega _n)}{\omega _n^2+\mathrm{\Delta }^2(\omega _n)+q^2}\mathrm{ln}\left(\frac{\overline{\mathrm{\Lambda }}}{\left|\omega _n^{}\omega _n\right|}\right)},`$ (42)
where $`\omega _n=\pi T(2n+1)`$ and $`q\stackrel{}{v}_F\stackrel{}{q}`$. We now use the constant (but temperature-dependent) gap approximation, $`\mathrm{\Delta }(\omega _n)\mathrm{\Delta }(T)`$ for all $`n`$. Taking $`n^{}=0`$ and converting the logarithm into integration, we get
$`\mathrm{\Delta }(T)={\displaystyle \frac{g_s^2}{18\pi }}T{\displaystyle \underset{n=\mathrm{}}{\overset{+\mathrm{}}{}}}{\displaystyle \frac{\mathrm{d}q}{2\pi }_0^{\overline{\mathrm{\Lambda }}^2}dx\frac{\mathrm{\Delta }(T)}{\omega _n^2+\mathrm{\Delta }^2(T)+q^2}\frac{1}{x+(\omega _n\omega _0)^2}}.`$ (43)
Using the contour integral, one can in fact sum up over all $`n`$ to get
$`1={\displaystyle \frac{g_s^2T}{36\pi ^2}}{\displaystyle dq_0^{\overline{\mathrm{\Lambda }}^2}dx\frac{1}{2\pi i}_C\frac{\mathrm{d}\omega }{1+e^{\omega /T}}\frac{1}{\left(\omega ^2q^2\mathrm{\Delta }^2\right)\left[(\omega _ni\omega _0)^2+x\right]}}.`$ (44)
Since the gap vanishes at the critical temperature, $`\mathrm{\Delta }(T_C)=0`$, after performing the contour integration in Eq. (44), we get
$`1`$ $`=`$ $`{\displaystyle \frac{g_s^2}{36\pi ^2}}{\displaystyle }\mathrm{d}q{\displaystyle _0^{\overline{\mathrm{\Lambda }}^2}}\mathrm{d}x\{{\displaystyle \frac{(\pi T_C)^2+xq^2}{\left[(\pi T_C)^2+xq^2\right]^2+(2\pi T_Cq)^2}}{\displaystyle \frac{\mathrm{tanh}\left[q/(2T_C)\right]}{2q}}`$ (45)
$`+{\displaystyle \frac{(\pi T_C)^2+q^2x}{\left[(\pi T_C)^2+q^2x\right]^2+(2\pi T_C)^2x}}{\displaystyle \frac{\mathrm{coth}\left[\sqrt{x}/(2T_C)\right]}{\sqrt{2}}}\}.`$
At high density $`\overline{\mathrm{\Lambda }}T_C`$, the second term in the integral in Eq. (45) is negligible, compared to the first term, and integrating over $`x`$, we get
$`1`$ $`=`$ $`{\displaystyle \frac{g_s^2}{36\pi ^2}}{\displaystyle _0^{\lambda _c}}dy{\displaystyle \frac{\mathrm{tanh}y}{y}}\left[\mathrm{ln}\left({\displaystyle \frac{\lambda _c^2}{(\pi /2)^2+y^2}}\right)+O\left({\displaystyle \frac{y^2}{\lambda _c^2}}\right)\right]`$
$`=`$ $`{\displaystyle \frac{g_s^2}{36\pi ^2}}[\left(\mathrm{ln}\lambda _c\right)^2+2A\mathrm{ln}\lambda _c+\mathrm{const}.]`$
where we have introduced $`yq/(2T_C)`$ and $`\lambda _c\overline{\mathrm{\Lambda }}/(2T_C)`$ and $`A`$ is given as
$$A=_0^1dy\frac{\mathrm{tanh}y}{y}+_1^{\mathrm{}}dy\frac{\mathrm{tanh}y1}{y}=\mathrm{ln}\left(\frac{4}{\pi }\right)+\gamma .$$
(46)
Therefore, we find, taking the Euler-Mascheroni constant $`\gamma 0.577`$,
$$T_C=\frac{e^A}{2}\mathrm{\Delta }1.13\mathrm{\Delta },$$
(47)
which shows that the ratio between the critical temperature and the Cooper-pair gap in color superconductivity is same as the BCS value, $`e^\gamma /\pi 0.57`$ .
## 5 More on CFL
As pointed out by Schäfer and Wilczek , the low-lying particle spectrum of the CFL phase at high density resembles that of low density hadron phase. Both phases have pions and kaons, arising from the chiral symmetry breaking. The baryons and mesons at high density have integral multiplet of the electron charge, the charge corresponding to the unbroken $`U(1)`$ gauge symmetry at high density. Since the diquark condensate provides additional baryon number $`B=2/3`$, quarks in color superconductor have baryon number $`B=1`$.
To describe the dynamics of pions and kaons, the chiral Lagrangian for the CFL phase at high density has been constructed and it is shown in that quarks in the CFL phase is realized as a topological soliton, called superqualiton, as baryons in the hadron phase at low density. Unlike the low density phase, the parameters in the chiral Lagrangian can be calculated from the microscopic theory. For instance, the mass of Nambu-Goldstone bosons is found to be
$$m_{NG}^2m_q^2\mathrm{\Delta }\overline{\mathrm{\Delta }}\mathrm{ln}(\mu ^2/\mathrm{\Delta }^2)/\mu ^2,$$
(48)
showing that mesons become massless at asymptotically large chemical potential, as the Dirac mass term, $`m_q\overline{\psi }_+\psi _{}(m_q^2/\mu )\psi _+^{}\psi _+`$, vanishes for infinite density. (See also .) This is confirmed subsequently . Another interesting feature of meson mass is that the mass hierarchy is reversed . For instance, if $`m_s>m_{u,d}`$,
$$m_K<m_\pi .$$
(49)
This inverse mass hierarchy is due to the fact that what we call a kaon in the CFL phase is the fluctuation of Cooper-pairs in the up and down flavor spaces,
$$U_{L}^{}{}_{a\alpha }{}^{}(x)\underset{yx}{lim}\frac{\left|xy\right|^{\gamma _m}}{\kappa }ϵ^{ij}ϵ_{abc}ϵ_{\alpha \beta \gamma }\psi _{Li}^{b\beta }(\stackrel{}{v}_F,x)\psi _{Lj}^{c\gamma }(\stackrel{}{v}_F,y),$$
(50)
where $`\gamma _m`$ is the anomalous dimension.
## 6 Conclusion
I have discussed the exciting recent developments in color superconductivity in high density quark matter in terms of an effective theory formalism. I have shown that the effective theory calculation reproduces recent results on the Cooper pair gap, the critical temperature, and on the ground state of high density QCD. It not only simplifies the calculation very much but also allows us to estimate the critical density.
I wish to thank the organizers of the TMU-Yale symposium for the wonderful meeting. I am grateful to T. Lee, D.-P. Min, V. Miransky, M. Rho, I. Shovkovy, L. C. R. Wijewardhana, and I. Zahed for the collaborations on the works described here and to S. Hsu, R. Pisarski, D. Rischke, and T. Schäfer for helping me to understand this subject. My research is supported by the academic research fund of Ministry of Education, Republic of Korea, Project No. BSRI-99-015-DI0114.
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# Hadronic matrix elements from the lattice 11footnote 1To appear in Proceedings of The Third International Conference on B Physics and CP Violation, Taipei, December 3-7, 1999, H. -Y. Cheng and W. -S. Hou, eds. (World Scientific, 2000).
## 1 Introduction
Hadronic matrix elements are of crucial importance for constraining the parameters of the Standard Model (SM) in conjunction with experimental information. Lattice approach has already attained considerable success in handling $`B`$-parameters and decay constants and to a lesser degree semi-leptonic form factors. Weak decays into purely hadronic final states, of which $`K\pi \pi `$ is the simplest example, are extremely problematic. After almost a decade the topic is again getting some attention due in part to important developments with regard to chiral symmetry on the lattice.
## 2 Heavy-light decay constants: effects of quenching.
Lattice methods have made considerable progress in pinning down, fairly accurately, $`f_B`$ and other heavy-light decay constants, in the quenched approximation (see Table 1). Unfortunately, the effects of quenching could be substantial; current estimates place them around $`(20\pm 10)\%`$, seriously limiting the phenomenological applications. Since the past few years there is heightened activity to accurately ascertain the effects of quenching. For technical reasons, the dynamical simulations are with $`N_f=2`$ i.e. with only two “light” sea-quarks rather than the three ($`u,d,s`$) in real life. Furthermore, the mass of the sea-quarks tends to be relatively heavy. The indications from these dynamical simulations is that $`f_B^{\mathrm{dynamical}}>f_B^{\mathrm{quenched}}`$.
## 3 Hints from Experiment: $`f_{D_S}^{\mathrm{expt}}>f_{D_S}^{\mathrm{QQCD}}`$ ?
Table 2 exhibits a compilation of the results for $`f_{D_S}`$ from several experiments. Curiously, the central value of all but one experiment is above the value for $`f_{D_S}`$ from quenched QCD simulations: $`f_{D_S}^{\mathrm{QQCD}}=225\pm 20`$ MeV. Since the errors in the existing experimental numbers are rather large we cannot draw strong conclusions; it seems plausible, nevertheless, that these experiments are also indicating that quenched QCD tends to underestimate the heavy-light pseudoscalar decay constants.
## 4 Precise studies of $`f_{D_S}`$
For experiment as well as for dynamical lattice simulations, a precise determination of $`f_{B_D}(f_{B_S})`$ is significantly more problematic than $`f_{D_S}`$; in fact experimentally direct determination of $`f_B`$ is not of immediate reach. Therefore, it would be very useful if the experimental as well as the lattice determinations of $`f_{D_S}`$ could be improved to 10–15% accuracy. A comparison between the two would then serve as a very useful guide for correcting $`f_{B_d}(f_{B_S})`$ from the lattice. In this regard the ratios $`f_{B_d}/f_{D_S}`$ and $`f_{B_S}/f_{D_S}`$ from the lattice would clearly be useful.
## 5 Heavy-light $`B`$ parameters
Lattice has had success with heavy-light $`B`$-parameters for a long time although calculations in the static approximation are still somewhat problematic. Table 3 presents a brief summary. Recall that a precise value of the ratio $`f_{B_S}(B_{B_S})^{1/2}/f_{B_d}(B_{B_d})^{1/2}`$ is needed in accurately deducing $`V_{td}/V_{ts}`$ once $`B_S`$-$`\overline{B}_S`$ oscillations get experimentally measured.
In view of the importance of the aforesaid ratio it may be useful to determine it also more directly from the SU(3) breaking ratio:
$$\frac{M_{B_S}(\mu )}{M_{B_d}(\mu )}=\frac{\overline{B}_{B_S}|(\overline{b}\gamma _\rho (1\gamma _5)s)^2|B_{B_S}}{\overline{B}_{B_d}|(\overline{b}\gamma _\rho (1\gamma _5)d)^2|B_{B_d}}$$
Many of the systematic errors should cancel in this ratio. However, in the attempts that have been so far made with this direct method, the errors are not yet small enough to make this method competitive with the traditional $`f_B^2B`$ method.
## 6 Semi-leptonic form factors, $`B\pi (\rho )\mathrm{}\nu `$.
The large $`b`$-quark mass presents a difficult computational problem. The cleanest lattice simulations are for ‘rest to rest’ i.e. both the initial and final meson at rest. Then the large $`B`$ mass forces $`q^2`$ ($`q=`$lepton 4-momentum$`=p_Bp_\pi `$) to be rather large. For some phenomenological applications the value of the form factor(s) are needed near $`q^20`$. This entails large extrapolations, introducing additional errors and model dependence. However, one important phenomenological application, namely deducing the mixing angle $`V_{ub}`$ from experimental data, only requires precise knowledge of the form factor at one value of $`q^2`$. This should work so long as the experiment has enough data so that even the differential rate around that region of $`q^2`$ is accurately determined, which is anticipated to be feasible at the $`B`$-factories. In this regard two recent approaches are noteworthy. Both of these efforts avoid the use of large extrapolations in $`q^2`$ and focus instead on accurate predictions in a limited region of $`q^2`$.
UKQCD focussed on near the end-point or the zero-recoil region where the lattice data tends to be cleanest. Furthermore, heavy quark symmetry also provides useful scaling relations in this region. Their result for the form factor $`f^+`$ as a function of $`q^2`$ is given below.
The FNAL group is also focussing on an approach towards the semi-leptonic form factors for $`B\pi \mathrm{}\nu `$, $`D\pi (K)\mathrm{}\nu `$ suitable for accurate determinations of mixing-angles in conjunction with high statistics data samples anticipated from experiments. The key idea here again is to concentrate directly on the differential decay spectrum in an interval with $`0.4<\stackrel{}{p}_\pi /\mathrm{GeV}<\mathrm{\hspace{0.33em}0.8}`$ thus avoiding the need for large extrapolation in $`q^2`$. The partial width over this interval can be computed on the lattice. From the experimental point of view this should have the advantage of using a range of $`q^2`$ wherein the decay rate is not small unlike near the end-point.
## 7 $`B_K`$
The kaon $`B`$-parameter, $`B_K`$, has been studied most extensively on the lattice. Two important limitations that lattice simulations still need to adequately address are SU(3) breaking (the ‘kaon’ on the lattice is a pseudoscalar made of degenerate quarks with $`m_{\mathrm{pseudoscalar}}m_K`$) and the quenched approximation. Both of these effects are expected to be rather small $`5`$–10% and we get $`\widehat{B}_K=.85\pm .13`$. This number is based on results of various groups , amongst which the one from JLQCD is the most precise. The error on the lattice number contains a guess-estimate of the SU(3) breaking and quenching errors. We note, in passing, that there are some preliminary indications that unquenching will increase or decrease $`B_K`$ by just a few percent. For now, we have not changed the central value of $`B_K`$ due to this effect; only the systematic errors are increased to reflect this possibility.
## 8 $`K2\pi `$ Decays and $`ϵ^{}/ϵ`$.
It was realized long ago that chiral perturbation theory (ChPT) can be used to simplify the problem so that $`\pi \pi |Q|K`$ can be obtained by computing on the lattice simpler entities: $`\pi |Q|K`$ and $`\mathrm{vac}|Q|K`$, where $`Q`$ is a 4-quark operator. The coefficients in this relation can be calculated using lowest order ChPT. Traditionally this strategy has been available for staggered fermions as they possess a remnant chiral symmetry. Since Wilson fermions explicitly break chiral symmetry this approach cannot be used with this discretization. The new development in this regard is that domain wall quarks (DWQ) are found to be quite practical for simulating QCD and possess excellent chiral behavior. Therefore, the $`K\pi `$ (and $`K`$vac) method for dealing with $`K2\pi `$ is amenable to this discretization as well.
Final state interactions (FSI) are of course a serious limitation of the $`K\pi `$ method. However, it should still be very instructive to quantify how well this concrete approximation works. Actually direct $`K2\pi `$ decays may also be amenable to the lattice due to an elegant application of the CPS symmetry. The restrictions of the Maiani-Testa theorem are bypassed here by working near the threshold. In any case the direct $`K2\pi `$ methods are computationally extremely intensive and for now there is not much to report based on these methods.
The credit for an extensive study of the $`K2\pi `$ in the $`K\pi `$ approach with staggered fermions goes to Pekurovsky and Kilcup (PK). The work presents a first study of lattice spacing dependence, finite size effects as well as quenching effects. Unfortunately PK used lattice weak coupling perturbation theory (LWCPT) to renormalize the operators and this seems to completely fail for staggered quarks in the case of LR operators such as $`Q_6`$ which are crucial for $`ϵ^{}/ϵ`$.
At $`\beta =6.0`$, PK observe a significant enhancement of the $`\mathrm{\Delta }I=1/2`$ channel over the 3/2; however at $`\beta =6.2`$, i.e. closer to the continuum limit (compared to $`\beta =6.0`$), they find that the central value of the enhancement weakens appreciably. While the large errors at $`\beta =6.2`$ do not allow a strong conclusion, more work is needed to unambiguously show that the $`\mathrm{\Delta }I=1/2`$ enhancement survives in the continuum limit.
PK’s calculation of $`ϵ^{}/ϵ`$ is seriously hampered as the one-loop LWCPT that they use for renormalization completely fails for $`Q_6`$. The perturbation theory corrections are several hundreds of percents showing extreme sensitivity to the renormalization scale as well as the quark mass (see Table 5). As PK themselves clearly emphasize their calculation of $`ϵ^{}/ϵ`$ is extremely fragile and not at all reliable.
It seems quite plausible that for the renormalization of the $`\mathrm{\Delta }S=1`$, 4-quark operators, as has also been known to some extent to be the case for quark bilinears, the breaking of flavor symmetry by the staggered approach is responsible for the failure of LWCPT. In any case, use of non-perturbative renormalization (NPR) improved actions and/or operators with the staggered approach are highly desirable in the context of $`K\pi \pi `$ decays.
Domain wall quarks are extremely attractive as at the expense of introducing a fictitious $`5^{th}`$ dimension they preserve the full SU($`N)\times `$SU($`N`$) chiral symmetries of the continuum theory in the limit of an infinite $`5^{th}`$ dimension. Quenched QCD numerical simulations showed that in practice for $`\beta >\mathrm{\hspace{0.33em}6.0}`$, 10–20 sites in the $`5^{th}`$ dimension may be sufficient to render very good chiral behavior. Early numerical studies also seem to indicate that the discretization errors are effectively $`0(a^2)`$; if substantiated this improved scaling behavior may off-set the extra cost of the $`5^{th}`$ dimension.
Calculation of $`K2\pi `$ and $`ϵ^{}/ϵ`$ in this method has been in progress for quite sometime. With DWQ, considerable progress has been made in non-perturbative renormalization of quark bilinears, $`\mathrm{\Delta }S=2`$ and $`\mathrm{\Delta }S=1`$ Hamiltonians and so far the method seems promising.
## Acknowledgments
I thank the organizers for a very interesting workshop. This work was supported in part by the U.S. DOE contract DE-AC02-98CH10886.
## References
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# Cyclotron effective masses in layered metals
## A Introduction
Quasi-two-dimesional metals such as the organic molecular crystals based on the BEDT-TTF molecule \[where BEDT-TTF is bis-(ethylenedithia-tetrathiafulvalene)\] and the layered perovskites Sr<sub>2</sub>RuO<sub>4</sub>, show properties which are consistent with a Fermi liquid description at low temperatures. Although transport properties of these materials show unconventional behaviour with temperature at high temperatures, at low temperatures (below about 20 K in the organics) the resistivity is quadratic with temperature, the thermopower is linear in temperature and a Drude peak is present in the optical conductivity. Furthermore, magnetic oscillations such as the de Haas - van Alphen effect is observed suggesting the presence of a well defined Fermi surface and quasiparticle excitations described by Fermi liquid theory. In order to understand the role of electron-electron interactions in these materials it is then necessary to quantify the strength of electron correlations and test how robust the Fermi liquid description is.
Cyclotron effective masses for the quasiparticles can be obtained from fitting the observed temperature dependence of the amplitude of magnetic oscillations to the Lifshitz-Kosevich form. The amplitude at a temperature $`T`$ is proportional to
$$R_T=\frac{X}{\mathrm{sinh}X}X=\frac{2\pi ^2k_BT}{\mathrm{}\omega _c^{}}$$
(1)
where $`\omega _c^{}=eB/m_c^{}`$ is the cyclotron frequency and $`m_c^{}`$ is the cyclotron effective mass, including many-body effects.
Typical values obtained for the cyclotron mass in these materials are in the range, $`m_c^{}/m_e17`$, (where $`m_e`$ is the free electron mass) suggesting the possibility that many-body effects may cause a significant enhancement of the effective mass. However, knowing $`m_c^{}/m_e`$ by itself is not sufficient to determine the size of many-body effects due to electron-electron and electron-phonon interactions. First, it is necessary to compute the cyclotron band mass, $`m_c`$, which takes into account the fact that electrons are not free but are moving in the presence of the periodic potential associated with the crystal lattice. Then, the ratio $`m_c^{}/m_c`$ can be used to estimate the importance of many-body effects. As we will see, estimates of $`m_c`$ deduced from band structure calculations, can vary by as much as a factor of three.
On the other hand, recent calculations of the transport properties of strongly correlated systems using dynamical mean-field theory to solve the Hubbard model on a frustrated hypercubic lattice, indicate that as the electronic correlations become stronger there is a clear crossover from a Fermi liquid at low temperatures to a “bad metal” with no quasiparticles at high temperatures. However, such a crossover and the associated signatures in transport properties (e.g., a peak in the temperature dependence of the thermopower and resistivity, and disappearance of the Drude peak in the optical conductivity) are only observed for sufficiently large values of the ratio: $`m_c^{}/m_c34`$. For smaller values, transport properties resemble the ones found in a nearly free-electron metal. Since the transport properties of the organic metals do show the signatures discussed above it is important to have accurate estimates of $`m^{}/m_c`$ in order to check the consistency of describing them as strongly correlated systems.
In this paper, we show that in a quasi-two-dimensional Fermi liquid there is a simple relation between the cyclotron mass and the density of states at the Fermi surface. This result, Eq. (7), holds for any dispersion relation for the quasiparticles. Using this relation, we compute the ratio of the cyclotron band masses associated with the $`\alpha `$ and $`\beta `$ orbits, $`m_c^\beta /m_c^\alpha `$, for a model band structure for $`\kappa `$-(BEDT-TTF)<sub>2</sub>X. The ratio is approximately 2, and varies by only about ten per cent even when the band structure parameters are varied significantly. This is in good agreement with values of $`m^\beta /m^\alpha `$, deduced from magnetic oscillation experiments, suggesting that the quasiparticle renormalisation factor does not vary significantly between different parts of the Fermi surface.
## B Different effective masses
We now briefly review several of the effective masses which can be defined for electrons or quasiparticles with a general dispersion relation $`ϵ(𝐤)`$.
Band mass tensor. This is defined as
$$m_{\nu \mu }^b\mathrm{}^2\left(\frac{^2ϵ(𝐤)}{k_\nu k_\mu }\right)^1$$
(2)
where $`\nu `$ and $`\mu `$ are Cartesian coordinates, and gives information of the band dispersion for any value of the electronic momentum. In particular, from the band mass tensor, the band dispersion of the electrons in all directions near the Fermi surface can be reconstructed.
Cyclotron mass. When a metal is in the presence of an external magnetic field the electrons undergo periodic orbits in both position and momentum space. The cyclotron frequency, $`\omega _c`$, associated with the periodic motion along these orbits on the Fermi surface is given by
$$\frac{1}{\omega _c}=\frac{\mathrm{}^2}{2\pi eB}\frac{dk}{(ϵ(𝐤))_{}}\frac{m_c}{eB}$$
(3)
where $`B`$ is the strength of the magnetic field and $`(ϵ(𝐤))_{}`$, is the gradient of the dispersion relation in the plane perpendicular to the field and the line integral is around the periodic orbit on the Fermi surface. The last relation has been used to define a cyclotron effective mass, $`m_c`$. Note that this effective mass involves an average of the dispersion relation along the periodic orbit. It determines the energy spacing of the Landau levels and can be extracted from the temperature dependence of the amplitude of magnetic oscillations, as discussed above.
Plasma frequencies. Reflectivity measurements can be used to determine the plasma frequency associated with collective oscillations of a charged Fermi liquid. Polarized light can be used to determine the anisotropy of these frequencies. For light polarized with the electric field in the $`\mu `$ direction in a metal with Fermi energy $`ϵ_F`$, the plasma frequency, $`\omega _{p\mu }`$, is given by,
$$\omega _{p\mu }^2=(\frac{e}{\pi \mathrm{}})^2d^3𝐤\frac{^2ϵ(𝐤)}{k_\mu ^2}\theta (ϵ_Fϵ_𝐤)\frac{ne^2}{m_{p\mu }}$$
(4)
where the integral runs over the first Brillouin zone and the last identity has been used to define an effective mass, $`m_{p\mu }`$, when $`n`$ is the total number of charge carriers. The above expression is derived from Lindhard’s dielectric function. Note that in contrast with the cyclotron mass in Eq. (3), which depends on electron states at the Fermi surface, the plasma mass includes all the occupied states, and not only those which are close to the Fermi energy. This is because the plasma oscillation is a collective process in which all the electrons participate.
For a parabolic dispersion relation, $`ϵ(𝐤)=\mathrm{}^2k^2/(2m_0)`$, all of the effective masses defined above will equal $`m_0`$. However, we stress that for a general dispersion relation they will not be equal and so caution is in order when trying to compare effective masses extracted from different measurements.
## C The cyclotron mass and the density of states
We know show how for a quasi-two-dimensional metal, the cyclotron mass defined by (3) is simply related to the density of states at the Fermi energy. First, following Ashcroft and Mermin, Eq. (3) can be rearranged to give
$$m_c=\frac{\mathrm{}^2}{2\pi }\frac{A(ϵ_F)}{ϵ_F}$$
(5)
where $`A(ϵ_F)`$ is the area of the cross section of the Fermi surface defined by the orbit described by an electron or hole, in the presence of a the magnetic field, $`B`$.
For a quasi-two-dimensional system with only one band that crosses the Fermi energy, and a magnetic field perpendicular to the layers, the area of the orbit is just the cross sectional area of the Fermi surface within a layer
$$A(ϵ_F)=4\pi ^2\underset{𝐤}{}\theta (ϵ_Fϵ(𝐤))$$
(6)
where k is the two-dimensional wavevector within a layer. Eq.(6) is just based on state counting and assumes that the interlayer dispersion can be neglected. Corrections due to a finite interlayer bandwidth will be of order $`t_c/ϵ_F`$ where $`t_c`$ is the interlayer hopping integral. For typical organic metals this ratio is less than 0.01. Taking the derivative of (6) with respect to $`ϵ_F`$ gives, for the cyclotron mass
$$m_c=2\pi \mathrm{}^2\rho _\sigma (ϵ_F)$$
(7)
where $`\rho _\sigma (ϵ_F)`$ is the density of states per spin at the Fermi energy, $`\rho _\sigma (ϵ_F)=_𝐤\delta (ϵ_Fϵ(𝐤))`$. We stress that this simple expression for the cyclotron band mass is only true for quasi-two-dimensional metals. In other cases, the reduction of the general expression (5) to (7) cannot be done. For example, for a three-dimensional system the area associated to an electron or hole orbit is not defined by Eq. (6), and, therefore, it is not possible to relate the cyclotron mass to the density of states at the Fermi energy. The result (7) was previously pointed out by Tamura et al. but its significance appears to have been completely overlooked. We will show below that as a consequence of Luttinger’s theorem it is also true in the presence of interactions.
For more general situations where the Fermi surface of the metal crosses several bands, the different cyclotron masses can be expressed in terms of the partial density of states associated with each of the bands. As an example, we will compute the band cyclotron masses for the $`\alpha `$ and $`\beta `$ orbits in the $`\kappa `$-(BEDT-TTF)<sub>2</sub>X family, for which several bands are present.
## D Model band structures
There are several approaches used for calculating the band structure of layered materials. Semi-empirical approaches such as the Hückel approximation use parametrized tight-binding Hamiltonians with parameters that are partially determined from experiment. In the case of $`\kappa `$-(BEDT-TTF)<sub>2</sub>X crystals, the effective tight-binding Hamiltonian which is used to model the interaction between the antibonding orbitals of BEDT-TTF dimers is
$$H=t_1\underset{ij}{}(c_i^{}c_j+h.c.)+t_3\underset{ik}{}(c_i^{}c_k+h.c.)+t_2\underset{il}{}(c_i^{}c_l+h.c.)$$
(8)
where $`c_i^{}`$, creates an electron in the antibonding orbital at site $`i`$ on a square lattice. $`t_1`$ and $`t_3`$ are nearest-neighbour hoppings, and $`t_2`$ is the next-nearest neighbour hopping amplitude along only one diagonal. The $`\kappa `$-(BEDT-TTF)<sub>2</sub>X materials have two dimers per unit cell and, because $`t_1`$ and $`t_3`$ can be slightly different the two dimers in each cell of the $`\kappa `$-(BEDT-TTF)<sub>2</sub>X materials are inequivalent. The relationship between the different hopping integrals and the geometrical arrangement of the BEDT-TTF molecules is shown in Fig. 1.
In Fig. 2 we show the stacking pattern for the $`\beta `$-(BEDT-TTF)<sub>2</sub>X family. In this case all the sites in the lattice are equivalent and there is only one dimer per unit cell.
If we diagonalize the Hamiltonian (8), we obtain the two dispersion relations
$$ϵ^\pm (𝐤)=t_2\mathrm{cos}(k_y)\pm \left(t_1^2+t_3^2+2t_1t_3\mathrm{cos}(k_x)\right)^{1/2}\mathrm{cos}(k_y/2)$$
(9)
The Fermi surface is shown in Fig. 3, for $`t_1t_3=0.05`$ and $`t_2=t_1`$. The $`\alpha `$ orbit is associated with the hole pocket in the Fermi surface and the unoccupied part of the lower band, $`ϵ(𝐤)`$, while the $`\beta `$ orbit (which occurs in large magnetic fields due to magnetic breakdown) contains parts from both the upper and lower band dispersions and corresponds to the outer orbit described with arrows in Fig. 3.
For the $`\beta `$-(BEDT-TTF)<sub>2</sub>X materials, due to the columnar stacking of the BEDT-TTF molecules and being $`t_1t_2t_3`$, there is only one dimer of BEDT-TTF molecules per site, and there is only one half-filled band which cuts the Fermi energy and is described by
$$ϵ(𝐤)=t_2\mathrm{cos}(k_y)+2t_1\mathrm{cos}(k_x/2)\mathrm{cos}(k_y/2)$$
(10)
For the $`\theta `$-(BEDT-TTF)<sub>2</sub>X materials, the geometrical arrangement is similar to that for $`\kappa `$-(BEDT-TTF)<sub>2</sub>X with each dimer replaced by a single BEDT-TTF molecule. It is then described by the dispersion relation (9) but the band is 3/4 filled.
We have evaluated the cyclotron band masses associated with the different orbits described along the Fermi surface for $`\kappa `$-(BEDT-TTF)<sub>2</sub>X. The area associated with the $`\alpha `$-orbit (see Fig. 3) is given by
$$A^\alpha (ϵ_F)=4\pi ^2\underset{𝐤}{}(1\theta (ϵ_Fϵ^{}(𝐤)))$$
(11)
and the cyclotron effective mass is, from Eq. (7)
$$m_c^\alpha =2\pi \mathrm{}^2\rho _\sigma ^{}(ϵ_F)$$
(12)
where $`\rho _\sigma ^{}(ϵ_F)`$, is the density of states per unit cell and spin associated with the $`ϵ^{}(𝐤)`$ band. Similarly, the area enclosed by the $`\beta `$-orbit is
$$A^\beta (ϵ_F)=4\pi ^2\underset{𝐤}{}(1\theta (ϵ_Fϵ^+(𝐤)))+4\pi ^2\underset{𝐤}{}(1\theta (ϵ_Fϵ^{}(𝐤)))$$
(13)
and the cyclotron mass is proportional to the total density of states:
$$m_c^\beta =2\pi \mathrm{}^2\rho _\sigma (ϵ_F)$$
(14)
where $`\rho _\sigma (ϵ_F)`$ is the total density of states per unit cell and spin. Note that a minus sign comes in the above expressions when we are considering the electron mass instead of the hole mass as $`m_e=m_h`$, where $`m_h`$ is the hole mass.
For the dispersion (9) with $`t_1=t_2=t_3=t`$, Ivanov, Yakushi, and Ugolkova have obtained analytical expressions for the density of states projected onto the upper and lower bands can be obtained. If all energies are in units of $`t`$, the total density of states per unit cell and spin is
$`\rho (3/2ϵ1)`$ $`=`$ $`{\displaystyle \frac{2}{\pi ^2q\sqrt{\tau }}}K({\displaystyle \frac{1}{q}})`$ (15)
$`\rho (1ϵ3)`$ $`=`$ $`{\displaystyle \frac{2}{\pi ^2\sqrt{\tau }}}K(q)`$ (16)
and, for the partial density of states associated with the lower band
$`\rho ^{}({\displaystyle \frac{3}{2}}ϵ1)`$ $`=`$ $`{\displaystyle \frac{2}{\pi ^2q\sqrt{\tau }}}K({\displaystyle \frac{1}{q}})`$ (18)
$`\rho ^{}(1ϵ1)`$ $`=`$ $`{\displaystyle \frac{2}{\pi ^2\sqrt{\tau }}}F(\mathrm{arcsin}\left({\displaystyle \frac{1}{2q}}\sqrt{{\displaystyle \frac{(5\tau ^2)(\tau +1)}{2}}}\right);q)`$ (19)
where $`q=\sqrt{1(\tau 1)^3(\tau +3)/(16\tau )}`$ with $`\tau =\sqrt{2ϵ+3}`$. $`K`$ and $`F`$ are the complete elliptic integral and the elliptic integral of the first kind, respectively.
From the above expressions and Eq. (7) we obtain the following cyclotron masses: $`m_c^\beta /m_e=0.23/t`$ and $`m_c^\alpha /m_e=0.11/t`$, with $`t`$ given in eV and we have used the intralayer unit cell area of $`A=104.\AA ^2`$. This gives $`m_c^\beta /m_c^\alpha =2`$ and it turns out that this ratio is relatively insensitive to variations in the band structure parameters. We have relaxed the condition on the hopping integrals $`t_1=t_2=t_3`$, and, we have numerically evaluated the partial density of states instead of using eq.(LABEL:Ivan1) and eq.(LABEL:Ivan2). The ratio of the cyclotron masses obtained from the effective dimer model for fixed $`t_1=t_3`$ but different values of $`t_2/t_1`$ is, $`m_c^\beta /m_c^\alpha `$ = 2.4, 2.2, and 2.0, for $`t_2/t_1=0.5,0.7,1.0`$, respectively.
In order to have a realistic description of the layered materials we use the hopping amplitudes obtained from quantum chemistry calculations using the Hückel approximation and, in some cases, results obtained from first-principle calculations. The hoppings of the effective dimer model, for which $`t_1=t_3`$ and $`t_1t_2`$, are given in Table I. A more detailed discussion of this model and the relationship between $`t_1`$ and $`t_2`$ and the intermolecular hoppings calculated in the Hückel approximation can be found in Reference . A minor point is that if we denote the Coulomb repulsion in each molecule by $`U_0`$, and the hopping amplitude between the molecules within one dimer by $`t_b`$, for $`U_0>>4t_b`$ (strongly correlated case), the hopping amplitudes should be corrected by a factor of $`1/\sqrt{2}`$ with respect to the ones obtained from the Hückel calculation. However, in the case $`U_04t_b`$, this factor is 0.92 and the effect of correlations to the matrix elements is small. Different calculations suggest that the ratio $`U_0/4t_b1`$, so that in Table I we multiply all the bare hoppings by 0.92.
In Table I, we also give the cyclotron masses obtained from eq.(7), where the density of states has been computed numerically for the different hoppings. It can be seen that the calculated cyclotron band masses are sensitive to the parameters and the values deduced from the parameters calculated by different groups for the same material can vary significantly. However, the calculated ratio, $`m_c^\beta /m_c^\alpha `$, is relatively insensitive to the parameters.
The band structures of the (BEDT-TTF)<sub>2</sub>X family have been calculated by several different techniques and some of the results for the density of states at the Fermi energy are compared in Table II. The Hückel method is the simplest and only considers the $`\pi `$ orbitals and neglects all $`\sigma `$ orbitals. The overlap integrals that are calculated are all scaled by some empirical parameter and then used as hopping integrals in a tight-binding band structure. It is generally acknowledged that this method gives a good qualitative description of electronic properties (such as the symmetry and ordering of states) but cannot give a quantitative description of electronic properties.
The extended Hückel method treats both $`\pi `$ and $`\sigma `$ orbitals. Although it is more quantitatively reliable than the Hückel approximation it still does not give a completely quantitative description of organic molecules. It has been used to calculate the band structure of a wide range of organic metals by Whangbo and coworkers.
The energy levels for a pair of BEDT-TTF dimers with the same geometrical arrangement as in $`\kappa `$-(BEDT-TTF)<sub>2</sub>Cu\[N(CN)<sub>2</sub>\]Br have been calculated by an ab initio method. The tight-binding parameters for a Hubbard model for the dimers is then evaluated by fitting the energy levels to the ab initio values. The resulting parameters are similar to those obtained by an extended Hückel calculation for the dimer pair. But the resulting density of states is more than twice the results of extended Hückel for the solid.
The most reliable method of calculating band structures is generally considered to be ab initio methods based on the local density approximation (LDA). Nevertheless, different groups still often obtain quite different results. For example, values obtained for the density of states at the Fermi energy in the fullerene metal, K<sub>3</sub>C<sub>60</sub>, differ by as much as fifty per cent. (Extended Hückel calculations do fall in this range). Due to the large number of atoms in a unit cell only a few ab initio calculations have been attempted for the (BEDT-TTF)<sub>2</sub>X materials.
Results for the density of states (and the corresponding cyclotron masses) obtained using the three methods are shown in Table II. Note the large variation in results for each of the materials. In particular, the Hückel method gives masses that are two to five times larger than those obtained by the other more sophisticated methods.
## E The cyclotron mass in the presence of interactions
The above treatment neglected the effect of interactions between the electrons. We now show that Eq. (7) has a natural generalization in the case of a Fermi liquid. The one-electron Green’s function in a general interacting electron system is
$$G(𝐤,\omega +i\eta )=\frac{1}{\omega +i\eta ϵ(𝐤)\mathrm{\Sigma }(𝐤,\omega )}$$
(21)
in momentum space, where $`\mathrm{\Sigma }(𝐤,\omega )`$ is the electron self-energy. In a Fermi liquid, near the quasiparticle poles, the Green’s function can be rewritten as
$$G(𝐤,\omega )=\frac{Z_𝐤}{\omega \stackrel{~}{ϵ}(𝐤)}$$
(22)
where $`\stackrel{~}{ϵ}(𝐤)`$ is the quasiparticle energy and $`Z_𝐤=\frac{1}{1\frac{\mathrm{\Sigma }(𝐤,\omega )}{\omega }|_{\omega =\stackrel{~}{ϵ}(𝐤)}}`$ is the residue at the quasiparticle pole. Note that the above expression is true for a Fermi liquid, and for electrons with momentum close to the Fermi surface for which Im$`\mathrm{\Sigma }(𝐤𝐤_𝐅,\omega ϵ_F)0`$. The spectral density is then given by
$$A(𝐤,\omega )=\frac{1}{\pi }\mathrm{Im}G(\omega +i\eta )=\frac{\delta (\omega \stackrel{~}{ϵ}(𝐤))}{1\frac{\mathrm{\Sigma }(𝐤,\omega )}{\omega }|_{\omega =\stackrel{~}{ϵ}(𝐤)}}$$
(23)
Thus, the quasiparticle density of states at the Fermi energy is
$$\stackrel{~}{\rho }(\stackrel{~}{ϵ}_F)=\underset{𝐤}{}\delta (\stackrel{~}{ϵ}_F\stackrel{~}{ϵ}(𝐤))=\underset{𝐤}{}(1\frac{\mathrm{\Sigma }(𝐤,\omega )}{\omega }|_{\omega =\stackrel{~}{ϵ}_F})A(𝐤,\stackrel{~}{ϵ}_F)$$
(24)
Müller-Hartman showed that, if the self energy is independent of momentum, then at zero temperature, $`_𝐤A(𝐤,\stackrel{~}{ϵ}_F)=\rho (ϵ_F)`$, the non-interacting density of states at the Fermi energy. So in this case, $`\stackrel{~}{\rho }(\stackrel{~}{ϵ}_F)=\rho (ϵ_F)/Z`$. Note that the quasiparticle density of states is always enhanced because for a Fermi liquid $`\frac{\mathrm{\Sigma }(𝐤,\omega )}{\omega }|_{\omega =\stackrel{~}{ϵ}_F}<0`$.
Some time ago, Luttinger showed that in an interacting system with Fermi liquid properties, the results of Lifshitz and Kosevich still describe the de Haas van Alphen oscillations provided that the relevant quasiparticle quantities are used. Thus, (5) is replaced by $`m_c^{}=\frac{\stackrel{~}{A}}{\stackrel{~}{ϵ}_F}`$ where a tilde denotes renormalised quantities. In a quasi-two-dimensional Fermi liquid, the area enclosed by the orbits of the quasiparticles is
$$\stackrel{~}{A}(\stackrel{~}{ϵ}_F)=4\pi ^2\underset{𝐤}{}\theta (\stackrel{~}{ϵ}_F\stackrel{~}{ϵ}(𝐤))$$
(25)
and so, we find that the cyclotron effective mass is
$`m_c^{}=2\pi \mathrm{}^2{\displaystyle \underset{𝐤}{}}\delta (\stackrel{~}{ϵ}_F\stackrel{~}{ϵ}(𝐤))=2\pi \mathrm{}^2\stackrel{~}{\rho }_\sigma (\stackrel{~}{ϵ}_F)`$ (26)
Again, equations (26) and (24) show the cyclotron mass enhancement produced by the factor appearing in eq.(24). The same enhancement also appears in the specific heat coefficient. A further simplification is obtained for the case of a momentum independent self-energy, as then the cyclotron effective masses reduce to
$$m_c^{}=2\pi ^2\mathrm{}^2\rho (ϵ_F)/Z=m_c/Z$$
(27)
where $`Z`$ is the quasiparticle weight, which, in terms of the self-energy, is: $`Z=(1\frac{\mathrm{\Sigma }(\omega )}{\omega }|_{\omega =\stackrel{~}{ϵ}_F})^1`$. In this case, the ratios of the cyclotron effective masses associated with the quasiparticles moving along different orbits, $`m_c^\beta /m_c^\alpha `$, should be the same as the ratios associated with the non-interacting system, $`m_c^\beta /m_c^\alpha `$.
A partial test of the momentum independent self energy is provided by comparing the measured ratios of the renormalized cyclotron masses in different orbits with the cyclotron band mass ratios. This is done in Table I. The relative consistency between the observed values of this ratio and the band structure values suggests that if there are sizeable renormalizations due to many-body effects then these renormalizations are not significantly different on the different parts of the Fermi surface. However, this consistency is only a necessary condition but not sufficient for having a momentum independent self energy, as cyclotron masses include averages over the Fermi surface and, therefore, cancellations of contributions from different parts of the Fermi surface may occur.
Furthermore, in Reference the effective masses for $`\kappa `$-(BEDT-TTF)<sub>2</sub>Cu(SCN)<sub>2</sub> were measured as the pressure was increased from 1 bar to 20 kbar. $`m_c^\beta /m_e`$ decreased from 6.5 $`\pm `$ 0.1 at 1 bar to 2.7 $`\pm `$ 0.1 at 16.3 kbar. $`m_c^\alpha /m_e`$ decreased from 3.5 $`\pm `$ 0.1 at 1 bar to 1.4 $`\pm `$ 0.1 at 16.3 kbar. However, the ratio $`m_c^\beta /m_c^\alpha `$ has a constant value of 1.9 within error.
## F Specific heat
Measurements of the electronic specific heat in the (BEDT-TTF)<sub>2</sub>X crystals and Sr<sub>2</sub>RuO<sub>4</sub> show a linear temperature dependence at low temperatures, consistent with a Fermi liquid description. The corresponding specific heat coefficient $`\gamma `$ is given in Table III for some of these materials. This coefficient is related to the quasiparticle density of states at the Fermi energy, $`\stackrel{~}{\rho }(\stackrel{~}{ϵ}_F)`$, (see Eq. (24)) by
$$\gamma =\frac{2\pi ^2k_B^2}{3}\stackrel{~}{\rho }(\stackrel{~}{ϵ}_F)$$
(28)
Since the quasiparticle density of states is also related to the cyclotron effective mass by Eq. (26) the measured specific heat coefficient can be used to calculate a corresponding cyclotron effective mass. This has been done in Table III for a range of organic materials. The values obtained for $`m_c^\beta /m_c`$ from specific heat measurements agree for $`\kappa `$-(BEDT-TTF)<sub>2</sub>I<sub>3</sub> and $`\beta `$-(BEDT-TTF)<sub>2</sub>I<sub>3</sub> but not for the materials with copper in the anion. Since this comparison does provide a quantitative test of a Fermi liquid description further careful measurements are justified, particularly on a wider range of materials.
Such a comparison was also done recently for Sr<sub>2</sub>RuO<sub>4</sub> in Reference , where relation (26) was implicitly assumed, presumably based on its validity for a parabolic dispersion relation. Our work provides a rigorous justification for this comparison. In Sr<sub>2</sub>RuO<sub>4</sub> there are three distinct Fermi surfaces and the associated cyclotron masses deduced from de Haas van Alphen oscillations were $`m_c^{}/m_e`$=3.4, 7.5 and 14.6 for the $`\alpha `$, $`\beta `$ and $`\gamma `$ orbits, respectively.From the above discussion, it follows that the specific heat coefficient of Sr<sub>2</sub>RuO<sub>4</sub> is related to the effective masses by
$$\gamma =\frac{\pi k_B^2}{3\mathrm{}^2}(m_c^\alpha +m_c^\beta +m_c^\gamma )$$
(29)
which comes from the fact that the total density of states is just the sum of the density of states of the different Fermi surfaces. Evaluating (29) we obtain a specific heat coefficient of 36.7 mJ/(K<sup>2</sup> mol), which agrees with the measured value of 37.4 mJ/(K<sup>2</sup> mol).
## G Conclusions
We now summarize our results and their implications. First, it was shown that in a quasi-two-dimensional metal in which the dispersion perpendicular to the layers can be neglected, the cyclotron effective mass for a particular orbit in a general band structure is simply related to the density of states at the Fermi energy associated with the relevant band. Second, it was shown that, due to Luttinger’s results for a Fermi liquid, a similar relationship holds in the presence of interactions.
These results have a number of general applications to layered metals which have Fermi liquid properties at low temperatures.
(i) In order to evaluate the effective mass from band structure it is not necessary to numerically evaluate the derivative in (5), as has been done previously by a number of authors. Instead (7) can be used together with the density of states at the Fermi energy. This eliminates the need to perform the cumbersome task of repeating the band structure calculations for many different Fermi energies.
(ii) We found that for model band structures describing the family $`\kappa `$-(BEDT-TTF)<sub>2</sub>X, the ratio of the effective mass for the $`\beta `$-orbit to the mass for the $`\alpha `$ orbit is fairly insensitive to the details of the band structure, having a value close to two.
(iii) Our results imply that a quantitative test of the Fermi liquid description of a layered metal is to compare measurements of the cyclotron effective mass to the linear coefficient in the specific heat.
(iv) The agreement between the ratio of the different measured cyclotron masses and the ratio calculated from band structure, suggests a momentum independent self-energy, although other experimental probes such as polarized Raman scattering, photoemission spectra, and angular dependent magnetoresistance oscillations are needed before making any definitive conclusion.
Based on comparison with a wide range of materials we conclude the following. First, the effective masses deduced from magnetic oscillations and specific heat are consistent for Sr<sub>2</sub>RuO<sub>4</sub> and for two out of four of the organic materials considered. For three out of four of the organic materials for which data is avalailable the measured ratio $`m_c^\beta /m_c^\alpha `$ is consistent with the band structure ratio $`m_c^\beta /m_c^\alpha `$. Furthermore, for the $`\kappa `$-(BEDT-TTF)<sub>2</sub> Cu(NCS)<sub>2</sub> this ratio does not change under pressure while the individual effective masses decrease by a factor of 2.5. This suggests that the self energy does not vary significantly over the different parts of the Fermi surface. We also note that the significant variation of the effective masses with pressure cannot be explained in terms of band structure; it predicts a small variation with pressure.
A comparison of the results of band structure calculations using a range of methods found that they produced a large range in values for the density of states (and thus the effective masses). The Hückel method has often been used to estimate the hopping integrals in tight binding band structures (as in Table I). It is less sophisticated than the extended Hückel method which in turn is less sophisticated than ab initio methods based on the local density approximation. We suggest that the Hückel method is producing hopping integrals which are too small by a factor of two to four. The best strategy to evaluate these integrals would be to fit a LDA band structure to a tight binding dispersion, such as (9). Such an approach was recently taken for Sr<sub>2</sub>RuO<sub>4</sub>.
We now come back to the central question of this paper: are the layered metals we have considered strongly correlated? A definitive answer is not possible because of the large variation in values for the band cyclotron masses that have been calculated by different band structure methods. However, we suggest that due to their greater sophistication, the local density approximation and extended Hückel approximation calculations give the most reliable values. We suggest that the appropriate values for the band cyclotron masses are those calculated by the local density approximation and extended Hückel approximation. The mass ratios given in Table II then imply that $`m_c^\beta /m_c2.54`$, suggesting appreciable quasiparticle renormalization due to many-body effects. This is consistent with the strong temperature dependence of the transport properties, discussed in detail in Reference .
###### Acknowledgements.
We thank R. McKinnon, N. Harrison, J. S. Brooks, J. Wosnitza, and E. Canadell for helpful discussions. This work was supported by the Australian Research Council.
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# Dynamical content of quantum diffeomorphisms in two-dimensional quantum gravity
## 1 Introduction
Our main goal is the construction of a space-time notion just from symmetry principles. Therefore, we consider space-time as a derived, or secondary, object in our theory. Assuming that the nature of the world is a quantum one, we want to start from the very beginning with a quantum theory in which, perhaps, there is not an explicit notion of space-time. We understand by such a quantum theory a unitary and irreducible representation of the algebra of physical operators and the space-time structure itself, in case such an structure makes any sense, should be built out of this quantum theory.
In order to fill the gap between the starting point (quantum theory) and the objective (space-time construction), we present a very simple model that we interpret as a quantum gravity (see ), in fact a two-dimensional one, even though the link with standard two-dimensional quantum gravity is not completely obvious.
The technical tool we are going to use is a Group Approach to Quantization (GAQ) (see and references therein), which starts from a given group of symmetries and constructs a quantum theory in the sense mentioned above. Among the elements of this approach, one is specially important for this work as it will play the role of physical guide in our search of space-time. This element is group (pseudo-)cohomology, which is a mathematical concept related to central extensions of the considered group. Its importance follows from the fact that it allows the classification of the physical operators in two sets:
* Dynamical operators, which appear in conjugated pairs, and give a central term in their commutators: $`[\text{ . },\text{ . }]1`$.
* Kinematical operators, which, broadly speaking, do not give this central term (the distinction between dynamical and kinematical operators can be rigorously characterised in terms of a pre-symplectic form constructed from the group two-cocycle ). Space-time should appear within this second set which, in the language of GAQ, is called characteristic subalgebra and denoted by $`𝒢_\mathrm{\Theta }`$.
## 2 The mathematical model
The first step in the construction of the model consists in selecting an appropiate starting group. We are choosing the Virasoro group, in abstract terms, for such a group. The reasons for this are, on the one hand, that it is a group simple enough in order to be handled and, at the same time, rich enough for giving non-trivial answers. On the other hand, Polyakov’s action for $`2D`$ gravity can be derived from it ().
The Virasoro algebra is defined ($`h=\frac{cc^{}}{24}`$) by:
$`[L_n,L_m]=(nm)L_{n+m}+{\displaystyle \frac{1}{12}}(cn^3c^{}n)\delta _{n,m}.`$
A formal group law ($`l^{\prime \prime m}=F^m(l^n,l^n)`$) can be derived from this algebra, and from this group law the rest of the elements of GAQ (quantization 1-form, left- and right-invariant infinitesimal generators…) can be computed. As we said before, we are particularly interested in the characteristic subalgebra, the kinematical degrees of freedom. Looking at the commutation relations, we find the two posible cases:
$`\text{i)}{\displaystyle \frac{c^{}}{c}}r^2,rZ,𝒢_\mathrm{\Theta }`$ $`=`$ $`\stackrel{~}{X}_{l^o}^L`$
$`\text{ii)}{\displaystyle \frac{c^{}}{c}}=r^2,rZ,𝒢_\mathrm{\Theta }`$ $`=`$ $`\stackrel{~}{X}_{l^r}^L,\stackrel{~}{X}_{l^o}^L,\stackrel{~}{X}_{l^r}^L.`$
Since we are looking for a space-time with one time dimension and at least one spatial dimension, the first case is excluded. Therefore, our first conclusion is that in order to find a space-time, we must fall in the second critical case in which $`𝒢_\mathrm{\Theta }`$ closes an $`sl(2,R)`$ algebra. In fact, unitarity imposes $`r=1`$, so that we have $`c=c^{}`$, which is the only case we consider from now on.
The quantum representation is achieved by taking the set of $`U(1)`$complex functions defined on the group as the Hilbert space, and the group acting on it via the regular representation. The main problem then is that this representation is highly reducible. In principle, this can be solved by taking advantage of the trivial commutation among left- ($`\stackrel{~}{X}_{l^n}^L`$) and right-invariant ($`\stackrel{~}{X}_{l^n}^R`$) vector fields. This allows us to implement the representation with one set of vector fields (for instance, the right-invariant ones: $`L_n=i(\stackrel{~}{X}_{l^n}^R)`$) and to reduce the representation in a consistent way by using the other set, via the introduction of polarization equations: $`\stackrel{~}{X}_{l^{n1}}^L=0`$.
However, in the general case for the Virasoro group, the resulting representation is still reducible. We have two solutions to this problem which prove to be equivalent. The first one benefits from the existence of a vacuum state and consists in taking the orbit of the group through this vacuum. The resulting representation is an irreducible highest-weight representation, characterised by: $`L_n0=0,(n1)`$ and $`\mathrm{\Psi }=L_{n_j}\mathrm{}L_{n_1}0,(n2)`$. The second solution, which is more natural in the framework of GAQ, consists in imposing further polarizations conditions using higher-order differential operators. This makes pseudo-differential operators come on the scene, something which is a source of technical difficulties, however. For the sake of simplicity we are not going to consider those cases in which higher-order polarizations appear. The equivalence between the two solutions was proved in .
With regards to unitarity, the values of $`c`$ and $`c^{}`$ that make unitary the representation are :
* $`c1`$, with $`\frac{cc^{}}{24}0`$.
* $`0<c<1`$ with: $`c=1\frac{6}{m(m+1)}`$ and $`\frac{cc^{}}{24}=\frac{[(m+1)rms]^21}{4m(m+1)}`$, where $`1srm1`$, and $`m,r,s`$ integers with $`m2`$.
But the discrete values of $`c`$ and $`c^{}`$ for $`0<c<1`$ are precisely those cases related with higher-order polarizations and are, therefore, disregarded.
Thus, what we have by now, is a representation of the Virasoro algebra in which $`c=c^{}`$ (in order to find a space-time) and $`c>1`$ (to have unitarity). Furthermore, we have classified the operators in two sets: space-time operators ($`L_1,L_0,L_1`$), and dynamical operators ($`L_n,n2`$), which we shall refer to as gravity operators in the sequel. This way, the starting point of the work, the quantum theory, is constructed and now we try to accomplish our main objective by identidying a space-time structure.
Firstly, we consider the reduction of the Virasoro Hilbert space, $`_{(c,c)}`$, under the kinematical $`SL(2,R)`$ subgroup, obtaining:
$`_{(c,c)}={\displaystyle \underset{N}{}}(D^{(N)}D^{(N1)})R_S^{(N)}`$
where,
* $`R_S^{(N)}`$ is a maximal-weight irreducible representation of $`SL(2,R)`$ with Casimir $`N(N`$$`1)`$. We denote the states in this representation by $`N,n`$.
* $`D^{(N)}`$ is the dimension of the Virasoro level $`N`$.
* The different representations $`R_S^{(N)}`$ are orthogonal. This will be important in the physical interpretation because it permits a standard quantum mechanical interpretation.
Once the representation has been reduced, we proceed to associate a space-time with each $`SL(2,R)`$ representation in the model; that is, as we have an infinite number of $`SL(2,R)`$ representations, an infinite number of space-times are realised simultaneosly in our model.
In order to make this association more concrete, we take an specific $`SL(2,R)`$ representation and construct a $`C^{}`$-algebra by considering all the products of the wave functions in the representation. At this point, we can apply a theorem by Gelfand and Naimark , which allows the reconstruction of a manifold from the $`C^{}`$-algebra. The problem with this approach is that the mentioned theorem is a rather abstract tool and the identification of the actual manifold under consideration is a difficult task. Fortunately, we can look at the problem in another way. For this, we consider again an isolated $`SL(2,R)`$ representation and notice that we know another system with the same Hilbert space but for which the configuration space is explicitly known. This system is a particle moving on a two-dimensional AdS space-time. Thus, we associate with each $`SL(2,R)`$ representation a one-sheet hyperboloid. In fact, this is what one expects to find from a $`SL(2,R)`$ group, which is isomorphic to AdS group in two dimensions when imposing the Casimir constraint.
Before entering into the physical interpretation, let us make some brief comments on the mathematical model we have just introduced. Firstly, in the framework of GAQ, it is natural to assign dimensions to the vector fields. In our case we fix $`[L_n]=(Length)^1`$. The commutation relations then imply $`[c]=Length`$ and $`[n]=[c^{}]=(Length)^1`$. But this represents a problem when trying to interpret expressions like $`c>1`$, since we need a scale. To solve this problem we introduce a constant $`a`$, such that $`[a]=Length`$, and redefine $`n\frac{n}{a}`$. Thus, the redefined integers are dimensionless.
Another important point is that the Virasoro algebra has a natural notion of classical limit which is obtained by making $`c\mathrm{}`$, or in redefined terms, $`\frac{c}{a}\mathrm{}`$. In this case, the constant $`\frac{a}{c}`$ behaves as a parameter of a perturbative series, playing the role of the Planck constant in our model. This suggests to redefine the Virasoro generators: $`H_n\frac{a}{c}L_n`$.
The previous introduction of a hyperboloid associated with a given representation does not provide a metric structure for space-time. The only natural metric we can find inside the model is the one induced from the Killing metric of $`SL(2,R)`$, thus providing an AdS space-time. In order to fully determine it, we have to introduce a scale: the radius of the hyperboloid. It is defined from the Casimir in terms of the redefined Virasoro generators:
$`{\displaystyle \frac{1}{R^2}}H_{0}^{}{}_{}{}^{2}{\displaystyle \frac{1}{2}}(H_1H_1+H_1H_1)=`$
$`=({\displaystyle \frac{a}{c}})^2{\displaystyle \frac{N(N1)}{a^2}}={\displaystyle \frac{N(N1)}{c^2}}.`$
## 3 Physical interpretation
Now we can provide a physical interpretation. As we have stressed, we associate an AdS<sub>2</sub> space-time of radius $`R=\frac{c}{\sqrt{N(N1)}}`$ with each $`SL(2,R)`$ representation.
Given an specific $`SL(2,R)`$ representation, we interpret each vector in it ($`N,n`$) as a state of the corresponding space-time. Since we are dealing with maximal-weight representations, a tower of states of the space-time is found, where the maximal-weight vector plays the role of space-time ground state. $`H_0`$ is interpreted as the energy and $`H_1`$ ($`H_1`$) as raising (lowering) space-time operators. When we consider the dynamical operators, $`H_{n2}`$, we notice that they do not leave invariant the $`SL(2,R)`$ representations and therefore they produce the effect of mixing the different space-times.
The whole picture provided by the model is as follows: we understand by Universe the entire ensemble of different space-times which are realised at the same time. A state of the Universe is a particular state in the Virasoro representation, that is, an specific linear combination of states in different $`SL(2,R)`$ representations or, in other words, a quantum superposition of different space-times. The coefficients of the linear combination define a weight distribution of space-times. The question about the radius of the Universe makes no real sense, since we have space-times of different radii. The meaningful question is about the probability for the Universe to have a certain radius, and then it is essential the orthogonality of the $`SL(2,R)`$ representations, which allows the definition of orthogonal projectors. Finally, the effect of gravity is that of mixing the different space-times, thus changing the weight distribution of space-times.
As far as the classical limit is concerned, we can consider the limit $`c\mathrm{}`$ and find that for every space-time the energy of the ground state tends to zero ($`Energy(N,0)=\frac{N}{c}0`$), and the radius to infinity ($`R=\frac{c}{\sqrt{N(N1)}}\mathrm{}`$). There is no physical way to distinguish between the different space-times in this limit and it makes sense, accordingly, to identify them. In order to do that, we define the equivalence relation, $`N,n`$ $`N^{},n`$, and take the quotient $`_{(c,c)}/`$. The problem with this quotient is that the dynamical operators are ill-behaved in it; in fact, they are multivalued. To avoid this trouble, we impose to these operators to act trivially on the states (this is only consistent with the commutation relations in the classical limit, as pointed out by S. Carlip after the talk): $`H_{n2}\mathrm{\Psi }=0`$.
Hitherto, we have considered the Virasoro group as an abstract one, but at this point we can look at it as a diffeomorphism group, so that the previous constraints are precisely the classical diffeomorphisms constraints. What we find is that diffeomorphism invariance is only recovered in the classical limit, while diffeomorphisms have a dynamical content at the quantum level.
## 4 Conclusions
* We have constructed a space-time notion from a quantum theory but only for the critical value of the Virasoro anomaly ($`c=c^{}`$). Out of this value space-time makes no sense, even though the quantum theory does exist.
* The model presents the Universe as a quantum superposition of different space-times which are mixed by gravity modes.
* We have implemented a model which exhibits a quantum breakdown of diffeomorphism invariance through the appearance of physical (non-gauge) degrees of freedom through an anomaly process. General covariance is recovered only in the classical limit.
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# Pomeron intercepts at colliders
## 1 Forward jet cross-section at HERA
The study of forward jets at colliders is considered as the milestone of QCD studies at high energies, since it provides a direct way of testing the perturbative resummations of soft gluon radiation. More precisely, the study of one forward jet (w.r.t. the proton) in an electron-proton collider seems to be a good candidate to test the energy dependence of hard QCD cross-sections. It is similar to the previous proposal of studying two jets separated by a large rapidity interval in hadronic colliders , for which only preliminary results are available . This test is also possible in $`\gamma ^{}`$-$`\gamma ^{}`$ scattering but here the statistics and the energy range are still insufficient to get a reliable determination of the physical parameters for hard QCD cross-sections. Indeed, the proposed (and favored for the moment being) set-up is to consider jets with transverse momentum $`k_T`$ of the order of the photon virtuality $`Q`$ allowing to damp the QCD evolution as a function of $`k_T`$ (DGLAP evolution ) in favor of the evolution in energy at fixed $`k_T`$ (BFKL evolution ).
In contrast to full Monte-Carlo studies we want to focus on the jet cross section $`d\sigma /dx`$ observable itself, by a consistent treatment of the experimental cuts and minimizing the uncertainties for that particular observable. Let us remark that our approach is not intended to provide a substitution to the other methods, since the Monte-Carlo simulations have the great merit of making a set of predictions for various observables. Hence, our method has to be considered as complementary to the others and dedicated to a better determination of the effective Pomeron intercept using the $`d\sigma /dx`$ data. As we shall see, it will fix more precisely this parameter, but it will leave less constrained other interesting parameters, such as the cross-section normalization.
The cross-section for forward jet production at HERA in the dipole model reads :
$`{\displaystyle \frac{d^{(4)}\sigma }{dxdQ^2dx_Jdk_T^2d\mathrm{\Phi }}}`$ $`=`$ $`{\displaystyle \frac{\pi N_C\alpha ^2\alpha _S(k_T^2)}{Q^4k_T^2}}f_{eff}(x,\mu _f^2)\mathrm{\Sigma }e_Q^2{\displaystyle _{\frac{1}{2}i\mathrm{}}^{\frac{1}{2}+i\mathrm{}}}{\displaystyle \frac{d\gamma }{2i\pi }}\left({\displaystyle \frac{Q^2}{k_T^2}}\right)^\gamma \times `$ (1.1)
$`\times `$ $`\mathrm{exp}\{ϵ(\gamma ,0)Y\}\left[{\displaystyle \frac{h_T(\gamma )+h_L(\gamma )}{\gamma }}(1y)+{\displaystyle \frac{h_T(\gamma )}{\gamma }}{\displaystyle \frac{y^2}{2}}\right]`$
where
$`Y`$ $`=`$ $`\mathrm{ln}{\displaystyle \frac{x_J}{x}}`$ (1.2)
$`ϵ(\gamma ,p)`$ $`=`$ $`\overline{\alpha }\left[2\psi (1)\psi (p+1\gamma )\psi (p+\gamma )\right]`$ (1.3)
$`f_{eff}(x,\mu _f^2)`$ $`=`$ $`G(x,\mu _f^2)+{\displaystyle \frac{4}{9}}\mathrm{\Sigma }(Q_f+\overline{Q_f})`$ (1.4)
$`\mu _f^2`$ $``$ $`k_T^2,`$ (1.5)
are, respectively, $`Y`$ the rapidity interval between the photon probe and the jet, $`ϵ(\gamma ,p)`$ the BFKL kernel eigenvalues, $`f_{eff}`$ the effective structure function combination, and $`\mu _f`$ the corresponding factorization scale. The main BFKL parameter is $`\overline{\alpha },`$ which is the (fixed) value of the effective strong coupling constant in LO-BFKL formulae. Note that we gave the BFKL formula not including the azimuthal dependence as we will stick to the azimuth-independent contribution with the dominant $`\mathrm{exp}\{ϵ(\gamma ,0)Y\}`$ factor.
The so-called “impact factors”
$`\left(\begin{array}{c}h_T\\ h_L\end{array}\right)={\displaystyle \frac{\alpha _S(k_T^2)}{3\pi \gamma }}{\displaystyle \frac{(\mathrm{\Gamma }(1\gamma )\mathrm{\Gamma }(1+\gamma ))^3}{\mathrm{\Gamma }(22\gamma )\mathrm{\Gamma }(2+2\gamma )}}{\displaystyle \frac{1}{1\frac{2}{3}\gamma }}\left(\begin{array}{c}(1+\gamma )(1\frac{\gamma }{2})\\ \gamma (1\gamma )\end{array}\right),`$ (1.10)
are obtained from the $`k_T`$ factorization properties of the coupling of the BFKL amplitudes to external hard probes. The same factors can be related to the photon wave functions within the equivalent context of the QCD dipole model .
Our goal is to compare as directly as possible the theoretical parametrization (1.1) to the data which are collected in experiments . The crucial point is how to take into account the experimentally defined kinematic cuts .
The main problem to solve is to investigate the effect of these cuts on the determination of the integration variables leading to a prediction for $`d\sigma /dx`$ from the given theoretical formula for $`d^{(4)}\sigma `$ as given in formula (1.1). The effect is expected to appear as bin-per-bin correction factors to be multiplied to the theoretical cross-sections for average values of the kinematic variables for a given $`x`$-bin before comparing to data (e.g. fitting the cross-sections) .
The experimental correction factors have been determined using a toy Monte-Carlo designed as follows. We generate flat distributions in the variables $`k_T^2/Q^2`$, $`1/Q^2`$, $`x_J,`$ using reference intervals which include the whole of the experimental phase-space (the $`\mathrm{\Phi }`$ variable is not used in the generation since all the cross-section measurements are $`\varphi `$ independent). In practice, we get the correction factors by counting the numbers of events which fulfill the experimental cuts given in Table I for each $`x`$-bin. The correction factor is obtained by the ratio to the number of events which pass the experimental cuts and the kinematic constraints, and the number of events which fullfil only the kinematic constraints,i.e. the so-called reference bin. The correction factors are given in reference .
Weperform a fit to the H1 and ZEUS data with only two free parameters. these are the effective strong coupling constant in LO BFKL formulae $`\overline{\alpha }`$ corresponding to the effective Lipatov intercept $`\alpha _P=1+4\mathrm{log}2\overline{\alpha }N_C/\pi `$, and the cross-section normalisation. The obtained values of the parameters and the $`\chi ^2`$ of the fit are given in Table III for a fit to the H1 and ZEUS data separately, and then to the H1 + ZEUS data together.
| fit | $`\overline{\alpha }`$ | $`\alpha _P`$ | Norm. | $`\chi ^2(/dof)`$ |
| --- | --- | --- | --- | --- |
| H1 | 0.17 $`\pm `$ 0.02 $`\pm `$ 0.01 | 1.44 $`\pm `$ 0.05 $`\pm `$ 0.025 | 29.4 $`\pm `$ 4.8 $`\pm `$ 5.2 | 5.7 (/9) |
| ZEUS | 0.20 $`\pm `$ 0.02 $`\pm `$ 0.01 | 1.52 $`\pm `$ 0.05 $`\pm `$ 0.025 | 26.4 $`\pm `$ 3.9 $`\pm `$ 4.7 | 2.0 (/2) |
| H1+ZEUS | 0.16 $`\pm `$ 0.01 $`\pm `$ 0.01 | 1.43 $`\pm `$ 0.025 $`\pm `$ 0.025 | 30.7 $`\pm `$ 2.9 $`\pm `$ 3.5 | 12.0 (/13) |
| D0 | 0.24 $`\pm `$ 0.02 $`\pm `$ 0.02 | 1.65 $`\pm `$ 0.05 $`\pm `$ 0.05 | | |
Table I- Fit results
The $`\chi ^2`$ of the fits have been calculated using statistical error only and are very satisfactory (about $`0.6perpoint`$ for H1 data, and $`1.perpoint`$ for ZEUS data). We give both statistical and systematic errors for the fit parameters. The values of the Lipatov intercept are close to one another and compatible within errors for the H1 and ZEUS sets of data, and indicate a preferable medium value ($`\alpha _P=1.41.5`$). We also notice that the ZEUS data have the tendency to favour a higher exponent, but the number of data points used in the fit is much smaller than for H1, and the H1 data are also at lower $`x`$. The normalisation is also compatible between ZEUS and H1. The fit results are shown in Figure 1 and compared with the H1 and ZEUS measurements.
## 2 Comparison with Tevatron results and prospects for LHC
The final result of our new determination of the effective pomeron intercept is $`\alpha _P=1.43\pm 0.025`$ (stat.) $`\pm 0.025`$ (syst.). Our method allows a direct comparison of the intercept values with those obtained in other experimental processes, i.e. $`\gamma ^{}\gamma ^{}`$ cross-sections at LEP , jet-jet cross-sections at Tevatron at large rapidity intervals , $`F_2`$ and $`F_2^D`$ proton structure function measurements . Let us first consider the known determinations of the effective intercepts in $`F_2`$ and $`F_2^D`$ measurements at HERA . It is known that the effective intercept determined in these measurements is rather low<sup>*</sup><sup>*</sup>*It is interesting to note that the “hard” Pomeron intercept obtained within the framework of two-Pomeron models fits with our determination. However our parametrization (1.1) corresponds to only one Pomeron.(1.2-1.3). This is the reason why these data can be both described by a DGLAP or a BFKL-LO fit Note that in the BFKL descriptions of these data , the effective intercept is taken to be constant, while the $`Q^2`$ dependence comes from the BFKL integration (see for instance formula (1.1)).
Now let us consider processes initiated by two hard probes which allow a more direct comparison between experiments and BFKL predictions. These processes suppress DGLAP evolution by selecting events with comparable hard scales for both hard probes. Recent data on $`\gamma ^{}\gamma ^{}`$ cross-section measurements at LEP lead to a BFKL description with a low effective intercept compatible with the one of $`F_2`$ and $`F_2^D`$ at HERA ($`\alpha _P`$=1.2-1.3 ) The statistics for these data is still very low. L3 and OPAL Collaborations have released the cuts used to enhance BFKL effects to get more statistics . These data can be both described by BFKL and DGLAP evolution equations.. The fact that similar values of the intercepts are found could be interpreted by sizeable higher order corrections to BFKL equation. On the other hand, it is interesting to note that our result based on forward jet measurement at HERA obtained in comparable $`Q^2`$ ($`Q^210`$ GeV<sup>2</sup>) and rapidity ($`Y`$ 3-4) domains is quite different. The value of the intercept is significantly higher.
Let us compare our results with the effective intercept we obtain from recent preliminary dijet data obtained by the D0 Collaboration at Tevatron . The measurement consists in the ratio $`R=\sigma _{1800}/\sigma _{630}`$ where $`\sigma `$ is the dijet cross-section at large rapidity interval $`Y\mathrm{\Delta }\eta `$ for two center-of-mass energies (630 and 1800 GeV), $`\mathrm{\Delta }\eta _{1800}=4.6`$, $`\mathrm{\Delta }\eta _{630}=2.4.`$ The experimental measurement is $`R=2.9\pm 0.3`$ (stat.) $`\pm 0.3`$ (syst.). Using the Mueller-Navelet formula , this measurement allows us to get a value of the effective intercept for this process
$`R={\displaystyle \frac{_{\frac{1}{2}i\mathrm{}}^{\frac{1}{2}+i\mathrm{}}\frac{d\gamma }{2i\pi \gamma (1\gamma )}e^{ϵ(\gamma ,0)\mathrm{\Delta }\eta _{1800}}}{_{\frac{1}{2}i\mathrm{}}^{\frac{1}{2}+i\mathrm{}}\frac{d\gamma }{2i\pi \gamma (1\gamma )}e^{ϵ(\gamma ,0)\mathrm{\Delta }\eta _{630}}}}.`$ (2.11)
We get $`\alpha _P`$=1.65 $`\pm `$ 0.05 (stat.) $`\pm `$ 0.05 (syst.), in agreement with the value obtained by D0 using a saddle-point approximation (see Table 1). This intercept is higher than the one obtained in the forward jet study.
Formula (2.11) is obtained after integration over the jet tranverse energies at 630 and 1800 GeV, $`E_{T_1}`$, $`E_{T_2}`$. We note that the non integrated formula
$`R(E_{T_1}/E_{T_2})={\displaystyle \frac{_{\frac{1}{2}i\mathrm{}}^{\frac{1}{2}+i\mathrm{}}\frac{d\gamma }{2i\pi }\left(\frac{E_{T_1}}{E_{T_2}}\right)^{2\gamma }e^{ϵ(\gamma ,0)\mathrm{\Delta }\eta _{1800}}}{_{\frac{1}{2}i\mathrm{}}^{\frac{1}{2}+i\mathrm{}}\frac{d\gamma }{2i\pi }\left(\frac{E_{T_1}}{E_{T_2}}\right)^{2\gamma }e^{ϵ(\gamma ,0)\mathrm{\Delta }\eta _{630}}}}`$ (2.12)
shows a sizeable dependence on $`E_{T_1}/E_{T_2}`$, which could be confronted with experiment. Let us show both the integrated and $`E_{T_1}/E_{T_2}`$ dependent cross-sections in Figure 2.
Prospects at LHC are quite appealing for this measurement due to the large rapidity intervals which can be reached. For instance, we estimate that it is possible to reach a rapidity interval $`\mathrm{\Delta }\eta _{14000}`$ of the order of 10 for a center-of-mass energy of 14 TeV. Taking as a reference value the maximal value at Tevatron ($`\mathrm{\Delta }\eta _{2000}5`$), which could be also reached at LHC by reducing the beam energies, we use Formulae (2.11, 2.12) with these values of rapidity ranges to get a prediction on the $`R`$ values (see figure 3). The BFKL prediction gives a high value of $`R`$ ($`R20`$), with a typical dependence on $`E_{T_1}/E_{T_2}`$ quite smaller than at Tevatron, which would be nice to test at LHC. If we consider two more close rapidity intervals (e.g. 10 and 8 corresponding to center-of-mass energies of 14 and 7 TeV), we get similar values of $`R`$ compared to Tevatron ($`R3.2`$). However, the $`E_{T_1}/E_{T_2}`$ dependence is much smaller than in Figure 2 for Tevatron. It favors the BFKL dynamics due to the high rapidity domains involved. It is thus important to perform this measurement and BFKL test at LHC.
The question arises to interpret the different values of the effective intercept for the different experimental processes. It could reasonably come from the differences in higher order QCD corrections for the BFKL kernel and/or in the impact factors depending on the initial probes ($`\gamma ^{}`$ vs. jets). In order to evaluate the approximate size of the higher order BFKL corrections, we will use their description in terms of rapidity veto effects . In formula (1.1), we make the following replacement
$`\mathrm{exp}(ϵ(\gamma ,0)Y)\mathrm{\Sigma }_{n=0}^{\mathrm{}}\theta (Y(n+1)b){\displaystyle \frac{\left[ϵ(\gamma ,0)(Y(n+1)b)\right]^n}{\mathrm{\Gamma }(n+1)}}.`$ (2.13)
The Heaviside function $`\theta `$ ensures that a BFKL ladder of $`n`$ gluons occupies $`(n+1)b`$ rapidity interval where $`b`$ parametrises the strength of NLO BFKL corrections. The value of the leading order intercept is fixed to $`\alpha _p=1.75(\alpha _S(Q^2=10)=0.28)`$, where $`Q^2=10`$ GeV<sup>2</sup> is inside the average range of $`Q^2`$ in the forward jet measurement. The fitted value of the $`b`$ parameter obtained using the forward jet data is found to be 1.28 $`\pm `$ 0.08 (stat.) $`\pm `$ 0.02 (syst.). Imposing the same value of $`\alpha _P`$ with Tevatron data gives $`b`$=0.21 $`\pm `$ 0.11 (stat.) $`\pm `$ 0.11 (syst.). Note that the theoretical value of $`b`$ for the NLO BFKL kernel is expected to be of the order 2.4, which is also compatible with the result obtained for the $`\gamma ^{}\gamma ^{}`$ cross-section. A contribution from the NLO impact factors is not yet known, and could perhaps explain the different values of $`b`$. The LHC would be also a very interesting testing ground for the study of rapidity veto and higher order BFKL effects.
## 3 Conclusion
To summarize our results, using a new method to disantangle the effects of the kinematic cuts from the genuine dynamical values we find that the effective pomeron intercept of the forward jet cross-sections at HERA is $`\alpha _P=1.43\pm 0.025`$ (stat.) $`\pm 0.025`$ (syst.). It is much higher than the soft pomeron intercept, and, among those determined in hard processes, it is intermediate between $`\gamma ^{}\gamma ^{}`$ interactions at LEP and dijet productions with large rapidity intervals at Tevatron, where we get $`\alpha _P`$=1.65 $`\pm `$ 0.05 (stat.) $`\pm `$ 0.05 (syst.).
Looking for an interpretation of our results in terms of higher order BFKL corrections expressed by rapidity gap vetoes $`b`$ between emitted gluons, we find a value of $`b=`$1.3 at HERA, and 0.21 at Tevatron. The HERA value is sizeable but less than the theoretically predicted value for the NLO BFKL kernel ($`b=`$2.4). The Tevatron value is compatible with zero. The observed dependence in the process deserves further more precise studies .
The LHC will open new possibilities to test different aspects of BFKL dynamics including higher order effects. Due to the large ranges in rapidity, large ratios of dijet cross-sections can be expected by using two different center-of-mass energies. We suggest to measure the dependence of the dijet cross-sections as a function of the jet transverse energies as a signal for BFKL pomeron at Tevatron run II, and at LHC. The Mueller Navelet jet study would also benefit from a lower energy run at LHC to allow a normalisation independence of the intercept determination and BFKL tests. LHC will thus allow a measurement of the pomeron intercept in a different kinematical domain more suited for BFKL dynamics, and a direct test of higher order effects.
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# Convective-Dynamical Instability in Radiation-Supported Accretion Disks
## 1 Introduction
Shakura & Sunyaev (1973) predicted that the inner portions of accretion disks that extend into relativistically-deep gravitational potentials should be radiation pressure-dominated when the accretion rate is greater than a modest fraction of the Eddington rate. In that regime, they found that disks could achieve hydrostatic balance in the vertical direction if the local dissipation rate were proportional to the local mass density. Given that assumption, upward radiation flux could support the disk matter against gravity if the density were essentially constant as a function of height (falling sharply to zero at the top surface) and the radiation pressure fell gradually from the disk midplane to the surface.
Soon after this equilibrium was discovered, it was found to suffer from several sorts of instabilities. Lightman & Eardley (1974) pointed out that if the viscous stress is proportional to the total pressure (in this case, dominated by radiation), perturbations with radial wavelengths long compared to the vertical thickness $`h`$, but short compared to a radius $`r`$, grow on the (comparatively long) viscous inflow timescale. Shakura & Sunyaev (1976) then observed that in these conditions perturbations in the same range of wavelengths would also grow on the (shorter) thermal timescale. Bisnovatyi-Kogan & Blinnikov (1977) noticed that if the radiation is locked to the gas even on short lengthscales (i.e., if, for the purpose of dynamics, the optical depth is treated as effectively infinite), such disks should be convectively unstable, for the specific entropy decreases outward; the linear growth rate for convective “bubbles” was worked out by Lominadze & Chagelishvili (1984). More recently, Gammie (1998) has demonstrated that a magnetic field in radiation-supported disks can catalyze a short-wavelength ($`kh1`$) overstable wave mode. In view of these instabilities, it has long been a puzzle just what sort of equilibrium would actually be found in Nature when the accretion rate is high enough that radiation pressure-domination might be expected (see, e.g. Shapiro, Lightman & Eardley 1976; Liang 1977; Coroniti 1981; Svensson & Zdziarski 1994; Szuszkiewicz & Miller 1997; Krolik 1998).
In this paper, we take a closer look at the nature of the short wavelength modes in radiation-supported disks without magnetic fields. Our goal (motivated by a companion work on radiation-hydrodynamics simulations of such disks: Agol & Krolik 2000b) is to examine more closely which modes can be expected to grow most quickly, what happens when finite optical depth permits some photon diffusion, and what role, if any, is played by associated perturbations in the local dissipation rate.
## 2 Problem Definition
We begin by writing down the equations of non-relativistic radiation hydrodynamics so that we may first describe the equilibrium in this language, and then discuss linear perturbations to this equilibrium. Because we are interested in accretion disks, it is convenient to write them in a rotating frame. The first is the usual equation of mass conservation:
$$\frac{\rho }{t}+\left(\rho \stackrel{}{v}\right)=0.$$
(1)
Our notation is the usual one, in which $`\rho `$ is the mass density and $`\stackrel{}{v}`$ is the fluid velocity. Next is the the fluid force equation:
$$\rho \frac{\stackrel{}{v}}{t}+\rho \stackrel{}{v}\stackrel{}{v}=p_g+\rho \stackrel{}{g}+(\kappa \rho /c)\stackrel{}{}+2\rho \stackrel{}{v}\times \stackrel{}{\mathrm{\Omega }}\rho v_r(\mathrm{\Omega }/\mathrm{ln}r)\widehat{\varphi }$$
(2)
where $`p_g`$ is the gas pressure, $`\stackrel{}{g}`$ is the local gravity, $`\kappa `$ is the opacity per unit mass, $`\stackrel{}{}`$ is the radiation flux, and $`\mathrm{\Omega }`$ is the rotation rate of the fluid. Although all the fluid quantities are defined in the rotating frame, the radiation quantities (e.g., $`\stackrel{}{}`$) are defined in the frame of the local fluid motion, i.e. including any departures from corotation (Mihalas & Mihalas 1984). Note that we have omitted magnetic forces.
For the two equations describing radiation energy density and momentum density, we follow Buchler (1979), but write the Lagrangian time derivatives explicitly, i.e., $`D/Dt=(/t+\stackrel{}{v})`$. The evolution of radiation energy density $`E`$ is described by:
$$\frac{E}{t}+\stackrel{}{v}E+\stackrel{}{}+𝐩_𝐫:\stackrel{}{v}+\frac{E}{c^2}\stackrel{}{v}+\frac{2}{c^2}\left(\frac{\stackrel{}{v}}{t}+\stackrel{}{v}\stackrel{}{v}\right)\stackrel{}{}=Q.$$
(3)
Here $`𝐩_𝐫`$ is the radiation pressure tensor and $`Q`$ is the net local emissivity. Finally, there is the equation describing the time-dependence of the radiation momentum density $`(1/c^2)`$:
$$\frac{1}{c}\frac{\stackrel{}{}}{t}+\frac{\stackrel{}{v}}{c}\stackrel{}{}+c𝐩_𝐫+\frac{1}{c}\left(\stackrel{}{}\stackrel{}{v}+\stackrel{}{}\stackrel{}{v}\right)+\frac{1}{c}\left(E𝐈+𝐩_𝐫\right)\left(\frac{\stackrel{}{v}}{t}+\stackrel{}{v}\stackrel{}{v}\right)=\stackrel{}{q}.$$
(4)
In this equation, $`𝐈`$ is the identity matrix and $`\stackrel{}{q}`$ is the net rate per unit volume at which photon momentum is created by radiation (usually negative in the fluid frame because photon momentum is lost due to opacity, while newly-created photons are usually isotropic in the fluid frame).
In the equilibrium, $`/t=\stackrel{}{v}=0`$. To isolate the effect of radiation support, we also take the extreme limit of $`p_gp_r`$. If we regard the rotation of the disk matter as cancelling the radial component of gravity, the only non-trivial remark to make about the equilibrium is that $`_z=cg/\kappa `$, where $`g=g_z(z,r)`$ is the local vertical component of gravitational acceleration \[in a thin disk, $`g_z(z,r)GMz/r^3`$ for central mass $`M`$\].
Now consider perturbed versions of equations (1) through (4). In order to write these perturbations in Fourier-transform form (i.e., for any quantity $`X`$, the perturbation is $`\delta \stackrel{~}{X}=\delta X\mathrm{exp}\{i(k_rr+k_zz)i\omega t\}`$), we will suppose that the wavevectors obey three conditions: that $`k=(k_z^2+k_r^2)^{1/2}\kappa \rho `$; that $`k_z1/h`$; and that $`k_r1/r`$. The first limit means that the diffusion approximation applies, i.e. $`𝐩_𝐫=p_r𝐈`$, so that $`E=3p_r`$. The second is the WKB approximation, as applied to variations in both the radial and vertical directions. Note that we further restrict our attention to axisymmetric perturbations. The condition $`k_zh1`$ also means that we can ignore any gradients in the gravity or equilibrium radiation flux. In addition, assuming $`k_r1/r`$ allows us to neglect the terms in vector divergences arising from cylindrical geometry. For example, after Fourier-transforming, $`\delta \stackrel{}{v}`$ becomes $`ik_r\delta v_r+\delta v_r/r+ik_z\delta v_zik_r\delta v_r+ik_z\delta v_z`$. We then find:
$`i\omega \delta \rho `$ $`+`$ $`\rho i\stackrel{}{k}\delta \stackrel{}{v}=0`$ (5)
$`i\omega \rho \delta v_r`$ $`=`$ $`(\kappa \rho /c)\delta _r+2\rho \mathrm{\Omega }\delta v_\varphi `$ (6)
$`i\omega \rho \delta v_z`$ $`=`$ $`(\kappa \rho /c)\delta _z`$ (7)
$`i\omega \rho \delta v_\varphi `$ $`=`$ $`(\kappa \rho /c)\delta _\varphi (1/2)\rho \mathrm{\Omega }\delta v_r`$ (8)
$`3i\omega \delta p_r3\rho g\delta v_z`$ $`+`$ $`ik_z\delta _z+ik_r\delta _r+4p_r\left(ik_z\delta v_z+ik_r\delta v_r\right)i{\displaystyle \frac{2\omega }{c^2}}\delta v_zF_z=\delta Q`$ (9)
$`i(\omega /c)\delta _r`$ $`+`$ $`ik_rc\delta p_r+ik_z(g/\kappa )\delta v_r4i(\omega /c)p_r\delta v_r=\kappa \rho \delta _r`$ (10)
$`i\left({\displaystyle \frac{\omega }{c}}\right)\delta _z+ik_zc\delta p_r`$ $`+`$ $`{\displaystyle \frac{g}{\kappa }}(2ik_z\delta v_z+ik_r\delta v_r)4i{\displaystyle \frac{\omega }{c}}p_r\delta v_z=cg\delta \rho \kappa \rho \delta _z`$ (11)
$`i(\omega /c)\delta _\varphi `$ $`+`$ $`i(k_z/c)_z\delta v_\varphi 4i(\omega /c)p_r\delta v_\varphi =\kappa \rho \delta _\varphi `$ (12)
The second and third equations in this set are the two components of the fluid force equation, with the vertical component reduced by the fact that the radiation flux exactly balances gravity in the equilibrium. The last three equations are the three components of the radiation flux equation. These equations may be further simplified by the assumption (verifiable post hoc) that $`\omega \kappa \rho c`$; that is, the wave frequency is very small compared to the mean time between photon scattering events.
With the further approximation that $`p_r\rho c^2`$, equations (5) – (12) may be manipulated to yield a dispersion relation. This relation is most conveniently displayed in terms of dimensionless quantities, so that $`\omega =\omega _{}\sqrt{g/h}`$ and $`\stackrel{}{k}=\stackrel{}{k}_{}/h`$. To satisfy the WKB approximation, $`k_z1`$ and $`k_rh/r`$. The dispersion relation is then seen to depend on six dimensionless combinations of parameters:
$`P`$ $`=`$ $`{\displaystyle \frac{p_r}{\rho gh}}`$ (13)
$`G`$ $`=`$ $`{\displaystyle \frac{\sqrt{gh}}{c}}`$ (14)
$`\mathrm{\Omega }_{}^2`$ $`=`$ $`h\mathrm{\Omega }^2/g`$ (15)
$`\tau `$ $`=`$ $`\kappa \rho h`$ (16)
$`𝒬_p`$ $`=`$ $`(Q/p_r)(h/g)^{1/2}`$ (17)
$`𝒬_\rho `$ $`=`$ $`(Q/\rho )(hg^3)^{1/2}`$ (18)
Note that here $`g`$ is $`g=g_z(z)`$, the local value of the vertical gravity. In any thin disk, the quantity $`G1`$, for it is of order the free-fall speed from the top of the disk to the midplane in units of $`c`$. The relative importance of rotational effects is given by $`\mathrm{\Omega }_{}^2`$, which is simply $`h/z`$ for Keplerian rotation. The parameter $`P`$ is also a function of height above the midplane. In an optically thick disk, $`P=\frac{1}{2}[1+1/\tau (z/h)^2]/(z/h)`$. Thus, the two parameters $`P`$ and $`\mathrm{\Omega }_{}^2`$ could be replaced by the single parameter $`z/h`$. The last two dimensionless parameters describe the sensitivity of the dissipation to the radiation pressure and the density, respectively. If one could extrapolate the “$`\alpha `$” prescription to local fluctuations, one might expect that $`𝒬_p\alpha `$.
## 3 The Dispersion Relation
### 3.1 General considerations
In terms of these dimensionless quantities, the dispersion relation is:
$`3\omega _{}^5+\left[ik_{}^2/(G\tau )3k_zG/\tau i𝒬_p\right]\omega _{}^4`$
$``$ $`\left(4k_{}^2P+3\mathrm{\Omega }_{}^2+ik_z^3/\tau ^2ik_zG𝒬_p/\tau \right)\omega _{}^3`$
$`+`$ $`\left[4k_z^3PG/\tau ik_{}^2\mathrm{\Omega }_{}^2/(G\tau )+i𝒬_p(ik_z+\mathrm{\Omega }_{}^2)ik_{}^2𝒬_\rho \right]\omega _{}^2`$
$`+`$ $`\left[4k_z^2P\mathrm{\Omega }_{}^23k_r^2+k_z^2G(ik_z𝒬_\rho 𝒬_p)/\tau \right]\omega _{}`$
$`+`$ $`k_zk_r^2\mathrm{\Omega }_{}^2/(G\tau )+k_z\mathrm{\Omega }_{}^2(ik_z𝒬_\rho 𝒬_p)=0.`$
In the limit of $`\tau \mathrm{}`$ and $`𝒬_p=𝒬_\rho =0`$, this dispersion relation simplifies to
$$\omega _{}^5\left[(4/3)k_{}^2P+\mathrm{\Omega }_{}^2\right]\omega _{}^3+\left[(4/3)k_z^2P\mathrm{\Omega }_{}^2k_r^2\right]\omega _{}=0.$$
(19)
One root is clearly $`\omega _{}=0`$. The other four are given by:
$$\omega _{}^2=\{\begin{array}{cc}(4/3)k_{}^2P\hfill & \\ \left[2k_z^2\mathrm{\Omega }_{}^2(3/2)k_r^2/P\right]/k_{}^2\hfill & \end{array}$$
(20)
The first pair of roots are the familiar radiation-supported sound waves. The second pair describe buoyancy behavior (cf. Balbus 1999). If $`k_z^2\mathrm{\Omega }_{}^2(3/4)k_r^2/P`$, there are two neutrally stable gravity (epicyclic) waves; on the other hand, if $`k_z^2\mathrm{\Omega }_{}^2<(3/4)k_r^2/P`$, one mode is damped, but the other grows exponentially with essentially no oscillation. It is this last mode that corresponds to convection. Although rotation tends to have a stabilizing effect, when $`\mathrm{\Omega }_{}>0`$ one can always find a mode with $`k_r/k_z`$ large enough to satisfy this criterion. As a consequence, when radiation pressure dominates gas pressure in a very optically thick disk, the standard equilibrium is always unstable to convection. However, relatively strong rotational effects (i.e., $`\mathrm{\Omega }_{}^2P(h/z)^2`$ relatively large, or location near the midplane) do diminish the range of wavevector directions that is unstable.
### 3.2 Dependence on parameters
As already remarked, both $`P`$ and $`\mathrm{\Omega }_{}`$ are determined by $`z/h`$, and both should be, except very near the midplane, $`1`$. Because the growth rate of the fastest growing mode in the optically thick limit is $`P^{1/2}`$, and $`P`$ decreases upward, convective instability should develop most rapidly near the top of the disk. This expectation is borne out when the full equations are solved: as predicted by equation 20, when $`k_rk_z`$, the growth rate is nearly independent of $`|\stackrel{}{k}|`$ and is $`7\times `$ greater at $`z/h=0.9`$ than at $`z/h=0.1`$.
The optical depth $`\tau `$ in a radiation pressure-dominated Shakura-Sunyaev disk is
$$\tau =\frac{4/3}{\alpha }\dot{m}^1x^{3/2}\left(\frac{\kappa }{\kappa _T}\right)^1\frac{R_zR_T}{R_R^2},$$
(21)
where $`\alpha `$ is the usual dimensionless stress parameter, $`\dot{m}`$ is the accretion rate in Eddington units for unit efficiency, $`x`$ is the radius in units of $`GM/c^2`$, $`\kappa /\kappa _T`$ is the opacity relative to pure Thomson opacity, and $`R_{R,T,z}`$ are the relativistic corrections to the dissipation rate, torque, and vertical gravity (Page & Thorne 1974; Abramowicz, Lanza & Percival 1997; see also Agol & Krolik 2000a). Because we expect $`\alpha ,\dot{m}<1`$, but $`x`$ must be $`>1`$, these disks should be optically thick. However, in the inner part of the disk, $`\tau `$ might be as little as $`10^2`$.
Fig. 1 shows how the growth rate and phase velocity of the unstable convective modes depend on wavenumber and $`k_z/k_r`$ when there is no perturbation to the dissipation rate and the optical depth is very large ($`\tau =10^6`$). To illustrate the dependence on optical depth, the same curves are shown in Fig. 2 for a case with the smallest optical depth one might expect, $`\tau =10^2`$. In both cases, the range of meaningful wavenumbers (and hence the range of wavenumbers displayed in the figures) is limited by two criteria. On the one hand, the WKB approximation demands that $`k_zh1`$; on the other, the diffusion approximation is only consistent with $`kh\tau `$. We present results extending out to $`kh\tau `$; at the highest wavenumbers shown, the diffusion approximation has only marginal validity.
For fixed optical depth, the growth rate depends on $`k_z/k_r`$; when $`k_z/k_r\stackrel{>}{}1.2`$, the unstable convective mode disappears. For smaller $`k_z/k_r`$ ($`<0.1`$), the growth rate is independent of $`k_z/k_r`$. Decreasing the optical depth reduces the extent of the “flat” portion of the growth rate curve for a given $`k_z/k_r`$. However, the largest value of the dimensionless growth rate, $`(\omega _i)_{max}`$, in the wavenumber range of physical interest is almost independent of $`\tau `$ over this range.
Figs. 1 and 2 also show the phase velocity of the unstable convective mode, normalized to the free fall speed in the disk. For large optical depths, the real frequency of convective modes is so close to zero that they can be considered as non-oscillatory; decreasing the optical depth towards the smallest value ($`\tau =10^2`$) expected for radiation pressure dominated disks leads to an increase in the real frequency of the mode, endowing it with a non-negligible phase velocity and, as a consequence, it becomes oscillatory.
In the absence of perturbations to the dissipation rate, finite (although large) optical depth primarily affects the radiation sound waves (cf. Agol & Krolik 1998). Not surprisingly, as the opacity falls, short wavelength sound waves become damped. Less intuitively, convective modes remain unstable even when the diffusion rate ($`k_{}^2c/(h\tau )`$) is larger than the characteristic frequency $`\sqrt{g/h}`$, i.e. when $`k_{}>\sqrt{G\tau }`$. In this regime, epicyclic motions cause convective modes to acquire an oscillatory character. Convection remains unstable even in the face of strong radiative diffusion because the equilibrium gas density is constant. A parcel that begins to fall because its specific entropy is too small is squeezed as it encounters higher pressure. Its radiation pressure rises, but because diffusion is so effective, only up to the ambient level. Meanwhile, however, the increased gas density means that the parcel’s density exceeds that of its neighbors, and it continues to fall.
Such impunity to the effects of radiation diffusion does not persist all the way to zero opacity, however. When $`k_{}>\tau `$, the diffusion approximation is no longer a valid description of radiation transfer. For wavelengths that short (or opacity that small), the radiation streams freely, entirely independent of gas motions. In that limit, if there is nothing to alter the distribution of dissipation, there is no perturbation to the flux as a function of position, and therefore, no perturbation to the force. The result, of course, is that convection disappears.
Again applying the Shakura-Sunyaev disk solution, we expect $`G`$ to be small, for it is given by
$$G=\frac{3}{2}\dot{m}\left(\frac{\kappa }{\kappa _T}\right)\left(\frac{z}{h}\right)^{1/2}x^{3/2}\frac{R_R}{R_z^{1/2}}.$$
(22)
$`G`$ becomes irrelevant over most of the wavenumber interval of interest when $`\tau \mathrm{}`$, as every factor in which it appears in the dispersion relation is divided by $`\tau `$; the value of $`G`$ can still have some influence even at large $`\tau `$ only for very large wavenumbers, $`k_{}\stackrel{>}{}0.01\tau `$, since some of those same terms of the dispersion relation actually are $`k_{}^2/\tau `$ or even $`k_{}^3/\tau `$. However, even when $`\tau `$ is as large as $`10^5`$, some dependence on $`G`$ begins to appear across the range of interesting wavenumbers. When the optical depth is no larger than this, diminishing $`G`$ tends to enhance instability, particularly at short wavelengths. When $`\tau 10^5`$, the growth rate of the very short wavelength convective modes increases by factors of several when $`G`$ falls from $`10^2`$ to $`10^4`$. When $`\tau \stackrel{<}{}10^4`$, and $`G10^4`$, a second mode, in addition to the convective mode, becomes unstable at short wavelengths. This mode is weakly oscillatory, and its very low phase velocity ($`2.5\times 10^2`$$`1.6\times 10^3`$ the free fall velocity) behaves similarly to that of damped radiation sound waves at the same wavelengths. Analysis of the corresponding amplitudes of the various physical quantitities’ perturbations shows that this mode has magnitude and especially phase relationships between perturbations in pressure, density and velocity components that are intermediate between those of convective modes and those characteristic of radiation sound waves.
It is important to note, however, that if the canonical disk solution applies, the product
$$G\tau =2\left(\frac{z}{h}\right)^{1/2}\alpha ^1\frac{R_TR_z^{1/2}}{R_R}.$$
(23)
In other words, the product $`G\tau `$ should be very nearly independent of radius (if $`\alpha `$ is). Even where the relativistic corrections are substantial, $`G\tau `$ hardly changes because $`R_R`$ and $`R_T`$ are very nearly proportional, and $`R_z`$ has a total range of at most a factor of two (Krolik 1999). Thus, the criterion for which scaled wavelengths should be substantially affected by photon diffusion (i.e., $`k_{}>\sqrt{G\tau }`$) is almost independent of radius, but does depend somewhat on $`z/h`$.
Given our estimate of $`G`$ (equation 22), we can also verify our expectation that $`\omega /c\kappa \rho 1`$. Rewritten in terms of our dimensionless parameters, this ratio is $`\omega _{}G/\tau `$. As we have just estimated, $`G1`$ and $`\tau 1`$ in radiation pressure-supported $`\alpha `$-disks. Consequently, unless $`\omega _{}`$ is extremely large, $`\omega /c\kappa \rho 1`$ is a very safe assumption.
As already pointed out, $`𝒬_p`$ is much like a local version of the vertically-averaged quantity $`\alpha `$, the ratio of the $`r`$$`\varphi `$ stress to the total pressure. On this basis we expect $`𝒬_p`$ to be somewhat less than unity. The (dimensional) dissipation perturbation coefficient $`Q_\rho (Q/\rho )`$, on the other hand, is most closely related to the ratio of the flux to the surface density. In a Shakura-Sunyaev disk, this quantity (reduced to a dimensionless number according to equation 18), is
$$\left(\frac{\stackrel{~}{F}}{\stackrel{~}{\mathrm{\Sigma }}}\right)\frac{}{\mathrm{\Sigma }}\frac{1}{(hg^3)^{1/2}}=\frac{1}{2}\left(\frac{3}{2}\right)^{3/2}\alpha \dot{m}^{3/2}\left(\frac{\kappa }{\kappa _T}\right)^{3/2}x^{3/2}\frac{R_R^{5/2}}{R_z^2R_T}.$$
(24)
Although we don’t know for certain that $`𝒬_\rho \stackrel{~}{F}/\stackrel{~}{\mathrm{\Sigma }}`$, it is a likely estimator. On this basis, we expect $`𝒬_\rho 1`$, and should be particularly small at larger radius or in disks with relatively low accretion rate.
Like $`\tau `$, varying $`𝒬_p`$ or $`𝒬_\rho `$ affects the radiation sound waves much more than the convective modes. The growth rate of the near-zero real frequency convective mode is almost independent of $`𝒬_p`$ and $`𝒬_\rho `$ over a large range in these parameters, but the longer-wavelength radiation sound waves become unstable when these parameters are positive. For example, when $`z/h=0.5`$, and $`\tau 1`$, the growth rates for positive $`𝒬_p`$ and null $`𝒬_\rho `$ are $`(0.10.16)𝒬_p`$, and, similarly, those for positive $`𝒬_\rho `$ and null $`𝒬_p`$ turn out to be $`(0.10.16)𝒬_\rho `$.
Unlike the convective mode, for radiation sound waves, growth persists to $`k_r/k_z<1`$. Although the maximum growth rate does not depend on the value of this ratio, the range of wavelengths for which growth occurs shrinks and moves towards larger wavelengths as the ratio $`k_r/k_z`$ decreases.
Smaller optical depth can also counteract the instability created by positive $`𝒬_p`$ or $`𝒬_\rho `$ (see also Agol & Krolik 1998 for a discussion of radiation diffusion damping of MHD waves). In the very large $`\tau `$ limit, the maximum growth rate for radiation sound waves is almost independent of the value of the optical depth, but for $`\tau \stackrel{<}{}10^4`$ it decreases with decreasing $`\tau `$. The radiation sound waves are all damped for $`\tau <10^3`$ or so; at larger optical depth, the shorter wavelength modes remain damped (because, of course, diffusion is most effective acting on them), but growth appears for longer wavelengths, with the range of growing wavelengths stretching as $`\tau `$ increases. Fig. 3 shows the behaviour of the imaginary part of the frequency and the normalised phase velocity for radiation sound waves with $`𝒬_p=1`$ and $`𝒬_\rho =0`$ as functions of wavenumber for a variety of optical depths.
Summarizing this discussion, we conclude that there are two sorts of instabilities of importance in radiation pressure-supported disks: convective instabilities, which grow on roughly the dynamical timescale, should be very nearly ubiquitous unless the disk becomes almost optically thin; and radiation sound waves, which can be driven unstable (at a slower rate) if $`𝒬_p`$ or $`𝒬_\rho `$ is positive.
### 3.3 Comparison with the Classical Convection Problem
In standard treatments of convection in rotating fluids (e.g., Chandrasekhar 1961), it is shown that convection begins when the “Rayleigh number”
$$=\frac{g\alpha _V\beta h^4}{\kappa _{th}\nu }$$
(25)
exceeds a critical value (generally $`1`$) which depends on the “Taylor number”
$$𝒯=4\frac{\mathrm{\Omega }^2h^4}{\nu ^2}.$$
(26)
Writing $`T`$ for the temperature, we define $`\alpha _V|\mathrm{ln}\rho /T|`$, $`\beta |dT/dz|`$, $`\kappa _{th}`$ is the thermal conductivity, and $`\nu `$ is the kinematic viscosity.
This analysis applies partially, but not entirely, to the circumstances of radiation-dominated accretion disks. The fact that, rather than entering the fluid solely at the bottom, heat is actually generated throughout the fluid probably alters the quantitative criteria for convective instability but should not change anything at a qualitative level. On the other hand, certain contrasts are of greater importance.
One that requires mention is that the pressure in these circumstances, unlike in classical fluids, depends only on the radiation, and is therefore independent of density. Consequently, thermal diffusion can alter the pressure with no change in density. In the classical analog, thermal diffusion can quench convective instability by adding heat to falling low-entropy fluid elements, thereby forcing them to expand until their density no longer exceeds the surrounding density. This occurs when $`\kappa _{th}`$ is large enough that $``$ falls below the critical value. By contrast, in the radiation-dominated case, the density remains unchanged in the face of thermal diffusion. This is why convection persists even when $`k_{}>\sqrt{G\tau }`$; i.e., the radiation analog of the Rayleigh number does not accurately predict the effect of thermal diffusion on the instability.
Perhaps the most important contrast with classical fluids is that there is no easy definition of the kinematic viscosity. Bulk orbital shear leads to angular momentum transfer via magnetic forces (Balbus & Hawley 1998); these are, of course, entirely absent in our treatment. However, stresses arising from other shears may be quite different, whether they are caused by magnetic forces or other mechanisms (for example, photon diffusion: Agol & Krolik 1998). For this reason, it is unclear whether the conventional equation of the $`r\varphi `$ shear stress with a nominal viscosity (as in the Shakura-Sunyaev $`\alpha `$ formalism) is appropriate here. In any event, none of these possible viscosity mechanisms is present in our equations; not even radiation viscosity is considered, for no radiation shear stress exists in the pure diffusion approximation. Thus, in a formal sense, neither $``$ nor $`𝒯`$ exists in our treatment of this problem.
Nonetheless, having expressed these caveats, it is still of some interest to estimate what $`𝒯`$ and $``$ may be in radiation-dominated accretion disks in order to make some contact with the classical theory. To do so we make the suspect identification discussed in the preceding paragraph: we suppose that the shear viscosity in all directions is given by the Shakura-Sunayev $`\alpha `$. If so, $`𝒯10\alpha ^2`$, independent of radius (or anything else). Within the radiation-dominated regime, the Rayleigh number is similarly constant: $`O(10)(\kappa /\kappa _T)^2`$. Interestingly, the critical $``$ for $`10<𝒯<10^5`$ is $`10^3`$$`10^4`$ (Chandrasekhar 1961). Thus, if the most appropriate analog to ordinary fluid viscosity is correctly estimated by this means, the strong convective instability we find might be quenched; however, as we have stressed, it is by no means clear that this analogy applies.
## 4 Implications
We have shown that if an accretion disk is placed in the Shakura-Sunyaev radiation-dominated equilibrium, the growth rate of the convective instability is hardly affected by large changes in parameters such as the opacity or the way in which the local dissipation might be perturbed. This behavior is, perhaps, not surprising, when one considers that the convective instability is essentially dynamical, whereas the opacity and dissipation rates pertain to the (rather more slowly accomplished) thermal balance.
Given that the Shakura-Sunyaev equilibrium is unstable to convection for all reasonable parameters, we can hardly expect that it describes the actual state of radiation pressure-dominated disks. One’s first guess might be that what happens instead is that entropy is redistributed so as to make the disks at most marginally unstable to convection (Bisnovatyi-Kogan & Blinnikov 1977). If the entropy distribution is at most weakly unstable to convection, then departures from hydrostatic equilbrium should be small, and the equilibrium could be determined by applying these two conditions (i.e. isentropy and hydrostatic balance). However, marginal stability to convection also suggests that the amount of heat carried by convection is small, so that radiative diffusion dominates energy exchange. Unfortunately, this condition is inconsistent with the first two. Any two of these three conditions suffice to determine the equilibrium; the solution is over-determined if all three are applied. The result must be, then, that in real disks none of them is satisfied exactly, but all three are approximately correct. The quantitative character of this balancing act can be determined only by detailed calculations that carry convection into the nonlinear regime (e.g. Agol & Krolik 2000b).
Deviation from constant specific entropy can certainly be expected near the disk surface. There, radiation diffusion should be able to consistently dominate convection as a heat-loss mechanism because the diffusion time is much shorter than in the body of the disk. Consequently, the specific entropy distribution near the surface must be clearly stable, i.e. rise with increasing altitude.
The sense of heat redistribution by convection is to raise the specific entropy near the top and diminish it below. For the same reason that radiative diffusion must dominate near the top, this redistribution can have the side effect of changing the overall rate at which heat is able to leave the disk. The mean specific entropy in equilibrium is therefore not necessarily the same as the value predicted by the original Shakura-Sunyaev equilibrium.
These remarks do not bear on whether the thermal and viscous instabilities continue to plague radiation-dominated disks. It is possible that convection might be able to remove extra heat generated in the course of the thermal instability, and might consequently be able to quench it. Whether that is so can only be determined by detailed calculation of heat transport by convection in the nonlinear regime (Agol & Krolik 2000b). Whether the viscous instability acts under these conditions depends on the character of dynamical coupling (see, e.g., Agol & Krolik 1998) between the radiation and the MHD turbulence that most likely accounts for angular momentum transport in disks (Balbus & Hawley 1998).
We have also shown that radiation sound waves may be driven unstable if the emission of radiation is proportional to either the gas density or the local radiation pressure (the latter possibility is suggested by a straight-forward extrapolation of the “$`\alpha `$-model). Although phenomenological guesses of this variety are plausible, we do not yet know enough about the physical mechanisms of dissipation in accretion disks to say whether this is what should actually happen. However, if it does occur, growth of radiation sound waves (and, presumably, magnetosonic waves) could provide an additional source of turbulence in disks, supplementing magneto-rotational instability and convection. Indeed, if dissipation leads to renewed stirring, there might be interesting twists to the character of the MHD turbulent cascade.
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# Massive star clusters in dwarf starburst galaxies
## 1. Introduction
During the last decade, the study of young massive stellar clusters, sometimes referred to as super star clusters (SSCs) has seen rapid progress. From a few known examples of “blue populous clusters” in the Large Magellanic Cloud (LMC) and SSCs in the star forming dwarf galaxies NGC 1705 (Melnick, Moles, & Terlevich 1985) and NGC 1569 (Arp & Sandage 1985), the numbers have grown, largely thanks to HST. SSCs have now been found in a variety of different environments, as have globular clusters (GCs) which are old massive clusters. In this paper I will review the massive cluster content of, in particular, low mass starburst galaxies like blue compact galaxies (BCGs). A review of massive clusters in starbursts naturally becomes biased towards young SSC like objects. However, such galaxies may indeed contain also rich populations of older clusters.
The study of SSCs in starburst galaxies gained momentum with the advent of the HST. Ultraviolet imaging with the aberrated HST/FOC (e.g. Meurer et al. 1995, Conti and Vacca 1994) of starbursts revealed that a significant fraction of the star formation activity took place in SSCs. Meurer et al. (1995) studied nine galaxies, finding SSCs in most of them. Several other starbursts have been imaged with the aberrated HST/FOC (e.g. Conti et al., unpublished) and SSCs are frequently encountered. Optical imaging with the aberrated HST has also discovered SSCs in many low mass starburst galaxies (e.g. Hunter et al. 1994). The greater sensitivity of WFPC2 as compared to FOC has multiplied the number of detected clusters in ESO 338-IG04 (Östlin, Bergvall, & Rönnback 1998) and He 2-10 (Johnson et al. 2000).
In giant galaxies there seems to be several ways to form SSCs/GCs, e.g. mergers (see Miller 2000), and bars (e.g. Kristen et al. 1997) and circum-nuclear rings in spiral galaxies (e.g. Barth et al. 1995). There are even indications of SSC formation in the discs of normal spirals (Larsen & Richtler 1999). In low mass galaxies some of these mechanisms, e.g. formation of SSCs in bar and resonant induced density enhancements, are not available. Mergers are certainly producing SSCs in some low mass galaxies (e.g. ESO 338-IG04, ESO 350-IG38, ESO 185-IG13), but there might be other mechanisms too. The origin of active star formation in galaxies like NGC 1569 and NGC 1705 are not yet well understood. There are also dwarf stabursts which do not contain luminous SSCs (e.g. IC10, see Grebel 2000).
Cluster destruction mechanisms (e.g. due to tidal shocks) are weaker in low mass galaxies giving SSCs a greater chance of survival. In ESO 338-IG04 the dominant GC population is $`3`$ Gyr old, and this population alone is enough to classify the galaxy as GC rich in terms of specific frequency. Low mass galaxies are in general metal-poor (typically \[O/H\]$`1`$, Kunth & Östlin 2000), which make them suitable for comparison with high redshift conditions and early GC formation. Dwarf galaxies have certainly been important ingredients in the hierarchial buildup of massive galaxies. Another virtue is that internal extinction in general is small, which for instance makes age dating more secure.
Thus, dwarf starbursts are good places to investigate the formation and evolution of massive star clusters. Even if one cannot be sure whether a SSC will evolve into a bona fide GC or dissolve, a proto-GC must look very much like a SSC (Kennicutt & Chu 1988) and the collective formation of a few SSCs will, by necessity, be associated with a starburst. Thus populations of GCs and survived SSCs trace former starbursts, and may be used to study the evolution of galaxies. For example, if the excess of faint blue galaxies seen in deep optical counts is due to starbursts originating in merging galaxies at intermediate redshift, one would expect these to form significant numbers of SSCs/GCs which should be visible as intermediate age GC populations in local galaxies.
## 2. SSC richness in galaxies
The richness of GCs in galaxies is often parametrized by the specific frequency, $`S_N`$, which is the number of GCs divided by the host galaxy luminosity (Harris & van den Bergh 1981). If counting only luminous clusters ($`M_V11`$) and relating to $`M_{V,host}`$, the total absolute $`V`$ magnitude of the host galaxy, one may define a specific frequency of luminous SSCs: $`S_{11}=N_{11}\times 10^{0.4(M_{V,host}+15)}`$, where $`N_{11}`$ is the number of clusters with $`M_V11`$. $`S_{11}`$ will be independent of the assumed distance to a galaxy, and will be biased towards young massive clusters. The luminosity limit excludes contamination by supergiant stars and makes it possible to compare galaxies at different distances studied at different depth.
In Table 1, $`S_{11}`$ is given for a selection of galaxies of different types, where $`N_{11}`$ can be obtained with reasonable accuracy (in nearby galaxies like M82 the large angular extent makes it difficult to estimate $`N_{11}`$). The table is by no means complete (see e.g. Miller 2000, for more references). The $`N_{11}`$ values do not take internal extinction in the galaxies into account, but since the same is true for $`M_{V,host}`$, this should be a second order effect. If the luminous SSCs follow a power law luminosity function, $`\varphi (L)L^\alpha `$, with $`\alpha =2`$, $`S_{11}`$ will be independent of the internal extinction as long as $`A_V`$ is equal for the SSCs and the integrated galaxy light. This is not allways true since young SSCs may suffer from high local extinction. If $`\alpha <2`$, internal extinction will effectively lower the observed $`S_{11}`$ values. All values in Table 1 assume $`H_0=75`$ km/s/Mpc. The $`M_{V,host}`$ values were assembled from different sources including both accurate CCD photometry and photographic magnitudes. Many apparent total V-magnitudes have been taken, or estimated, from NED. Thus there might be an inhomogeneity on the 0.5 magnitude level (corresponding to 50% uncertainty on $`S_{11}`$), which one should be aware of when comparing galaxies.
Values $`S_{11}0.5`$ are found for giant mergers, while the luminous BCGs discussed below have $`S_{11}0.6`$. There are indications that these galaxies are actually dwarf mergers. The nearby star forming galaxies NGC 1569, NGC 1705 and NGC 1140, have intermediate $`S_{11}`$ values. For LMC, no cluster is bright enough to qualify without correcting for internal extinction. Most nearby dwarf irregular (dI) galaxies with modest star formation do not contain luminous SSCs (see Grebel 2000). In the mergers, the number of luminous SSCs decrease with the estimated age of the merger remnant as was found by Schweizer et al. (1996).
## 3. A few case studies
### 3.1. ESO 338-IG04 (= Tol 1924-416)
Ground based imaging had revealed an overdensity of faint blobs around this well known blue compact galaxy, which was the motivation behind HST/WFPC2 followup observations (Östlin et al. 1998). These observations revealed a starburst region composed of numerous blue star clusters and a swarm of surrounding GCs. In all, the number of star clusters (after correction for contamination by foreground stars, supergiant stars and background galaxies) is above 100 (completeness limit $`M_V9`$ to $`10`$). Spectral synthesis modelling of $`U,B,V,I`$ colors indicated the presence of distinct peaks in the cluster formation history. In addition to the ongoing event there are old ($`10`$ Gyr) GCs, and in particular a very prominent population of intermediate age (2 to 5 Gyr old) GCs. This intermediate age population contains the most massive cluster candidates and among them an object (no. 34 in the outer sample of Östlin et al. 1998) with an estimated mass in excess of $`10^7M_{}`$ (for a variety of different IMFs). ESO 338-IG04 is an intrinsically luminous ($`M_V=19.3`$) metal poor (\[O/H\]$`=1`$) BCG. The dynamics and morphology suggest that the galaxy is the product of a dwarf galaxy merger (Östlin et al. 1999, 2000). The current specific frequency, $`S_N=2`$, is predicted to rise to $`S_N10`$ in one Gyr as the starburst ceases and fades. The most luminous SSCs were found already by Meurer et al. (1995).
### 3.2. He 2-10 (= ESO 495-G21)
HST/FOC observations of He2-10 revealed “knots”, some of which appeared to be resolved with diameters $`<10`$ pc (Conti & Vacca 1994). It was suggested that these might be young GCs, but the use of a single UV passband made it hard to say more (e.g. addressing masses) than that the objects were young. More recent WFPC2 observations (Johnson et al. 2000) confirm the FOC results, and multiply the detected number of clusters by reaching fainter and redder objects. The faint, red objects may be intermediate age GCs or reddened SSCs. The galaxy is one of the most nearby BCGs and a candidate dwarf merger.
### 3.3. Markarian 996
Markarian 996 is a blue compact dwarf ($`M_V=17.2`$) with regular elliptical outer isophotes and intense star formation in a very compact central Hii region (Thuan, Izotov, & Lipovetsky 1996). The central Hii region is bright enough to saturate the WFPC2 images, but it could be a luminous ($`M_V<12`$) young SSC. Intererestingly there are plenty of old GCs, asymetrically distributed around Mrk 996, with a luminosity function similar to that of Galactic GCs. Mrk 996 is, together with ESO 338-IG04, one of the rare known examples of a BCG posessing an old GC population. Mrk 996 has a GC specific frequency $`S_N>5`$, similar to low-luminosity/dwarf ellipticals (Miller et al. 1998).
## 4. The Malkan et al. (1998) snapshot survey of AGN
Malkan, Gorjian, & Tam (1998) conducted an HST/WFPC2 snapshot survey of AGN, but included also a comparison sample consisting of 50 galaxies classified as having Hii-activity. These “Hii” galaxies can be divided in two broad classes: irregular galaxies and spiral galaxies with nuclear star formation.
Of galaxies having irregular or distorted morphology 3/4 appears to contain at least a few SSCs, while only 1/3 of the spirals do. Galaxies classified as irregular here are generally rather luminous and not typical dI galaxies like those in the local group. It is anyway obvious that the irregular galaxies, in view also of a lower average luminosity, appear to be more efficient SSC formers than spirals. Part of this effect could be explained by higher average extinction in spirals than irregulars. The survey used only one pass band (F606W) and rather short exposures which were not cosmic ray split, making it of limited use for quantitative studies. However, luminous SSCs can be easily identified and approximate photometry may be otained from archive images.
The survey includes ESO 185-IG13 and ESO 350-IG38 which have been studied dynamically by Östlin et al. (1999,2000). Both dynamics and morphology give strong support for a merger induced origin of the strong starbursts seen in these galaxies, which have similar properties to ESO 338-IG04. These galaxies are also among the most cluster rich in the whole survey. Extracting photometry for the SSCs in the HST archive images results in $`S_{11}=0.9`$ and $`S_{11}=1.5`$ for ESO 350-IG38 and ESO 185-IG13 respectively, see Table 1. ESO 350-IG38 contains several SSCs with $`M_V15`$.
Another example, perhaps the best one in the whole survey, of a very SSC rich starburst is Mrk 930, for which $`S_{11}=1.9`$, see Table 1. Also Mrk 930 has a very irregular morphology. The Malkan et al. (1998) survey also has a few objects in common with the study of Meurer et al. (1995).
## 5. Conclusions and Perspectives
Among dwarf and low mass galaxies we encounter both galaxies that appear to be very efficient formers of massive clusters, and galaxies that appear totally devoid of such objects.
Some luminous BCGs have specific frequncies of luminous ($`M_V11`$) SSCs $`S_{11}1`$, which is an order of magnitude larger than most of the prototypical SSC factories: the “Antennae”, NGC 7252 and NGC 1275. Thus there is a tendency for $`S_{11}`$ to increase when going to low mass starburst galaxies. A similar trend has been found for the specific frequency of GCs among dwarf ellipticals (Miller et al. 1998). The purpose of comparing $`S_{11}`$ values was to show that low mass, metal-poor, starburst galaxies are excellent hunting grounds for luminous SSCs, and in addition problems with extinction are much smaller than in giant mergers. The BCGs with the highest $`S_{11}`$ are believed to be the product of dwarf galaxy mergers (Östlin et al. 1999, 2000). A possible explanation to the higher $`S_{11}`$ values is that the starburst timescales are shorter in systems with lower mass, whereas in a giant merger one expects a more extended starburst. There are also BCGs and low mass starbursts which do not contain luminous SSCs. Although the expected number of SSCs in galaxies of very low luminosity will always be small and subject to statistical fluctuations, there migh be a connection to the triggering mechanism of the starbursts. Merging dwarfs might be more efficient SSC formers than non-merging ones.
A couple of BCGs (ESO 338-IG04 and Mrk 996) in addition contains rich populations of older GCs. There is no a priori reason to believe that old GC systems are intrinsically rare among BCGs. Rather few BCGs have been studied at sufficient depth and spatial resolution to unveil faint old GCs. The properties of relatively old GCs in BCGs may provide important information about the nature of the host galaxy.
An unbiased survey of star forming dwarf galaxies with HST to characterize the frequency of star clusters, their colors and host galaxy properties, would allow to quantitatively study cluster formation in low mass galaxies. The Malkan et al. (1998) survey do not fullfil these criteria but show that such a program could be very rewarding. A better understanding of the ultimate fate of SSCs is also required. Dynamical mass estimates are still rare and often result in masses of the right order of magnitude but surprisingly low mass to light ratios (see Smith & Gallagher 2000). If this is due to flat or top heavy IMFs it would make it harder for young SSCs to survive and become GCs.
#### Acknowledgments.
I would like to thank Kelsey Johnson and her collaborators for allowing me to discuss their work on He2-10 ahead of publication. I am also grateful to the organizers for arranging such a nice workshop. This work made frequent use of the NASA/IPAC extragalactic database (NED).
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## Discussion
U. Fritze von Alvensleben: I wonder if the star clusters you see in dwarf galaxies could really be young globular clusters. Burst strengths and star formation efficiencies found for samples of blue compact dwarfs are $`<`$ 1 %, whereas GC formation modelling by Brown et al. (1995, ApJ 440, 666) requires star formation efficiencies of $`10`$%.
G. Östlin: What matters is really the local star formation efficiency (SFE) and not e.g. how large a fraction of the gas is converted into stars on a global scale. Moreover, the SFE limit ($`10`$%) you refer to is the one required for a young GC to survive the Galactic tidal field, however tidal fields are weaker in low mass galaxies. But of course, formation of numerous young globular clusters would require high global SF efficiencies and burst strengths. This is also the case for the examples I am discussing, which all are galaxies with strong starbursts (i.e. the gas depletion timescale, and/or the timescale to build up the observed stellar mass is much shorter than a Hubble time). It is true that many galaxies termed “blue compact” in reality have very modest star formation rates, but it is not in these galaxies we predominantly find young GC candidates.
J. Gallagher: If most dwarfs make SSCs in reasonable numbers, why do we then see so few in WFPC2 images of “normal” dI galaxies? - Isn’t this a problem?
G. Östlin: Normal dIs have mostly “normal” star formation rates with respect to the mass of the galaxy and are not experiencing starbursts (in terms of star formation timescales, see previous comment). Moreover, their luminosities are often so low that we would not expect many SSCs if they had the same specific frequency of SSCs as e.g. ESO338-IG04 or giant mergers. Nevertheless, it may be that formation of bound clusters require special conditions (like high pressure or gas densities) which are not fulfilled in dwarfs except in the case of external triggers, e.g. mergers and strong interactions. The SSC rich dwarfs discussed above have perturbed morphology and signs of mergers and/or interactions.
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# Radiative corrections for the correlator of (0++,1-+) light hybrid currents
## 1 Introduction
Mesons with exotic quantum number have attracted a great deal of attention in low energy strong interaction physics. They constitute another intrinsic construction of matter beyond the quark model in QCD. Although these mesons have not been confirmed yet, recent experiments indeed give some evidence for their possible existence. The E852 Collaboration at BNL has reported a $`J^{PC}=1^+`$ isovector resonance $`\widehat{\rho }(1405)`$ in the reaction $`\pi ^{}p\eta \pi ^0n`$, with a mass of $`1370\pm 16_{30}^{+50}`$ MeV and a width of $`385\pm 40_{105}^{+65}`$ MeV . This state appears to have been confirmed by the Crystal Barrel Collaboration in $`p\overline{p}`$ annihilation with a mass of $`1400\pm 20\pm 20`$ MeV and a width of $`310\pm 50_{30}^{+50}`$ MeV . E852 lays claim to another $`J^{pc}=1^+`$ isovector state $`\widehat{\rho }(1600)`$ in the reaction $`\pi ^{}p\pi ^+\pi ^{}\pi ^{}p`$, with a mass and width of $`1593\pm 8`$ MeV and $`168\pm 20`$ MeV, resp., which decays into $`\rho \pi `$ . We will have to wait for further confirmation of these states.
The mass values for the exotic $`1^+`$ state reported by the different experiments disagree with most theoretical predictions. The flux-tube model predicts the lowest-lying $`1^+`$ hybrid meson to have a mass of 1.9 GeV , which is consistent with lattice QCD studies which predict the lightest exotic hybrid $`1^+`$ to have a mass of 2.0 GeV . Besides, the flux-tube model also predicts that the $`b_1\pi `$ and $`f_1\pi `$ modes dominate in $`1^+`$ hybrid meson decay , in contradiction with the experiments in which the mode $`\rho \pi `$ is dominant. Differing from the flux-tube model predictions calculations based on the QCD sum rule (QCDSR) approach seem to be closer to experimental results . The preliminary results of QCDSR show that the lightest exotic hybrid meson has a mass around $`1.6`$ GeV and a dominant $`\rho \pi `$ decay mode . This is consistent with the second $`1^+`$ state claimed by E852. The discrepancies between the different model predictions may come from different sources. The flux-tube model uses a non-relativistic linear potential model, which is not so suitable for the light-mass system. On the other hand, many effects, such as higher order terms in the OPE and radiative corrections, may affect the QCDSR predictions. Further theoretical studies are obviously necessary.
In this paper, we calculate the radiative corrections to the current-current correlator of the hybrid current $`\overline{q}(x)\gamma _\nu igG_{\mu \nu }^aT^aq(x)`$. Using this new result, we recalculate the masses of the $`1^+`$ and $`0^{++}`$ hybrids via the standard QCDSR method. We find that, including the radiative corrections, QCDSR have less room to accommodate the recent experimental data.
## 2 Renormalization of the current $`j_\mu =\overline{q}\gamma _\nu igG_{\mu \nu }^aT^aq`$
The operator-mixing problems associated with the renormalization of composite operators have been discussed sometimes ago . For instance, a given gauge invariant operator can mix with other gauge invariant operators, with non-gauge invariant operators which vanish by equations of motion and with operators containing ghosts. The mixing operators must have the same CP quantum number and the same dimension as the original one. In our case, the complete set of operators which can mix with $`\overline{q}\gamma _\nu igG_{\mu \nu }^aT^aq`$ is given by
$`j_\mu ^1`$ $`=`$ $`\overline{q}\gamma _\nu igG_{\mu \nu }q,`$
$`j_\mu ^2`$ $`=`$ $`\overline{q}(\stackrel{}{D}_\mu \stackrel{}{D/}D/D_\mu )q,`$
$`j_\mu ^3`$ $`=`$ $`\overline{q}(\gamma _\mu \sigma _{\alpha \beta }gG_{\alpha \beta }gG_{\alpha \beta }\sigma _{\alpha \beta }\gamma _\mu )q,`$
$`j_\mu ^4`$ $`=`$ $`\overline{q}(\gamma _\mu \stackrel{}{D/}\stackrel{}{D/}D/D/\gamma _\mu )q,`$ (1)
$`j_\mu ^5`$ $`=`$ $`\overline{q}(\gamma _\mu igA/\stackrel{}{D/}+D/igA/\gamma _\mu )q,`$
$`j_\mu ^6`$ $`=`$ $`\overline{q}(igA_\mu \stackrel{}{D/}+D/igA_\mu )q.`$
Note that there is no dimension-five operator containing ghost fields in the set (2). In (2) we have defined $`\sigma _{\alpha \beta }=\frac{i}{2}[\gamma _\alpha ,\gamma _\beta ]`$ and the covariant derivatives $`\stackrel{}{D}_\mu =\stackrel{}{}_\mu +igA_\mu `$, $`D_\mu =_\mu igA_\mu `$, which act on right and left fields, respectively. The fields and couplings in (2) are bare. Only $`j_\mu ^1`$ and $`j_\mu ^3`$ are physical currents. The other currents correspond to so-called nuisance operators which vanish by the equations of motion.
Renormalizing the composite operator corresponding to $`j_\mu ^1`$ we obtain
$$[j_\mu ^1]=Z_1j_\mu ^1+Z_2j_\mu ^2+Z_3j_\mu ^3+Z_4j_\mu ^4+Z_5j_\mu ^5+Z_6j_\mu ^6,$$
(2)
In order to determine the coefficients $`Z_i(i=1,6)`$, we insert the currents into 1PI diagrams and extract the ultraviolet divergences. We find that the divergences associated with the $`\overline{q}qg`$-vertex are sufficient to determine all counterterms. The relevant diagrams are shown in Fig.1. In Feynman gauge, we obtain
$`Z_1`$ $`=`$ $`1+{\displaystyle \frac{113}{18}}{\displaystyle \frac{g^2}{16\pi ^2}}{\displaystyle \frac{1}{ϵ}},`$
$`Z_2`$ $`=`$ $`{\displaystyle \frac{4}{9}}{\displaystyle \frac{g^2}{16\pi ^2}}{\displaystyle \frac{1}{ϵ}},`$
$`Z_3`$ $`=`$ $`{\displaystyle \frac{49}{72}}{\displaystyle \frac{g^2}{16\pi ^2}}{\displaystyle \frac{1}{ϵ}},`$
$`Z_4`$ $`=`$ $`{\displaystyle \frac{8}{9}}{\displaystyle \frac{g^2}{16\pi ^2}}{\displaystyle \frac{1}{ϵ}},`$ (3)
$`Z_5`$ $`=`$ $`{\displaystyle \frac{19}{72}}{\displaystyle \frac{g^2}{16\pi ^2}}{\displaystyle \frac{1}{ϵ}},`$
$`Z_6`$ $`=`$ $`{\displaystyle \frac{35}{36}}{\displaystyle \frac{g^2}{16\pi ^2}}{\displaystyle \frac{1}{ϵ}},`$
where we use dimensional regularization.
## 3 Next-to leading order correction to the current-current correlator
Let us consider the current-current correlator
$`\mathrm{\Pi }_{\mu \nu }(q^2)`$ $`=`$ $`i{\displaystyle d^4xe^{iqx}0|T\{j_\mu (x),j_\nu ^+(0)\}|0}`$
$`=`$ $`({\displaystyle \frac{q_\mu q_\nu }{q^2}}g_{\mu \nu })\mathrm{\Pi }_v(Q^2)+{\displaystyle \frac{q_\mu q_\nu }{q^2}}\mathrm{\Pi }_s(q^2)`$
where $`j_\mu (x)=\overline{u}(x)\gamma _\nu igG_{\mu \nu }^aT^au(x)`$. The invariants $`\mathrm{\Pi }_v(Q^2)`$ and $`\mathrm{\Pi }_s(Q^2)`$ correspond to the contributions from $`1^+`$ and $`0^{++}`$ states and their excited states, respectively. The leading order contribution to (3) including the quark and gluon condensate contributions has already been given in
$`\mathrm{\Pi }_v^0(q^2)`$ $`=`$ $`[{\displaystyle \frac{\alpha _s}{240\pi ^3}}(q^2)^3+{\displaystyle \frac{1}{36\pi }}q^2(\alpha _sG^2+8\alpha _sm\overline{u}u)]\mathrm{ln}({\displaystyle \frac{q^2}{u^2}})`$
$`+\left[{\displaystyle \frac{4\pi }{9}}\alpha _s\overline{u}u^2+{\displaystyle \frac{1}{192\pi ^2}}g^3G^3{\displaystyle \frac{83\alpha _s}{1728\pi }}m\overline{u}Gu\right],`$
$`\mathrm{\Pi }_s^0(q^2)`$ $`=`$ $`\left[{\displaystyle \frac{\alpha _s}{480\pi ^3}}(q^2)^3+({\displaystyle \frac{\alpha _s}{3\pi }}m\overline{u}u+{\displaystyle \frac{\alpha _sG^2}{24\pi }})q^2+{\displaystyle \frac{m^2}{8\pi }}\alpha _sG^2+{\displaystyle \frac{11\alpha _s}{72\pi }}m\overline{u}Gu\right]\mathrm{ln}({\displaystyle \frac{q^2}{u^2}})`$
$`+8\pi \alpha _s\overline{u}u^2.`$
There is no difference for isovector and isoscalar currents in this order. The next-to-leading order correction to the perturbative part of $`\mathrm{\Pi }_v(q^2)`$ and $`\mathrm{\Pi }_s(q^2)`$ can be obtained by calculating the Feynman diagrams in Fig.2, where Figs.2m and 2n only contribute to the isoscalar states. The technique of the calculation which we use here was firstly proposed in .
Let us now briefly comment on the radiative corrections to the correlator of the current $`g\overline{q}\gamma _5\gamma _\nu iG_{\mu \nu }^aT^aq`$ with $`0^{}`$ and $`1^+`$ quantum numbers. The isovector current correlator has the same radiative correction as the isovector current correlator of the ($`0^{++}`$, $`1^+`$) current. However, the isoscalar current correlator has different radiative corrections since diagrams Figs.2m and 2n, which correspond to the mixing with pure gluonic states, now give zero contributions.
Each of the diagrams Figs.2i-2n is gauge-parameter independent by itself because the current is antisymmetric in the Lorentz indices. We have checked on gauge invariance for the sum of diagrams Figs.2a-2h by doing the calculation in a general covariant gauge and found that the result is gauge-parameter independent. Explicit Feynman gauge results for the diagrams are listed in the appendix.
In the present application only the isovector currents are of interest. In the $`\overline{\mathrm{MS}}`$-scheme, the next-to leading corrections to the correlator of isovector currents is given by
$`\mathrm{\Pi }_v^{1a}(q^2)`$ $`=`$ $`{\displaystyle \frac{\alpha _s(\mu )}{240\pi ^3}}(q^2)^3\mathrm{ln}({\displaystyle \frac{q^2}{\mu ^2}})\left[({\displaystyle \frac{53}{4}}{\displaystyle \frac{76}{45}}n_f){\displaystyle \frac{\alpha _s(\mu )}{\pi }}({\displaystyle \frac{35}{24}}{\displaystyle \frac{1}{4}}n_f){\displaystyle \frac{\alpha _s(\mu )}{\pi }}\mathrm{ln}({\displaystyle \frac{q^2}{\mu ^2}})\right]`$ (6)
$`\mathrm{\Pi }_s^{1a}(q^2)`$ $`=`$ $`{\displaystyle \frac{\alpha _s(\mu )}{480\pi ^3}}(q^2)^3\mathrm{ln}({\displaystyle \frac{q^2}{\mu ^2}})\left[({\displaystyle \frac{2017}{216}}{\displaystyle \frac{229}{180}}n_f){\displaystyle \frac{\alpha _s(\mu )}{\pi }}({\displaystyle \frac{35}{24}}{\displaystyle \frac{1}{4}}n_f){\displaystyle \frac{\alpha _s(\mu )}{\pi }}\mathrm{ln}({\displaystyle \frac{q^2}{\mu ^2}})\right]`$
where we take the light quark mass to be zero for convenience. Although (3) include condensates which are proportional to the light quark mass, their contributions are very small compared to those of operators with the same dimension. Throughout this calculation taking zero quark mass is a good approximation.
Eq.(6) does not contain the complete next-to leading order correction. one also has to include the contribution from the renormalization of the current. By inserting the renormalized currents (2) into the correlator
$$id^4xe^{iqx}2Z_i^{\alpha _s}0|T\{j_\mu ^1(x),j_\nu ^{i+}(0)\}|0,$$
(7)
we obtain
$`\mathrm{\Pi }_v^{1b}(q^2)`$ $`=`$ $`{\displaystyle \frac{\alpha _s(\mu )}{240\pi ^3}}(q^2)^3\mathrm{ln}({\displaystyle \frac{q^2}{\mu ^2}})(({\displaystyle \frac{91}{16}}+{\displaystyle \frac{39}{40}}n_f){\displaystyle \frac{\alpha _s(\mu )}{\pi }}+({\displaystyle \frac{35}{36}}{\displaystyle \frac{1}{6}}n_f){\displaystyle \frac{\alpha _s(\mu )}{\pi }}\mathrm{ln}({\displaystyle \frac{q^2}{\mu ^2}}))`$ (8)
$`\mathrm{\Pi }_s^{1b}(q^2)`$ $`=`$ $`{\displaystyle \frac{\alpha _s(\mu )}{480\pi ^3}}(q^2)^3\mathrm{ln}({\displaystyle \frac{q^2}{\mu ^2}})(({\displaystyle \frac{679}{144}}+{\displaystyle \frac{97}{120}}n_f){\displaystyle \frac{\alpha _s(\mu )}{\pi }}+({\displaystyle \frac{35}{36}}{\displaystyle \frac{1}{6}}n_f){\displaystyle \frac{\alpha _s(\mu )}{\pi }}\mathrm{ln}({\displaystyle \frac{q^2}{\mu ^2}})).`$
In Eq.(7) we have introduced the notation $`Z_1^{\alpha _s}=Z_1Z_g1`$ and $`Z_i^{\alpha _s}=Z_i(i=2,\mathrm{..6})`$. For the original hybrid current $`j_\mu ^1`$ one also has to include the counterterm $`Z_g`$ resulting from the renormalization of the bare coupling constant $`g`$. The final result for the radiative corrections to the correlator is finally obtained by adding the contributions of (6) and (8).
## 4 Sum rules for $`1^+`$ and $`0^{++}`$ hybrid mesons
For sum rule applications we need the spectral density associated with the current-current correlator. The spectral density $`\rho _v(s)=Im\mathrm{\Pi }_v(s)`$ is defined via the standard dispersion relation
$$\mathrm{\Pi }_v(q^2)=\frac{(q^2)^n}{\pi }_0^{\mathrm{}}𝑑s\frac{\rho _v(s)}{s^n(sq^2)}+\underset{k=0}{\overset{n1}{}}a_k(q^2)^k,$$
(9)
where the $`a_k`$ are appropriate subtraction constants to render Eq.(9) finite.
From (3), (6) and (8) we obtain
$`\rho _v(s)`$ $`=`$ $`{\displaystyle \frac{\alpha _s(\mu )}{240\pi ^2}}s^3(1+{\displaystyle \frac{1301}{240}}{\displaystyle \frac{\alpha _s(\mu )}{\pi }}{\displaystyle \frac{17}{36}}{\displaystyle \frac{\alpha _s(\mu )}{\pi }}\mathrm{ln}({\displaystyle \frac{s}{\mu ^2}}))+{\displaystyle \frac{1}{36}}s(\alpha _sG^2+8\alpha _sm\overline{u}u)`$ (10)
$`\rho _s(s)`$ $`=`$ $`{\displaystyle \frac{\alpha _s(\mu )}{480\pi ^2}}s^3(1+{\displaystyle \frac{6979}{2160}}{\displaystyle \frac{\alpha _s(\mu )}{\pi }}{\displaystyle \frac{17}{36}}{\displaystyle \frac{\alpha _s(\mu )}{\pi }}\mathrm{ln}({\displaystyle \frac{s}{\mu ^2}}))`$
$`({\displaystyle \frac{\alpha _s}{3}}m\overline{u}u+{\displaystyle \frac{\alpha _sG^2}{24}})s{\displaystyle \frac{m^2}{8}}\alpha _sG^2{\displaystyle \frac{11\alpha _s}{72}}m\overline{u}Gu,`$
where we have set $`n_f=3`$.
On the other hand, the spectral density is saturated by narrow physical resonances and the continuum. We therefore write
$$(\frac{q_\mu q_\nu }{q^2}g_{\mu \nu })\rho _v(s)=\underset{R}{}0|j_\mu |RR|j_\nu |0\pi \delta (sm_R^2)+continuum,$$
(12)
where we have assumed that the mass $`m_R`$ is much larger than the decay width of the hybrid, so that the imaginary part of the propagator has been replaced by $`\pi \delta (sm_R^2)`$.
In order to extract information on the lowest-lying resonance, it is usually assumed that the lowest-lying resonance dominates the spectral density. The contribution of higher excited states can be suppressed by applying the Borel transformation $`\widehat{L}_M`$ to both sides of Eq. (9). One thus has
$$R_0(\tau )=M\widehat{L}_M\mathrm{\Pi }_v(q^2)=\frac{1}{\pi }_0^{\mathrm{}}e^{s/M}\rho _v(s)𝑑s.$$
(13)
The upper limit of the intergral can be replaced by a finite number $`s_0`$. The contributions beyond the threshold $`s_0`$ are considered to result from the continuum. $`R_0(\tau )`$ is the zeroth moment. Higher order moments are defined by $`R_k=(M^2\frac{}{M})^kR_0(M)`$. Resonance masses can be obtained by taking the ratio $`m_R^2=\frac{R_{k+1}}{R_k}`$ with the assumption that only a single narrow resonance dominates. In principle any value of $`k`$ can be chosen to determine the resonance masses. However, since we have to truncate the series of the power expansion, using higher moments will damage the convergence of the OPE. Besides, it is also arbitrary to define the scalar function $`\mathrm{\Pi }_v^k(q^2)=\frac{1}{(q^2)^k}\mathrm{\Pi }_v(q^2)`$ by extracting a general tensor factor
$$(q^2)^k(\frac{q_\mu q_\nu }{q^2}g_{\mu \nu })$$
(14)
from the correlator $`(\text{3})`$. Although (13) can be considered as a higher moment of $`\mathrm{\Pi }_v^k(q^2)`$, to the order of OPE that we are considering, the sum rules for $`\mathrm{\Pi }_v^k(q^2)`$ and $`\mathrm{\Pi }_v(q^2)`$ are obviously different. For instance, as pointed out in , the dimension-six operators of (3) do not contribute to sum rule (13), while they play an important role in stabilizing the sum rule of $`\mathrm{\Pi }_v^k(q^2)`$ in the case of $`k=1`$. We shall consider the two cases $`k=0`$ and $`k=1`$ in turn.
The single particle matrix elements contributing to (12) are parametrized as
$`0|j_\mu |V`$ $`=`$ $`iϵ_\mu f_vm_v^3.`$ (15)
$`0|j_\mu |S`$ $`=`$ $`ip_\mu f_sm_s^2.`$
In the narrow resonance approximation the sum rules are independent of the matrix element (15). When the decay width of the resonance is comparable with its mass, the approximation (12) is no longer valid. Then the information contained in (15) may become important. We will comment on this later on. By using (10)-(15) we obtain
$$m_{v,s}^2=\frac{_0^{s_0}e^{s\tau }s\rho _{v,s}(s)𝑑s}{_0^{s_0}e^{s\tau }\rho _{v,s}(s)𝑑s}$$
(16)
for the sum rule (13) and
$`m_v^2`$ $`=`$ $`{\displaystyle \frac{_0^{s_0}e^{s\tau }\rho _v(s)𝑑s}{_0^{s_0}e^{s\tau }\rho _v(s)𝑑s/s\frac{4\pi ^2}{9}\alpha _s\overline{u}u^2\frac{1}{192\pi }g^3G^2+\frac{83}{1728}\alpha _sm\overline{u}Gu}}`$ (17)
$`m_s^2`$ $`=`$ $`{\displaystyle \frac{_0^{s_0}e^{s\tau }\rho _s(s)𝑑s}{_0^{s_0}e^{s\tau }\rho _s(s)𝑑s/s8\pi ^2\alpha _s\overline{u}u^2}}`$
for the sum rule for $`\mathrm{\Pi }_v^k(q^2)`$ with $`k=1`$ , where the spectral densities $`\rho _{v,s}(s)`$ are given in (10). The various parameters entering in (16)-(17) are specified as
$`\mathrm{\Lambda }_{QCD}=0.25GeV,`$ $`m=0.01GeV,`$ $`m\overline{u}u={\displaystyle \frac{1}{4}}f_\pi ^2m_\pi ^2`$ (18)
$`\alpha _sG^2=0.04Gev^4,`$ $`g^3G^3=1.1GeV^2\alpha _sG^2,`$ $`f_\pi =0.132GeV`$
$`\alpha _s(\mu )={\displaystyle \frac{4\pi }{9\mathrm{ln}(\frac{\mu ^2}{\mathrm{\Lambda }_{QCD}^2})}},`$ $`\mu =2GeV,`$ $`g\overline{u}Gu=1.5Gev^2\overline{u}u`$
In Fig.3 and Fig.4 we show a mass plot of the $`1^+`$ state in its dependence on the Borel parameter $`M`$ in the two sum rule ratios. The second sum rule gives a smaller mass which we will take as the lower bound.
The sensitivity of the mass to the choice of the threshold value $`s_0`$ is obvious. The mass will go to infinity when both $`s_0`$ and $`M`$ go to infinity, because, when $`M`$ goes to infinity, the Borel measure (13) no longer suppresses the continuum. In order to give a reasonable estimate, we set $`s_0`$ around $`4`$ GeV<sup>2</sup>. Fig.5 shows that $`m_v`$ takes values around $`3`$ GeV when $`s_0`$ is set to infinity. In Fig.6 one cannot find any stable point in the two-dimensional ($`s_0`$,$`M`$) space. Therefore, before fixing $`s_0`$, we cannot make any precise prediction for the hybrid mass. However, if one believes that the $`1^+`$ hybrid mass lies around 2 $`GeV`$, $`s_0`$ should be larger than 4 $`GeV^2`$. This results in a lower bound for the $`1^+`$ mass of 1.55 GeV(see Fig.2). The radiative corrections enhance the lower bound which gives QCDSR less room to accommodate recent experimental data.
The prediction of the mass is also sensitive to the form of the spectral density. This is mostly due to a truncated OPE. The contributions from higher dimension operators are very likely not small. The uncertainty from the narrow resonance approximation in (12) does not seem to reduce this discrepancy. This can be checked by replacing the narrow resonance form $`\pi \delta (sm_R^2)`$ in (12) by the Breit-Wigner form
$$\frac{\mathrm{\Gamma }_Rm_R}{(sm_R^2)^2+\mathrm{\Gamma }^2m_R^2}.$$
(19)
$`\mathrm{\Gamma }_R`$ is the width of the $`1^+`$ hybrid. When we choose $`\mathrm{\Gamma }_R=200MeV`$ and a parametrization of the matrix element (15) in the form $`0|j_\mu |V=iϵ_\mu m_R\overline{f}_vs^x`$ with $`\overline{f}_v`$ constant and $`x=0.5÷1.5`$, the hybrid mass is not sensitive to using the full Breit-Wigner form propagator. We show the change in the Figs.3 and 4 by dotted lines. The discrepancy between the two cases k=0 and k=1 is still big and the mass predictions become somewhat larger.
The sum rule for the $`0^{++}`$ hybrid is shown in Figs.7 and 8 for the two cases, k=0 and k=1 resp., where we set $`s_0=7`$ GeV<sup>2</sup>. The radiative corrections reduce the discrepancy between the two cases. Similar to the $`1^+`$ case, when $`s_0`$ goes to infinity, the prediction of mass is around $`3`$ GeV. It means that the contribution of the continuum dominates over that of the resonance s. Therefore, the value chosen for the threshold $`s_0`$ is important.
## 5 Summary
In summary, we have calculated the next-to-leading order corrections to the two point correlator of the current $`g\overline{q}\gamma _\nu iG_{\mu \nu }^aT^aq(x)`$. We recalculated the masses of the $`1^+`$ and $`0^{++}`$ hybrids. We find that the radiative corrections reduce the lower bound of the $`1^+`$ mass and leave less room for QCD sum rules to fit the recent experimental data.
Note added in proof: While preparing this paper for publication, we became aware of a recent paper by K. Chetyrkin and S. Narison which addresses similar problems.
Acknowledgment We would like to thank S. Groote, A. A. Pivovarov and K. Chekyrkin for very useful discussions. We would also like to thank K. Chekyrkin for providing us with intermediate results of the calculation. The work of H.Y. J. is supported by the Alexander von Humboldt foundation.
Appendix In this appendix we list the results of calculating the diagrams Fig.2 in the Feynman gauge for the correlator (3).
$`Fig.2a`$ $`:`$ $`C\left[({\displaystyle \frac{637}{3840}}{\displaystyle \frac{3}{128}}\mathrm{ln}({\displaystyle \frac{q^2}{\mu ^2}}))g_{\mu \nu }+({\displaystyle \frac{143}{640}}+{\displaystyle \frac{9}{256}}\mathrm{ln}({\displaystyle \frac{q^2}{\mu ^2}})){\displaystyle \frac{q_\mu q_\nu }{q^2}}\right]`$
$`Fig.2b`$ $`:`$ $`C\left[({\displaystyle \frac{673}{172800}}+{\displaystyle \frac{1}{1920}}\mathrm{ln}({\displaystyle \frac{q^2}{\mu ^2}}))g_{\mu \nu }+({\displaystyle \frac{9}{1600}}{\displaystyle \frac{1}{1280}}\mathrm{ln}({\displaystyle \frac{q^2}{\mu ^2}})){\displaystyle \frac{q_\mu q_\nu }{q^2}}\right]`$
$`Fig.2c`$ $`:`$ $`C\left[{\displaystyle \frac{1}{480}}g_{\mu \nu }+({\displaystyle \frac{307}{3840}}+{\displaystyle \frac{3}{256}}\mathrm{ln}({\displaystyle \frac{q^2}{\mu ^2}})){\displaystyle \frac{q_\mu q_\nu }{q^2}}\right]`$
$`Fig.2d`$ $`:`$ $`C\left[({\displaystyle \frac{157}{1600}}+{\displaystyle \frac{9}{640}}\mathrm{ln}({\displaystyle \frac{q^2}{\mu ^2}}))g_{\mu \nu }+({\displaystyle \frac{867}{6400}}{\displaystyle \frac{27}{1280}}\mathrm{ln}({\displaystyle \frac{q^2}{\mu ^2}})){\displaystyle \frac{q_\mu q_\nu }{q^2}}\right]`$
$`Fig.2e`$ $`:`$ $`C\left[({\displaystyle \frac{583}{21600}}+{\displaystyle \frac{1}{240}}\mathrm{ln}({\displaystyle \frac{q^2}{\mu ^2}}))g_{\mu \nu }+({\displaystyle \frac{31}{800}}{\displaystyle \frac{1}{160}}\mathrm{ln}({\displaystyle \frac{q^2}{\mu ^2}})){\displaystyle \frac{q_\mu q_\nu }{q^2}}\right]`$
$`Fig.2f`$ $`:`$ $`C\left[{\displaystyle \frac{1}{640}}g_{\mu \nu }+{\displaystyle \frac{1}{640}}{\displaystyle \frac{q_\mu q_\nu }{q^2}}\right]`$
$`Fig.2(g+h+i)`$ $`:`$ $`C[({\displaystyle \frac{79}{1440}}{\displaystyle \frac{19}{2700}}n_f({\displaystyle \frac{1}{128}}{\displaystyle \frac{1}{960}}n_f)\mathrm{ln}({\displaystyle \frac{q^2}{\mu ^2}}))g_{\mu \nu }`$
$`+({\displaystyle \frac{97}{1280}}+{\displaystyle \frac{31}{3200}}n_f+({\displaystyle \frac{3}{256}}{\displaystyle \frac{1}{640}}n_f)\mathrm{ln}({\displaystyle \frac{q^2}{\mu ^2}})){\displaystyle \frac{q_\mu q_\nu }{q^2}}]`$
$`Fig.2j`$ $`:`$ $`C\left[({\displaystyle \frac{1}{24}}+{\displaystyle \frac{1}{144}}\mathrm{ln}({\displaystyle \frac{q^2}{\mu ^2}}))g_{\mu \nu }+({\displaystyle \frac{25}{216}}{\displaystyle \frac{1}{48}}\mathrm{ln}({\displaystyle \frac{q^2}{\mu ^2}})){\displaystyle \frac{q_\mu q_\nu }{q^2}}\right]`$
$`Fig.2k`$ $`:`$ $`C\left[{\displaystyle \frac{1}{720}}g_{\mu \nu }+{\displaystyle \frac{13}{3240}}{\displaystyle \frac{q_\mu q_\nu }{q^2}}\right]`$
$`Fig.2l`$ $`:`$ $`C\left[({\displaystyle \frac{83}{28800}}{\displaystyle \frac{1}{1920}}\mathrm{ln}({\displaystyle \frac{q^2}{\mu ^2}}))g_{\mu \nu }+({\displaystyle \frac{289}{86400}}{\displaystyle \frac{1}{1920}}\mathrm{ln}({\displaystyle \frac{q^2}{\mu ^2}})){\displaystyle \frac{q_\mu q_\nu }{q^2}}\right]`$
$`Fig.2m`$ $`:`$ $`C\left[{\displaystyle \frac{1}{8640}}g_{\mu \nu }{\displaystyle \frac{1}{1440}}{\displaystyle \frac{q_\mu q_\nu }{q^2}}\right]`$
$`Fig.2n`$ $`:`$ $`C\left[({\displaystyle \frac{71}{5400}}{\displaystyle \frac{1}{480}}\mathrm{ln}({\displaystyle \frac{q^2}{\mu ^2}}))g_{\mu \nu }+({\displaystyle \frac{401}{21600}}+{\displaystyle \frac{1}{320}}\mathrm{ln}({\displaystyle \frac{q^2}{\mu ^2}})){\displaystyle \frac{q_\mu q_\nu }{q^2}}\right]`$
We use the abbreviation $`C=\frac{\alpha _s(\mu )^2}{\pi ^4}(q^2)^3\mathrm{ln}(\frac{q^2}{\mu ^2})`$.
Figure captions
Fig. 1 Feynman diagrams for the renormalization of the hybrid current. Dots stand for the current vertices.
Fig. 2. Feynman diagrams for the next-to-leading calculation. Dots stand for the current vertices.
Fig. 3. $`1^+`$ hybrid mass $`m_v`$ versus Borel variable $`M`$ for the first sum rule Eq.(16). Dashed line gives the result for the leading order calculation. Solid line includes radiative corrections. Dotted line gives the results when using the Breit-Wigner resonance propagator.
Fig. 4. $`1^+`$ hybrid mass $`m_v`$ versus Borel variable $`M`$ for the second sum rule Eq.(17)
Fig. 5. $`1^+`$ hybrid mass $`m_v`$ versus Borel variable $`M`$ for the second sum rule Eq.(17) when $`s_0`$ goes to infinity.
Fig. 6. Three-dimensional figure of $`1^+`$ hybrid mass $`m_v`$ vs. the Borel variable $`M`$ and the threshold parameter $`s_0`$.
Fig. 7. $`0^{++}`$ hybrid mass $`m_s`$ versus Borel variable $`M`$ for the first sum rule Eq.(16). Solid line inculdes radiative corrections.
Fig. 8. $`0^{++}`$ hybrid mass $`m_s`$ versus Borel variable $`M`$ for the second sum rule Eq.(17). Solid line includes radiative correction.
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# 1. Introduction
## 1. Introduction
It is now widely accepted that electromagnetic duality provides a powerful and useful new principle which is valid in a large class of physically interesting quantum field theories in a Minkowski space-time of four dimensions. The prototype for the quantum version of this idea was furnished by a proposal of Montonen and one of the present authors, \[Montonen and Olive 1977\], who considered the context of a special sort of spontaneously broken $`SU(2)`$ gauge theory, later realised to be naturally supersymmetric \[D’Adda, Di Vecchia and Horsley 1978, Witten and Olive 1978, Osborn 1979\]. This can be regarded as a semi-realistic theory of unified particle interactions as it is the same sort of theory as the standard model even though it differs in some crucial respects. The manner in which duality is realised on particle states requires magnetically charged states arising as solitons solutions and, in addition, quantum bound state of these \[Sen 1994\].
This acceptance, which now extends to superstring theories, where solitons occur as higher branes, has been achieved despite the fact that no sort of proof has been found, even in the case of the $`SU(2)`$ gauge theory with the highest allowable degree of supersymmetry, $`N=4`$. It is this situation for which most evidence has steadily accumulated that the idea is exactly true. This supporting evidence has been facilitated bythe acquisition of new mathematical techniques that have enhanced our understanding of quantum field theory. For example, the Atiyah-Singer index theorem, which is related to the theory of axial anomalies, plays a ubiquitous role. A modifed version of this theorem is crucial in determining both classical and quantum properties of self-dual monopole solutions.
Since the quantum electromagnetic duality transformations combine to form a group related to the modular group (or its generalisations) it is reasonable to expect some version of the theory of modular forms to become increasingly important. Conversely the idea of electromagnetic duality has led to breakthroughs in the classification theory of four-manifolds (playing the role of Euclidean space-times for twisted supersymmetric gauge theories).
It would be nice to have a simpler, toy model which, although less realistic physically, could, in compensation, be more tractable mathematically. Indeed, such a theory exists. It is simply free Maxwell theory with only putative couplings, not realised in practice. Thus particle states possess what could be called a “Cheshire cat” existence, that will be become clearer later. Such a theory would be too trivial on flat space time and, in order to obtain some worthwhile structure, the ambient space-time manifold has to be allowed to be fairly general. It is taken to be smooth, compact and oriented. Hence it is a four-manifold which obeys the topological symmetry known as Poincaré duality, and this will be a crucial ingredient forestablishing electromagnetic duality here. This toy model follows the original proposals by E Verlinde and E Witten in 1995.
What is interesting about this model is that much of the same sort of mathematical structure as mentioned above again comes into play. In particular there enter modular invariant theta functions of a rather general nature. The Atiyah-Singer index theorem plays a mysterious role interlocked with the modular group and its subgroups of index three and the nature of space-time and its possible spin structures. This is despite the significant differences between the two situations (no hint of supersymmetry, curved Riemannian space-time in one case, $`N=4`$ supersymmetry, flat, non-compact Minkowski space-time in the other). This points to the conclusion that electromagnetic duality is indeed a rather general phenomenon.
Moreover it is likely that this picture extends to any space-time with Minkowski metric and dimension which is a multiple of four. If this dimension is denoted $`4k`$ then the putative coupling of Maxwell field strengths with charged particles is replaced by a coupling of a $`2k`$-form field strength to $`(2k2)`$-branes.
In section 2 we start with the naive idea of electromagnetic duality as a classical symmetry of the energy-momentum tensor of Maxwell theory in Minkowski space with respect to rotations between the electric and magnetic fields. It is shown how this idea can be extended to a larger group of $`SL(2,𝐑)`$ transformations acting on the energy momentum tensor. It is explained how this leads in turn to consideration of the Feynman path integral, of, rather surprisingly, the exponential of the Euclidean action. A special case of this Euclidean path integral is the Minkowski space partition function.
The effect of the Dirac quantisation condition for magnetic fluxes exhibits some extra subtleties in space-times of four dimensions, particularly when complex spinor wave functions are considered. This is reviewed in section 3. The effect is to break the continuous $`SL(2,𝐑)`$ group to a discrete subgroup related to the modular group in a way that is explained in sections 4 and 5. It is noteworthy that nothing like the Zwanziger-Schwinger quantisation condition for dyonic charges plays any role.
Owing to the complicated topology of the four-manifold of space-time, the possible magnetic fluxes are related to a lattice formed by the free part of the second homology group. This lattice is unimodular by virtue of Poincaré duality and its even or oddness properties are related to the presence or absence of spin structures on the space-time four manifold. In fact four-manifolds can be separated into three distinct types whose properties are described in section 3.
In section 4 properties of integral lattices are reviewed. Starting from an odd unimodular lattice a general construction is given of an even integral lattice, leading sometimes to new unimodular lattices which can be both odd or even. Furthermore, also in section 4 the relevance to the Dirac quantisation condition and the question of spin structures is explained.
In section 5, following the arguments of E Verlinde and E Witten , the “extended partition functions” are evaluated explicitly using the Dirac quantisation conditions and the semiclassical method. The results are proportional to generalised theta functions associated with the unimodular lattice formed by the free part of the second homology group. Particular attention is paid to the different possibilities afforded by the compatibility of either scalar or spinor complex wave functions on four-manifolds of type II, i.e. when the relevant lattice is odd.
In section 6 a more general construction is presented that associates theta functions with any integral lattice, not necessarily unimodular, whether or not the scalar product is positive definite. Such integral lattices are contained as subgroups of the lattices reciprocal to them and so define a finite number of cosets, to each of which corresponds a theta function. An action of the modular group is defined on these theta functions when the lattice is even and of the Hecke subgroup when the lattice is odd. Careful analysis of the self-consistency of this action furnishes a proof of “Milgram’s formula”, valid for any even integral lattice. This expresses the signature of the lattice, mod eight, in terms of coset properties.
In section 7, the theta function construction of section 6 is applied to the even integral lattice associated with an odd unimodular lattice by the construction of section 4. The result is to associate up to four theta functions with a given odd unimodular lattice, though there are usually linear dependences. When the odd unimodular lattice corresponds to that associated with any type II four-manifold, two of these theta functions enter the distinct Maxwell partition functions for fluxes supporting either scalar or spinor complex wave functions, provided the electric charges they carry coincide. Although each partition function is individually covariant with respect to only a subgroup of $`SL(2,𝐙)`$ of index three, they are related to each other by the missing transformations. In this way the full $`SL(2,𝐙)`$ group of electromagnetic duality transformations is restored for space-time four-manifolds of type II.
## 2. Abelian gauge fields and electromagnetic duality
We shall usually be considering a single abelian Maxwell field strength described, in exterior calculus notation, by a closed two form $`F`$ on a space-time manifold $`_4`$ which is closed, compact, connected, smooth and oriented. Nevertheless much of the argument extends to higher dimensional space-times of the same type, (which will be denoted $`_{4k}`$), as long as their dimension is a multiple of four and the field-strength $`F`$ is a closed $`2k`$-form, that is a mid-form. It would then be a generalised Kalb-Ramond field, \[Kalb and Ramond 1974\], that could couple to the world-volume of a $`2k2`$-brane.
For reasons that will become clear it is important to allow the space-time manifold to be topologically complicated. Sometimes we shall suppose that $`_{4k}`$ be such that it can be endowed with a Minkowski metric (that is with one time component) and this would require that its Euler number $`\chi (_{4k})`$ vanish. We may further require that this metric can always be “Wick rotated” to a Euclidean metric (with no time components) by an analytic continuation. This may impose further constraints on the topological properties of $`_{4k}`$.
With either sort of metric the Hodge star operation, $``$, can be defined, converting $`p`$-forms on $`_{4k}`$ to $`4kp`$-forms. Acting on $`2k`$-forms, that is, mid-forms, the repeated action of the Hodge star operation yields
$$=(1)^t,$$
$`(2.1)`$
where $`t`$ equals the number of time components, that is, one or zero. So, in the Minkowski case, $``$ has eigenvalues $`\pm \mathrm{i}`$ and hence no real eigenfunctions. It is this Minkowski situation, rather than the Euclidean one, in which electromagnetic duality can apply.
The basic idea of a resemblance between the parts of $`F`$ thought of as the electric and magnetic fields, $`\underset{¯}{E}`$ and $`\underset{¯}{B}`$, is very old and was reinforced by Maxwell’s discovery of his equations governing their behaviour in vacuo. “Duality rotations” between $`\underset{¯}{E}`$ and $`\underset{¯}{B}`$ provide a symmetry of the Minkowski energy density $`(E^2+B^2)/2`$ (this particular expression applies when the space-time is flat), and more generally the complete energy momentum tensor. Moreover the rotations map between solutions of the equations even though the action does change.
This idea was extended in the context of supergravity theories by Gaillard and Zumino , building on ideas of Cremmer and Julia . The following action can be defined on any of the space-times mentioned:
$$W=\frac{1}{2\tau _2}__4F\widehat{\tau }F,$$
$`(2.2)`$
where, in Minkowski space,
$$\widehat{\tau }=\tau _1+\tau _2=\frac{\theta }{2\pi }+\frac{2\pi \mathrm{}}{q^2},$$
$`(2.3)`$
so that this action, (2.2), is indeed real. In suitably extended supergravity theories $`\tau `$ will depend upon scalar fields related to the metric tensor by supersymmetry transformations. As far as this paper is concerned, there are no scalar fields and no supersymmetry. $`\tau _1`$ and $`\tau _2`$ are simply dimensionless quantities parametrising the theory. It will be convenient to combine them as real and imaginary parts of a single complex variable
$$\tau =\tau _1+i\tau _2=\frac{\theta }{2\pi }+\mathrm{i}\frac{2\pi \mathrm{}}{q^2}$$
$`(2.4)`$
The conventional Maxwell term is $`\frac{1}{2}FF`$. The other term, called the theta term, has no apparent effect classically as it affects neither the Euler-Lagrange equations, nor the value of the energy-momentum tensor. However it can affect the quantum phase when the topology of the background space-time is sufficiently non-trivial. In circumstances to be explained, $`(\tau _2)^1__4FF`$ is quantised so that the exponentiated quantum action, $`\mathrm{exp}(\frac{\mathrm{i}}{\mathrm{}}W)`$, depends upon the parameter $`\theta `$ in a periodic manner. But the dependence upon $`\theta `$ disappears altogether when the topology is too trivial, that is, when the second Betti number of $`_4`$ vanishes.
Following Gaillard and Zumino we consider the effect of the linear transformations:
$$\widehat{\tau }F\widehat{\tau }^{}F^{}=A\widehat{\tau }F+BF$$
$`(2.5a)`$
$$FF^{}=C\widehat{\tau }F+DF,$$
$`(2.5b)`$
where $`A,B,C`$ and $`D`$ are real constants.
Then the dimensionless complex coupling constant variable, $`\tau `$, (2.4) undergoes the fractional linear transformation
$$\tau \tau ^{}=\frac{A\tau +B}{C\tau +D}.$$
$`(2.6a)`$
while its imaginary part $`\tau _2`$ undergoes
$$\tau _2\tau _2^{}=\frac{\tau _2(ADBC)}{|C\tau +D|^2}.$$
$`(2.6b)`$
The transformations (2.5) evidently provide symmetries of the two equations of motion $`dF=0`$ and $`d\widehat{\tau }F=0`$. Now we consider the effect on the symmetric energy momentum tensor $`T_{\mu \nu }`$, obtained from (2.2) by variation of the metric and written in the Sugawara form,
$$T_{\mu \nu }=\frac{1}{2}(F_{\mu \lambda }g^{\lambda \sigma }F_{\sigma \nu }+^{}F_{\mu \lambda }g^{\lambda \sigma }{}_{}{}^{}F_{\sigma \nu }^{}),$$
$`(2.7)`$
where $`F=F_{\mu \nu }dx^\mu dx^\nu /2`$ and $`{}_{}{}^{}F=^{}F_{\mu \nu }dx^\mu dx^\nu /2`$. The result is
$$T_{\mu \nu }T_{\mu \nu }^{}=|C\tau +D|^2T_{\mu \nu }.$$
If we restrict the transformations (2.5) to the subgroup leaving $`\tau `$ unchanged in (2.6), evidently $`\tau _2`$ is also unchanged and so by (2.6b), the result is
$$T_{\mu \nu }T_{\mu \nu }^{}=(ADBC)T_{\mu \nu }.$$
Hence the energy momentum tensor is invariant under the subgroup if
$$ADBC=1$$
$`(2.8)`$
and this yields a $`U(1)`$ subgroup comprising the duality rotations previously mentioned. The symmetry of the energy momentum tensor (2.7) can be enlarged from this $`U(1)`$ to the full $`SL(2,𝐑)`$ if the transformations (2.5) are modified in the following natural way. Let us substitute for the physical field strengths $`F`$ more geometrical quantities $`G`$, by $`F=\mathrm{}G/q`$. $`G`$ is more geometrical in the sense that its fluxes, unlike those of $`F`$ are dimensionless. In terms of $`G`$ the dimensionless action is given by
$$\frac{W}{\mathrm{}}=\frac{1}{4\pi }__4G\widehat{\tau }G$$
and so is explicitly a function of $`\tau `$. Now the energy momentum tensor (2.7) reads
$$T_{\mu \nu }=\frac{\mathrm{}\tau _2}{4\pi }(G_{\mu \lambda }g^{\lambda \sigma }G_{\sigma \nu }+^{}G_{\mu \lambda }g^{\lambda \sigma }{}_{}{}^{}G_{\sigma \nu }^{}).$$
$`(2.7^{})`$
so that a dependence on $`\tau _2`$ is made explicit. Now substitute $`G`$ for $`F`$ in (2.5), thereby defining new duality transformations, still leading to (2.6). Under these transformations acting on both $`G`$ and $`\tau `$ in the energy momentum tensor $`T_{\mu \nu }`$ in the form (2.7’), the preceding calculations show that $`T_{\mu \nu }`$, is invariant providing only that (2.8) holds.
Thus now the dimensionless complex coupling $`\tau `$ changes whilst preserving the positive nature of $`\tau _2`$. This is appropriate as $`\tau _2`$ is the inverse of the fine structure constant and hence intrinsically positive.
These transformations form the three dimensional non compact group $`SL(2,𝐑)`$, or what is the same by a group theory isomorphism, the symplectic group $`Sp(2,𝐑)`$.
The transformations also map between solutions of the free Maxwell equations and the key question will be to what extent these transformations continue to provide symmetries when there is a possibility of electrically charged particles (or branes) being present, subject to the rules of quantum theory. Because it is the Minkowski energy, rather than the action, which is invariant classically, the natural quantity to consider in the quantum theory is the partition function constructed with this energy:
$$Z(\tau )=Tr\left(\mathrm{e}^{E(\tau )}\right).$$
$`(2.9)`$
It is this partition function that will be the candidate for quantum electromagnetic duality, just as it is the partition function that displays the Kramers-Wannier duality of the Ising model \[Kramers and Wannier 1941\]. Indeed it will be found that (2.9) is invariant under the transformations (2.6) provided they are restricted to a discrete subgroup, isomorphic to the modular group. The discreteness is a consequence of the Dirac quantisation condition that the magnetic fluxes have to satisfy in order to permit complex wave functions.
For rather general, nonlinear dynamical systems with a finite number of degrees of freedom and the property that the action includes only terms quadratic, linear and independent of velocities there is a Feynman path integral expression for the partition function (2.9).
$$Z(\tau )=\mathrm{}\delta A\mathrm{e}^{\frac{\mathrm{i}}{\mathrm{}}W_{EUCLIDEAN}}.$$
$`(2.10)`$
The Euclidean action $`W_{EUCLIDEAN}`$ is obtained from the original action by what can be thought of as a “Wick rotation” whereby velocities are multiplied by $`\mathrm{i}`$ and time by $`\mathrm{i}`$. As a result, $`iW_{EUCLIDEAN}`$ has an imaginary part linear in velocities, and a real part that is negative definite if the original energy is positive. Consequently the path integral is highly convergent. Because of the trace in (2.9), the paths integrated over are closed paths traversed in configuration space in unit time with distinguished end points. This result is known as the Feynman-Kac formula \[Feynman and Hibbs, Feynman\].
The presence of a complex phase factor in (2.10) due to terms linear in velocity appears to contradict the manifest reality of the partition function as defined in (2.9). But this is illusory because the space of closed paths in configuration space that are integrated over possess a $`Z_2`$ symmetry with respect to the interchange of pairs of identical paths differing only in the sense of time evolution along the path. Under this interchange the two contributions to $`\mathrm{exp}\frac{\mathrm{i}}{\mathrm{}}W_{EUCLIDEAN}`$ are related by complex conjugation. As a result the sum of these two complex contributions is indeed real.
Because it is quadratic in field strengths, something similar happens with the more complicated action (2.2) under consideration here. As a result of a similar argument, the Maxwell partition function can be expressed in the form (2.10) where now $`W_{EUCLIDEAN}`$ is obtained from $`W`$, (2.2), by a “Wick rotation” of the metric tensor. This tensor only enters the Maxwell term as the theta term is “topological” and hence independent of the metric, and so unaffected by the Wick rotation. The result is that $`W_{EUCLIDEAN}`$ is given by the same expression as before, (2.2), when it is understood that $`\tau `$ is replaced by
$$\tau _{EUCLIDEAN}=\tau +\mathrm{i}\tau _2=\frac{\theta }{2\pi }+\mathrm{i}\frac{2\pi \mathrm{}}{q^2}.$$
$`(2.11)`$
The metric dependence is encoded in the Hodge $``$ operator which, by (2.1), now has unit square and hence eigenvalues $`\pm 1`$. This has two consequences. One is that $`\mathrm{i}W_{EUCLIDEAN}`$ is complex when the field strengths are real and that its real part is negative definite, thereby ensuring convergence of the integral over gauge potentials $`A`$ in (2.9). The other consequence is that if $``$ is regarded as imaginary in Minkowski space-time and real in Euclidean space, in view of its eigenvalues, then $`\tau `$ has the same complex structure in either case. In this sense it is unaffected by the Wick rotation. Accordingly the complex variable $`\tau `$ given by the expression (2.4) not involving the Hodge $``$ is equally relevant with either metric.
In evaluating this partition function the space-time four manifold has to be considered as $`_4=S_1\times _3`$, where “time” is the coordinate around the circle, (periodic because of the trace), and $`_3`$ an appropriate section of $`_4`$. If the metric on $`S_1\times _3`$ factorises correspondingly it is easy to see that the partition function will be again real because reversing the sense of time around this circle will effect a complex conjugation of the two quantum amplitude contributions.
It will turn out to be highly instructive to consider what, by abuse of terminology, is often also called a “partition function”. This expression is given by the path integral (2.10) but with the integral in the action being over the four manifold given by the full space-time $`_4`$, instead of $`S_1\times _3`$. This path integral can be defined even for four manifolds with non-vanishing Euler number, that is ones for which a Minkowski metric is impossible. There is no reason for this new quantity to be real but it will turn out to have an interesting response to the electromagnetic duality transformations (2.6) (as pointed out by E Witten and E Verlinde). So these extended partition functions do have interesting mathematical properties as we shall see in more detail and it would be interesting to understand what, if any, physical significance they have. We shall henceforth refer to the real partition functions associated with $`S_1\times _3`$ as “strict partition functions”.
We shall see that both the Euler number, $`\chi (_4)`$, and the Hirzebruch signature, $`\eta (_4)`$, vanish for manifolds $`S_1\times _3`$. These are the two topological invariants of a $`_{4k}`$ manifold that are “local” in the sense that they can be expressed as integrals of closed forms over the manifold. Linear combinations of these topological invariants, namely $`(\chi \pm \eta )/2`$, will specify in a precise way how the extended partition functions deviate from satisfying exact electromagnetic duality.
These conclusions will depend on the explicit evaluation of the functional integral expression (2.10) for the partition function on any $`_4`$, and this is facilitated by taking account of another aspect of the quantum theory. If electrically charged particles are to be treated quantum mechanically, the background field strengths must satisfy certain Dirac flux quantisation conditions in order to allow the possible presence of complex wave functions for them. It is this that imposes a discrete structure that converts (2.10) to a sum rather than an integral, at least in the semiclassical approximation, which is very likely exact.
As we shall see, the resultant expression for the partition function (2.9) is proportional to an infinite sum forming a generalised sort of theta function associated with the lattice of homology classes of two-cycles in the space-time four-manifold $`_4`$. As explained below, this lattice is what is known as the free part of $`H_2(_4,𝐙)`$ and is intimately connected to the Dirac quantised fluxes.
If $`_4`$ is orientable, smooth, closed and compact, its topological structure satisfies the symmetry known as Poincaré duality. This implies that the following relation between the five Betti numbers of $`_4`$:
$$b_0=b_4,b_1=b_3.$$
$`(2.12)`$
Hence the Euler number is given by
$$\chi (_4)=2(b_0b_1)+b_2.$$
$`(2.13)`$
Furthermore, the aforementioned lattice, which has dimension $`b_2`$, is unimodular with respect to the scalar product furnished by the intersection number. It is this which is the origin of the covariance of (2.9) with respect to the $`S`$-transformation of electromagnetic duality:
$$S=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)$$
$`(2.14)`$
sending $`\tau `$ to $`1/\tau `$, a especially interesting example of (2.6). Before explaining this we must review the Dirac quantisation condition for fluxes in more detail, paying particular attention to the extra subtleties associated with spinning particles carrying electric charge.
## 3. The quantisation condition on four-manifolds
As already explained, we consider Maxwell theory in space-times that are compact, connected and oriented four-dimensional manifolds $`_4`$, whether or not their Euler number vanishes. These spaces are particularly convenient because they satisfy Poincaré duality, which is a topological property closely related to electromagnetic duality. The main mathematical tool needed in understanding the implications of the global topology of $`_4`$ on the duality properties of Maxwell theory is homology and cohomology theory. This discipline is described in many textbooks (see for example \[Schwarz 1994\]), and a short introduction to the relevant ideas was given in our earlier paper \[M Alvarez and Olive 1999\]. We shall follow the notations of the latter without full explanation.
The physical relevance of homology and cohomology theory is that it provides the natural mathematical language for the ideas of Faraday, Maxwell and, later, Dirac concerning electromagnetic theory. Important physical quantities are the (magnetic) fluxes of the field strength $`F`$ through a complete set of two-cycles $`\mathrm{\Sigma }_1,\mathrm{}\mathrm{\Sigma }_{b_2}`$ within the space-time four manifold $`_4`$.
According to Poincaré’s lemma, the gauge potential $`A`$, satisfying $`F=dA`$, can be constructed by integration locally, in topologically trivial neighbourhoods of $`_4`$, up to a gauge transformation. $`A`$ is needed to define the electromagnetic coupling to complex wave-functions of electrically charged particles. If the wave function is scalar, corresponding to a boson with charge $`q_B`$, quantum mechanical consistency of the patching procedure for it requires the fluxes to be quantised \[Dirac 1931, Wu and Yang 1975, O Alvarez 1985\]:
$$\frac{q_B}{2\pi \mathrm{}}_\mathrm{\Sigma }F=m(\mathrm{\Sigma })𝐙$$
$`(3.1)`$
By virtue of ordinary Stokes’ theorem and that the facts that $`\mathrm{\Sigma }`$ and $`F`$ are both closed, the value of the flux is unchanged either if $`\mathrm{\Sigma }`$ is replaced by $`\mathrm{\Sigma }+\mathrm{\Pi }`$, with $`\mathrm{\Pi }`$ a three-dimensional chain, or if $`F`$ is replaced by $`F+dB`$. Also, integer linear combination of cycles that satisfy the quantisation condition $`\mathrm{\Sigma }`$ also satisfy the same condition.
These statements can be summarised by saying that the cycles $`\mathrm{\Sigma }`$ are free elements of the integer homology class $`H_2(_4,𝐙)`$, and $`F`$ is in a cohomology class $`H^2(_4,𝐙)`$. Now $`H_2(_4,𝐙)`$ is an abelian group (with respect to the natural addition operation) and it possesses a unique subgroup built of elements of finite order, called the torsion group $`T_2(_4,𝐙)`$. The quotient group
$$F_2(_4,𝐙)H_2(_4,𝐙)/T_2(_4,𝐙)$$
$`(3.2)`$
is “free” and consists of $`b_2`$ copies of the integers, where $`b_2`$ is the second Betti number. This is the same as saying that $`F_2(_4,𝐙)`$ is a lattice of dimension $`b_2`$.
A slightly more general version of (3.1) is the “quantum Stokes’ relation”:
$$\mathrm{e}^{\mathrm{i}\frac{q_B}{\mathrm{}}_\mathrm{\Sigma }F}=\mathrm{e}^{\mathrm{i}\frac{q_B}{\mathrm{}}_\mathrm{\Sigma }A}.$$
$`(3.3)`$
Now $`\mathrm{\Sigma }`$ is allowed to have a non-vanishing boundary $`\mathrm{\Sigma }`$ and hence be a two-chain rather than a two-cycle. In the limiting case when the boundary vanishes, (3.1) is recovered (and so is necessary for the validity of (3.3)). The quantity on the right hand side of (3.3) is Dirac’s path dependent phase factor \[Dirac 1955\]. It is well defined in the situation described even though the exponent is not. Such phase factors are relevant in several different contexts \[Bohm-Aharonov, Wilson\] and go by several other names (Wilson loop, $`U(1)`$ holonomy etc).
However, many electrically charged particles, such as the electron, also carry spin and, as a result (3.1) and (3.3) may have to be modified if the topology of space-time is sufficiently complicated. If the complex wave function to which $`A`$ couples is spinor rather than scalar, and the associated fermionic particle carries charge $`q_F`$, (3.3) is modified by the presence of a possible minus sign \[M Alvarez and Olive 1999\].
$$\mathrm{e}^{\mathrm{i}\frac{q_F}{\mathrm{}}_\mathrm{\Sigma }F}=(1)^{w(\mathrm{\Sigma })}\mathrm{e}^{\mathrm{i}\frac{q_F}{\mathrm{}}_\mathrm{\Sigma }A},$$
$`(3.4)`$
at least for two-chains $`\mathrm{\Sigma }`$ whose boundary is an even cycle:
$$\mathrm{\Sigma }=2\alpha .$$
$`(3.5)`$
There are essentially two possibilities for this when $`\mathrm{\Sigma }`$ is odd. Either $`\alpha `$ vanishes and $`\mathrm{\Sigma }`$ is a closed surface, or not. If the latter, $`\mathrm{\Sigma }`$ could be the real projective plane in two dimensions, and so not orientable. Notice that although $`2\alpha `$ is closed, $`\alpha `$ itself is not. Hence the one-cycle $`\alpha `$ is what is known as a torsion cycle.
The sign factor $`(1)^{w(\mathrm{\Sigma })}`$ in (3.4) arises unambiguously in the procedure of patching together the neighbourhoods that make up $`\mathrm{\Sigma }`$, precisely when $`\mathrm{\Sigma }`$ satisfies (3.5), \[M Alvarez and Olive 1999\]. However, when $`\mathrm{\Sigma }`$ is closed, $`w(\mathrm{\Sigma })`$ can be constructed independently as the integer specifying the self-intersection number of $`\mathrm{\Sigma }`$ with itself. This is possible because $`\mathrm{\Sigma }`$ is a two-cycle in a closed oriented four-manifold. The equivalence of the two notions (mod $`2`$) can be deduced from the Atiyah-Singer index theorem on $`_4`$. When $`\mathrm{\Sigma }`$ is neither closed nor even, yet satisfies (3.5), it is not oriented and its self intersection number can only be defined mod $`2`$, and not absolutely. This still matches the previous definition, according to Wu’s formula. The important point is that the sign factor depends only on the topology of the background space-time.
Let us temporarily set the charges $`q_B`$ and $`q_F`$ equal. When the sign factor $`(1)^{w(\mathrm{\Sigma })}`$ equals $`1`$ equations (3.3) and (3.4) agree but when it equals $`1`$ they appear to differ. However the discrepancy is illusory as the gauge potentials in each equation differ as they are constructed by gauge inequivalent patching procedures.
On the other hand, when we consider the limiting case in which $`q_B`$ and $`q_F`$ both vanish the two versions of “quantum Stokes’”, (3.3) and (3.4), reduce to $`1=1`$ and $`1=(1)^{w(\mathrm{\Sigma })}`$, respectively. Now the second equation is manifestly a contradiction if the sign factor is negative. What this means is that a false assumption has been adopted, namely that it is possible to place an electrically neutral spinor wave function on $`_4`$. Clearly this is forbidden if there is any $`\mathrm{\Sigma }`$ satisfying (3.5) for which the sign factor $`(1)^{w(\mathrm{\Sigma })}`$ is negative. Mathematicians are familiar with this phenomenon and say that such an $`_4`$ lacks a “spin structure”. They recognise $`w(\mathrm{\Sigma })`$, as the Stiefel-Whitney class. It is an element of $`H^2(_4,𝐙_2)`$ and its nontriviality provides an obstruction to spin structures. See Appendix A of \[Lawson and Michelsohn\].
In our previous paper we found it convenient to separate all the four-manifolds under consideration into three types: I, II and III.
A four-manifold is of Type I if the sign factor $`(1)^{w(\mathrm{\Sigma })}`$ is plus one in all cases (3.5). The the flux quantisation condition reads
$$\frac{q_F}{2\pi \mathrm{}}_\mathrm{\Sigma }F=m(\mathrm{\Sigma })𝐙.$$
$`(3.6)`$
This is the same as (3.1) with $`q_F`$ replacing $`q_B`$. It follows that all cycles have even self-intersection number.
The intersection numbers of pairs of two-cycles endow the lattice $`F_2(_4,𝐙)`$ of free cycles with an integral scalar product. The Poincaré duality that applies to the four-manifolds under consideration implies that this scalar product is unimodular. Unimodular lattices fall naturally into two classes even, or odd. It follows from our remarks that the unimodular lattice $`F_2(_4,𝐙)`$ is even for Type I manifolds.
A four-manifold is of Type II if the sign factor $`(1)^{w(\mathrm{\Sigma })}`$ is minus one for at least one two-cycle $`\mathrm{\Sigma }`$. Such a cycle has an odd self-intersection number and consequently the unimodular lattice $`F_2(_4,𝐙)`$ is odd.
The flux quantisation implied by (3.4) reads
$$\frac{q_F}{2\pi \mathrm{}}_\mathrm{\Sigma }F=m(\mathrm{\Sigma })+\frac{w(\mathrm{\Sigma })}{2};m(\mathrm{\Sigma })𝐙.$$
$`(3.7)`$
Thus, when $`w(\mathrm{\Sigma })`$ is odd, the flux is fractional rather than integral and, in particular, can never vanish, unlike the integral fluxes.
The remaining possibility, Type III, arises when the sign factor $`(1)^{w(\mathrm{\Sigma })}`$ equals plus one for all cycles but minus for at least one of the open two-chains satisfying (3.5). As the self-intersection numbers of all cycles are even so is the unimodular lattice $`F_2(_4,𝐙)`$.
Thus, in summary, type I manifolds are the only ones that support spin structures, that is electrically neutral spinor wave functions. All three types support “spin<sub>C</sub>” structures, that is electrically charged complex spinor wave functions, provided the backgound fluxes satisfy the appropriate quantisation conditions.
These conditions imply that all fluxes are integral if the unimodular lattice $`F_2(_4,𝐙)`$ is even, that is for types I and III, but that some fluxes at least must be fractional when the lattice is odd, that is for Type II manifolds. Standard examples of the three types of four-manifold are, for type I the torus, $`T_4`$, the sphere $`S_4`$ and $`K(3)`$, for type II, the complex projective space, $`CP(2)`$ and for type III, $`S_2\times S_2/Z_2`$.
## 4. Integral lattices
Unimodular lattices are a sort of integral lattice with special structural features that will play a role in the picture of flux quantisation that we have begun to explain. Here we pause to explain some of the relevant concepts.
A lattice $`\mathrm{\Lambda }`$ of dimension $`n`$ is a discrete subgroup of $`𝐑^n`$ defined as
$$\mathrm{\Lambda }=\left\{\underset{i=1}{\overset{n}{}}n_ie_i,n_i𝐙\right\},$$
$`(4.1)`$
where $`e_1,e_2,\mathrm{},e_n`$ are elements of $`\mathrm{\Lambda }`$ spanning $`𝐑^n`$ and so providing a basis for $`\mathrm{\Lambda }`$.
The vector space $`𝐑^n`$ may be endowed with a real, symmetric scalar product, denoted $`xy`$, which is nonsingular but not necessarily positive definite. If this scalar product has $`b^\pm `$ positive (negative) eigenvalues the signature $`\eta `$ of $`\mathrm{\Lambda }`$ is
$$\eta (\mathrm{\Lambda })=b^+b^{},\text{where}b^++b^{}=n.$$
$`(4.2)`$
For example, when $`F_2(_4,𝐙)`$ is endowed with the scalar product given by the intersection numbers, its signature is known as the Hirzebruch signature of $`_4`$, and denoted $`\eta (_4)`$.
Given a lattice and a scalar product we can define another lattice, known as the reciprocal lattice:
$$\mathrm{\Lambda }^{}=\{x𝐑^n:yx𝐙y\mathrm{\Lambda }\}.$$
$`(4.3)`$
Obviously $`\mathrm{\Lambda }^{}=\mathrm{\Lambda }`$. $`\mathrm{\Lambda }`$ is said to be integral if
$$\mathrm{\Lambda }\mathrm{\Lambda }^{}$$
$`(4.4)`$
as a subgroup. Then the quotient group is well-defined and abelian
$$Z(\mathrm{\Lambda })=\mathrm{\Lambda }^{}/\mathrm{\Lambda }.$$
$`(4.5)`$
$`\mathrm{\Lambda }`$ is unimodular if its order, $`|Z(\mathrm{\Lambda })|`$, is one. This is the same as saying that the lattices $`\mathrm{\Lambda }`$ and $`\mathrm{\Lambda }^{}`$ coincide.
A convenient choice of fundamental domain in $`\mathrm{\Lambda }`$ consists of $`X=_ix_ie_i`$ with $`0x_i1`$. The volume of this fundamental domain is
$$V(\mathrm{\Lambda })=\sqrt{det(e_ie_j)}=|det(e_1,e_2,\mathrm{},e_n)|.$$
A standard basis for the reciprocal lattice is the reciprocal basis $`f_1,f_2\mathrm{},f_n`$ satisfying $`e_if_j=\delta _{ij}`$. Then
$$V(\mathrm{\Lambda }^{})V(\mathrm{\Lambda })=|det(f_1,f_2,\mathrm{},f_n)det(e_1,e_2,\mathrm{},e_n)|=|det(e_if_j)|=det(\delta _{ij})=1.$$
But $`|Z(\mathrm{\Lambda })|`$ copies of the fundamental domain of $`\mathrm{\Lambda }^{}`$ make up a fundamental domain for $`\mathrm{\Lambda }`$. Hence
$$V(\mathrm{\Lambda })=(V(\mathrm{\Lambda }^{}))^1=\sqrt{|Z(\mathrm{\Lambda })|}.$$
$`(4.6)`$
Thus $`\mathrm{\Lambda }`$ is unimodular if it is integral and the volume of its fundamental domain equals unity. This is what was used above.
In fact any choice of fundamental domain has the same volume. So changes of basis form the infinite discrete group $`SL(n,𝐙)`$.
If $`x`$ is an element of an integral lattice $`xx`$ is automatically an integer. If it is always an even integer $`\mathrm{\Lambda }`$ is said to be even. Otherwise $`\mathrm{\Lambda }`$ is odd.
Of course this applies to unimodular lattices but these possess an extra feature, the existence of elements called characteristic elements. $`c\mathrm{\Lambda }`$ is a characteristic vector if
$$cx+xx2𝐙x\mathrm{\Lambda }.$$
$`(4.7)`$
It is easy to establish that such quantities always exist and that there is an ambiguity of precisely mod $`2\mathrm{\Lambda }`$. It follows that $`cc`$ is uniquely defined, mod $`8`$ for any unimodular lattice. In fact it is known that
$$cc=\eta (\mathrm{\Lambda })+8𝐙$$
$`(4.8)`$
This will follow from our discussion of theta functions but a direct proof can be found in \[Milnor and Husemoller\]. The zero element is always a characteristic vector for an even unimodular lattice. Hence (4.8) implies that for such lattices the signature is a multiple of eight. A famous example is the $`E_8`$ root lattice.
The fact that an odd unimodular lattice $`\mathrm{\Lambda }`$ possesses a characteristic vector $`c`$ will facilitate some constructions which will be relevant to the analysis of type II four manifolds. First:
$$\mathrm{\Lambda }_{TOTAL}=\mathrm{\Lambda }(\mathrm{\Lambda }+c/2)$$
$`(4.9)`$
defines a lattice. Obviously $`\mathrm{\Lambda }_{TOTAL}/\mathrm{\Lambda }=Z_2`$, so that, by a slight extension of (4.6), $`V(\mathrm{\Lambda }_{TOTAL})=1/2`$. Furthermore we can split
$$\mathrm{\Lambda }=\mathrm{\Lambda }_{EVEN}\mathrm{\Lambda }_{ODD}$$
$`(4.10)`$
where $`\mathrm{\Lambda }_{EVEN/ODD}`$ consists of those elements of $`\mathrm{\Lambda }`$ with even/odd squared length. Now $`\mathrm{\Lambda }/\mathrm{\Lambda }_{EVEN}=Z_2`$ and so $`V(\mathrm{\Lambda }_{EVEN})=2`$. Thus $`c`$ itself will be in $`V(\mathrm{\Lambda }_{EVEN/ODD})`$ according as $`c^2`$ is even or odd, or according to (4.8), as the signature (4.2) is even or odd. It is easy to see that $`\mathrm{\Lambda }_{EVEN}`$ and $`\mathrm{\Lambda }_{TOTAL}`$ are a pair of reciprocal lattices. Furthermore we have
$$Z(\mathrm{\Lambda }_{EVEN})=\mathrm{\Lambda }_{TOTAL}/\mathrm{\Lambda }_{EVEN}=\{\begin{array}{cc}Z_4\hfill & \text{if }c\mathrm{\Lambda }_{ODD}\text{;}\hfill \\ Z_2\times Z_2\hfill & \text{if }c\mathrm{\Lambda }_{EVEN}\text{.}\hfill \end{array}$$
$`(4.11)`$
This follows because the quotient group is determined by the addition rules for the relevant cosets in the decomposition
$$\mathrm{\Lambda }_{TOTAL}=\mathrm{\Lambda }_{EVEN}(\mathrm{\Lambda }_{EVEN}+c/2)\mathrm{\Lambda }_{ODD}(\mathrm{\Lambda }_{ODD}+c/2)$$
$`(4.12)`$
It follows from (4.11) that, when $`c\mathrm{\Lambda }_{EVEN}`$,
$$\mathrm{\Lambda }^{}\mathrm{\Lambda }_{EVEN}(\mathrm{\Lambda }_{ODD}+c/2),\mathrm{\Lambda }^{\prime \prime }\mathrm{\Lambda }_{EVEN}(\mathrm{\Lambda }_{EVEN}+c/2)$$
$`(4.13)`$
are both lattices and that $`V(\mathrm{\Lambda }^{})=V(\mathrm{\Lambda }^{\prime \prime })=V(\mathrm{\Lambda }_{EVEN})/2=1`$. Hence, when $`\mathrm{\Lambda }^{}`$ and $`\mathrm{\Lambda }^{\prime \prime }`$ are both integral i.e. when $`c^2`$ is a multiple of four, they are unimodular, being even or odd according as $`c^2/4`$ is. This is a powerful result as odd unimodular lattices can be found for all signatures, simply by considering products of integers (hypercubic lattices). In fact this is the simplest construction of an even unimodular lattice and it works precisely when the signature $`\eta `$ is a multiple of eight since (4.8) can be checked explicitly for the hypercubic lattices.
The lattice of free two-cycles $`F_2(_4,𝐙)`$, endowed with the scalar product furnished by the intersection number of pairs of two-cycles, was unimodular. It was odd when $`_4`$ was Type II and even otherwise. In the former case it must have a characteristic vector (4.7). This turns out to be related to the quantity $`w(\mathrm{\Sigma })`$ introduced via the quantum Stokes relation (3.3). Recall that
$$w(\mathrm{\Sigma })=I(\mathrm{\Sigma },\mathrm{\Sigma })\text{ mod }2\mathrm{\Sigma }F_2(_4,𝐙).$$
$`(4.14)`$
Now expand $`\mathrm{\Sigma }`$ in terms of a basis $`\mathrm{\Sigma }=_{j=1}^{b_2}n^j\mathrm{\Sigma }_j`$ and introduce the notations
$$I(\mathrm{\Sigma }_i,\mathrm{\Sigma }_j)=(Q^1)_{ij},w(\mathrm{\Sigma }_i)=w_i,$$
$`(4.15)`$
where the matrix $`Q`$ and its inverse both have integer entries as it has determinant $`\pm 1`$. Now (4.14) reads
$$\underset{i,j,k=1}{\overset{b_2}{}}n^i(Q^1)_{ij}Q^{jk}w_k=\underset{i,j=1}{\overset{b_2}{}}n^i(Q^1)_{ij}n^j.$$
This means that the characteristic vector of the lattice $`F_2(_4,𝐙)`$ can be represented by the two-cycle
$$\gamma =\underset{j,k=1}{\overset{b_2}{}}w_jQ^{jk}\mathrm{\Sigma }_k.$$
$`(4.16)`$
Actually this argument is the reverse of that used in our previous paper.
What will be more important is the lattice structure of the quantised fluxes through these two-cycles. A scalar product is provided by the theta term contribution to the action (2.2). This is due to a generalisation of the Riemann bilinear identity applicable to any pair of closed forms. Here it reads
$$__4FF^{}=\underset{i,j=1}{\overset{b_2}{}}_{\mathrm{\Sigma }_i}FQ^{ij}_{\mathrm{\Sigma }_j}F^{}.$$
$`(4.17)`$
Considering first the quantised flux background necessary for $`_4`$ to support a complex scalar wave function, it is natural to define field strengths $`F^1,F^2,\mathrm{},F^{b_2}`$ referring to a basis reciprocal to $`\mathrm{\Sigma }_i`$:
$$\frac{q_B}{2\pi \mathrm{}}_{\mathrm{\Sigma }_i}F^j=\delta _i^j.$$
$`(4.18)`$
The general solution to the quantisation condition (3.1) is
$$F=\underset{i=1}{\overset{b_2}{}}m_iF^i,m_i𝐙$$
$`(4.19)`$
(or something cohomologous). These de Rham cohomology classes form a lattice with unimodular scalar product given by
$$\left(\frac{q_B}{2\pi \mathrm{}}\right)^2__4FF=\underset{i,j=1}{\overset{b_2}{}}m_iQ^{ij}m_j.$$
$`(4.20)`$
Thus the fluxes too form a unimodular lattice which is odd if $`_4`$ is of Type II and even otherwise.
Now consider the quantised flux backgrounds necessary for $`_4`$ to support a complex spinor wave function and define a basis similar to (4.18) but with $`q_F`$ replacing $`q_B`$. Then the general solution to the quantisation condition (3.6) is
$$F=\underset{i=1}{\overset{b_2}{}}\left(m_i+\frac{w_i}{2}\right)F^i,$$
$`(4.21)`$
(or something cohomologous), with scalar product
$$\left(\frac{q_F}{2\pi \mathrm{}}\right)^2__4FF=\underset{i,j=1}{\overset{b_2}{}}\left(m_i+\frac{w_i}{2}\right)Q^{ij}\left(m_j+\frac{w_j}{2}\right).$$
$`(4.22)`$
It follows from (4.16) that $`_iw_iF^i`$ represents the characteristic vector for the unimodular flux lattice. Thus when $`_4`$ is of Type I or III so that all $`w_i`$ vanish (mod $`2`$), the fluxes lie on an even unimodular lattice. But when $`_4`$ is of Type II the fluxes lie on an odd unimodular lattice displaced by half its characteristic vector, that is the non trivial coset of (4.9).
In particular, if $`q_B`$ and $`q_F`$ are equal, the choice between fluxes corresponding to the two terms in the decomposition (4.9) of $`\mathrm{\Lambda }_{TOT}`$ depends on whether the complex wave function is scalar or spinor. Notice that according to (4.7) one half of the expression (4.22) differs from $`cc/8`$ by an integer. Then (4.8) would follow from the integrality of the index of the Dirac operator on $`_4`$ using the version of the Atiyah-Singer index theorem quoted in \[M Alvarez and Olive\].
## 5. The Maxwell partition function and theta functions
We are now in a position to return to the evaluation of the partition function (2.9), (2.10) in terms of the lattice structures just described.
The basic idea of the semi-classical approximation to (2.10) is to expand the integrand about the stationary points of the exponent which is given in terms of the classical action (2.2), that is, solutions to the Euler-Lagrange equations, here the Maxwell equation $`dF=0`$. This together with $`dF=0`$ means that classical solutions $`F`$ are harmonic two-forms on $`_4`$. Because the relevant metric on $`_4`$ is Euclidean, Hodge’s theorem is valid and states that there really is one and only one harmonic two-form, i.e. stationary point, in each cohomology class. As we saw these classes are labelled by the $`b_2`$ magnetic integers $`m_1,m_2,\mathrm{},m_{b_2}`$. It is convenient now to suppose that the representative field strengths $`F^1,F^2,\mathrm{},F^{b_2}`$ satisfying (4.18) are these harmonic ones. Then the classical solutions in each class are given precisely by (4.19) or (4.21), as the case may be. The space of harmonic two-forms on $`_4`$ divides into a direct sum of two subspaces consisting of self-dual and anti-self-dual harmonic two-forms. The dimensions of these subspaces are the same numbers $`b^+`$ and $`b^{}`$ previously defined in purely topological terms by means of the $`F_2(_4,𝐙)`$ intersection matrix $`Q^1`$.
The result is a sum over stationary points, that is the points of the lattices described in the previous section. The contribution of each of these stationary points consists of the exponential of $`\frac{\mathrm{i}}{\mathrm{}}`$ times the Euclidean action $`W_{EUCLIDEAN}`$ evaluated at the stationary point, all multiplied by a determinantal factor $`\mathrm{\Delta }(\tau )`$ formed by the Gaussian integral of the quadratic fluctuations about it as well as zero modes. According to the arguments of E Verlinde and Witten this determinantal factor is common to all the terms of the sum. Furthermore, Witten argued that it takes the form
$$\mathrm{\Delta }(\tau )=Z_0(\tau _2)^{\frac{b_11}{2}}$$
$`(5.1)`$
where $`Z_0`$ is independent of $`\tau `$ and $`b_1`$ denotes the first Betti number. $`\tau _1`$ and $`\tau _2`$ are as defined in (2.3) but with $`q`$ replaced by $`q_B`$ or $`q_F`$, as appropriate. Hence the partition function is given by this factor $`\mathrm{\Delta }`$ times a sum over a lattice of an exponential of an expression quadratic in the coordinates of the lattice points. This is recognisable as a sort of theta function associated with the flux lattice which we shall examine in more detail. Because the action (2.2) is quadratic in field strengths the result of the procedure is exact and this is what makes this version of Maxwell theory mathematically tractable.
First we evaluate $`\mathrm{exp}\left(\frac{\mathrm{i}}{\mathrm{}}W_{EUCLIDEAN}\right)`$ at the stationary points (4.19) relevant to complex scalar wave functions on any four manifold and to complex spinor wave functions on four-manifolds of Types I and III. The contribution of the theta term to the exponent is evaluated using (4.20):
$$\frac{\mathrm{i}\tau _1}{2\mathrm{}\tau _2}__4FF=\mathrm{i}\pi \tau _1m^TQm.$$
$`(5.2)`$
The evaluation of the Maxwell term relies on the fact that, if $`F^i`$ is harmonic, so is its dual $`F^i`$. Hence there exists a matrix $`G`$ whereby
$$F^i=G^{ij}(Q^1)_{jk}F^k.$$
$`(5.3)`$
Because of (2.1) and the Euclidean nature of the metric
$$\left(GQ^1\right)^2=1.$$
$`(5.4)`$
Hence
$$\frac{1}{2\mathrm{}}__4FF=\pi \tau _2m^TGm.$$
$`(5.5)`$
In the course of this argument it becomes clear that the $`b_2\times b_2`$ matrix $`G`$ is symmetric and positive definite, reflecting the nature of the metric upon which it depends. Finally, evaluated at the classical solution (4.19),
$$\mathrm{e}^{\frac{i}{\mathrm{}}W_{EUCLIDEAN}}=\mathrm{e}^{\mathrm{i}\pi m^T(\mathrm{\Omega }(\widehat{\tau }))m}.$$
$`(5.6)`$
where
$$\mathrm{\Omega }(\tau )=\tau _1Q+i\tau _2G,$$
$`(5.7)`$
Hence, for the backgrounds considered so far, the partition function
$$Z(\tau )=\mathrm{\Delta }(\tau )\mathrm{\Theta }(\mathrm{\Omega }(\tau )),$$
$`(5.8)`$
where the richest structure resides in the theta function factor which has the general form
$$\mathrm{\Theta }(\mathrm{\Omega })=\underset{m_i𝐙}{}\mathrm{e}^{\mathrm{i}\pi m^T\mathrm{\Omega }m}.$$
$`(5.9)`$
The sum converges if the complex symmetric matrix $`\mathrm{\Omega }`$ has imaginary part which is positive definite. This is certainly so in our case as both $`\tau _2`$ and $`G`$ are positive in view of their physical interpretations.
As a consequence of the Poisson summation formula, the theta function (5.9) obeys the property
$$\mathrm{\Theta }(\mathrm{\Omega }^1)=\sqrt{det(\mathrm{i}\mathrm{\Omega })}\mathrm{\Theta }(\mathrm{\Omega }),$$
$`(5.10a)`$
where the positive sign of the root is understood. Furthermore
$$\mathrm{\Theta }(\mathrm{\Omega })=\mathrm{\Theta }(A\mathrm{\Omega }A^T)=\mathrm{\Theta }(\mathrm{\Omega }+B)$$
$`(5.10b)`$
where $`AGL(b_2,𝐙)`$ and $`B`$ is a symmetric matrix with integer entries which are even on the diagonal. These symmetries generate a subgroup of $`Sp(2b_2,𝐙)`$ with finite index. We do not know whether this has any physical significance but it does contain a subgroup recognisable as consisting of discrete electromagnetic duality transformations acting on $`\tau `$ when we take account of some special properties possessed by the matrix (5.7), namely
$$\mathrm{\Omega }(1/\tau )=Q\mathrm{\Omega }(\tau )^1Q,\mathrm{\Omega }(\tau +1)=\mathrm{\Omega }(\tau )+Q,$$
$`(5.11a)`$
and
$$\sqrt{det\left(\mathrm{i}\mathrm{\Omega }(\tau )\right)}=e^{\frac{2\pi \mathrm{i}\eta }{8}}\tau ^{b^+/2}(\tau ^{})^{b^{}/2},$$
$`(5.11b)`$
checked using (5.4) and the fact that $`GQ^1`$ has $`b^\pm `$ eigenvalues $`\pm 1`$. As before, $`\eta =b^+b^{}`$ is the signature of $`Q`$ and hence the Hirzebruch signature of $`_4`$.
Regarding the theta functions as functions of $`\tau `$, and denoting them accordingly as $`\mathrm{\Theta }(\tau )`$, it follows from (5.10) and (5.11) that
$$\mathrm{\Theta }(1/\tau )=e^{\frac{2\pi \mathrm{i}\eta }{8}}\tau ^{b^+/2}(\tau ^{})^{b^{}/2}\mathrm{\Theta }(\tau )$$
$`(5.12a)`$
$$\mathrm{\Theta }(\tau +1)=\mathrm{\Theta }(\tau )\text{if }Q\text{ is even},$$
$`(5.12b)`$
and
$$\mathrm{\Theta }(\tau +2)=\mathrm{\Theta }(\tau )\text{otherwise},$$
$`(5.12c)`$
We have already met the $`S`$-transformation (2.14). Together with
$$T=\left(\begin{array}{cc}1& 1\\ 0& 1\end{array}\right),$$
$`(5.13)`$
sending $`\tau \tau +1`$, it generates the discrete group $`Sp(2,𝐙)`$, a subgroup of the $`Sp(2,𝐑)`$ duality group of section 2. On the other hand, $`S`$ and $`T^2`$ generate a subgroup of $`Sp(2,𝐙)`$ called the Hecke group, $`\mathrm{\Gamma }_\theta `$, one of the three distinct subgroups of index three.
So, roughly speaking, the partition function associated with backgrounds of integrally quantised fluxes does manifest a property of electromagnetic duality in that it transforms simply under certain discrete subgroups of $`Sp(2,𝐑)`$ acting on the dimensionless variable $`\tau `$. This discrete subgroup is $`Sp(2,𝐙)`$ or $`\mathrm{\Gamma }_\theta `$ depending on whether the lattice of integrally quantised fluxes is even or odd. But it is precisely in the latter case that real spinor wave functions are forbidden.
The existence of complex spinor wave functions requires fractionally quantised fluxes, as explained above, and hence a different partition function. Repetition of the above calculation using (4.21) rather than (4.19) yields a similar result but with a different theta function, namely
$$\mathrm{\Theta }(\tau )_{w/2}=\underset{m_i𝐙}{}\mathrm{e}^{\mathrm{i}\pi (m+\frac{1}{2}w)^T\mathrm{\Omega }(\tau )(m+\frac{1}{2}w)},$$
$`(5.14)`$
We shall show that this transforms nicely under the action of another subgroup of $`Sp(2,𝐙)`$, also of index three, like $`\mathrm{\Gamma }_\theta `$. Furthermore, if $`q_F=q_B`$, (5.14) is related to the previous theta function, (5.9), by an element $`ST`$ of $`Sp(2,𝐙)`$, outside the two subgroups. Thus there is a sense in which the action of the full modular group $`PSL(2,𝐙)Sp(2,𝐙)/Z_2`$ of electromagnetic duality transformations can be realised taking into account the difference in the background of quantised fluxes needed to support complex scalar and spinor wave functions respectively on four-manifolds of type II.
In order to explain this in more detail we need more developments in formalism.
Notice that although $`\mathrm{\Theta }(\mathrm{\Omega })`$ is holomorphic in $`\mathrm{\Omega }`$, $`\mathrm{\Theta }(\tau )`$ is not holomorphic in $`\tau `$ unless $`b^{}`$ vanishes. For then $`G=Q`$ and $`\mathrm{\Omega }(\tau )=\tau Q`$. Similarly, if $`b^+`$ vanishes, $`\mathrm{\Omega }(\tau )=\tau ^{}Q`$.
## 6. Theta functions and integral lattices
In this section the theta function construction above is extended in a way that is similar to that described in the book \[Green, Schwarz and Witten 1987\] but which goes further. It seems to be novel and intrinsically interesting. A theta function is associated to each element of the group $`Z(\mathrm{\Lambda })`$, (4.5), defined by an integral lattice $`\mathrm{\Lambda }`$, whatever the signature of its scalar product. These $`|Z(\mathrm{\Lambda })|`$ theta functions support an action of the group $`Sp(2,𝐙)`$ (i.e. $`SL(2,𝐙)`$) if $`\mathrm{\Lambda }`$ is even, and its Hecke subgroup $`\mathrm{\Gamma }_\theta `$ if it is odd (or, more properly, their metaplectic extensions).
Careful analysis of how the effect of the generators $`S`$ and $`T`$ of $`Sp(2,𝐙)`$ satisfy the relation $`(ST)^3=I`$ will yield “Milgram’s formula” expressing the signature (mod $`8`$) in terms of the structure of $`Z(\mathrm{\Lambda })`$, when $`\mathrm{\Lambda }`$ is even. The construction (4.11) means that odd unimodular lattices can also be dealt with, leading to a proof of (4.8). The relevance of all this to electromagnetic duality will be explained in the following section.
We start with Poisson’s summation formula in the form:
$$\underset{n_i𝐙}{}f(x_i+n_i)=\underset{m_i𝐙}{}\mathrm{e}^{2\pi \mathrm{i}m_jx_j}\stackrel{~}{f}(m_i),\text{where}\stackrel{~}{f}(k)=d^nx\mathrm{e}^{2\pi \mathrm{i}k_jx_j}f(x)$$
$`(6.1)`$
denotes the Fourier transform and the summation convention is understood. The sums and integrals converge if
$$f(x)=\mathrm{e}^{\pi \mathrm{i}x_j\mathrm{\Omega }_{jm}x_m}\text{so}\stackrel{~}{f}(k)=\frac{1}{\sqrt{det(\mathrm{i}\mathrm{\Omega })}}\mathrm{e}^{\pi \mathrm{i}k_j(\mathrm{\Omega }^1)_{jm}k_m}$$
and $`\mathrm{\Omega }`$ is a complex symmetric matrix with positive definite imaginary part, as before. Hence
$$\underset{n_i𝐙}{}e^{\pi \mathrm{i}(x_j+n_j)\mathrm{\Omega }_{jm}(x_m+n_m)}=\frac{1}{\sqrt{det(\mathrm{i}\mathrm{\Omega })}}\underset{m_i𝐙}{}\mathrm{e}^{2\pi \mathrm{i}m_jx_j}\mathrm{e}^{\pi \mathrm{i}m_j\mathrm{\Omega }_{jk}^1m_k}.$$
Now suppose that, as in (4.1), $`n_j`$ are the coordinates of the point $`l`$ of the lattice $`\mathrm{\Lambda }`$ with respect to the basis $`e_i`$, and that $`m_j`$ are the coordinates of the point $`l^{}`$ of the reciprocal lattice $`\mathrm{\Lambda }^{}`$ with basis $`f_j`$. Thus
$$l=\underset{j=1}{\overset{n}{}}n_je_j,l^{}=\underset{j=1}{\overset{n}{}}m_jf_j\text{while}X=\underset{j=1}{\overset{n}{}}x_je_j.$$
Then, if we denote
$$\widehat{\mathrm{\Omega }}=f_j\mathrm{\Omega }_{jk}f_k^T,\text{then}(\widehat{\mathrm{\Omega }})^1=e_j(\mathrm{\Omega }^1)_{jk}e_k^T$$
and the Poisson summation formula reads
$$\underset{l\mathrm{\Lambda }}{}\mathrm{e}^{\pi \mathrm{i}(X+l)\widehat{\mathrm{\Omega }}(X+l)}=\frac{1}{\sqrt{det(\mathrm{i}\mathrm{\Omega })}}\underset{l^{}\mathrm{\Lambda }^{}}{}\mathrm{e}^{2\pi \mathrm{i}l^{}X}\mathrm{e}^{\pi \mathrm{i}l^T\widehat{\mathrm{\Omega }}^1l^{}}.$$
Now suppose that the lattice $`\mathrm{\Lambda }`$ is integral and that the $`Z(\mathrm{\Lambda })`$ coset decomposition of $`\mathrm{\Lambda }^{}`$ can be written
$$\mathrm{\Lambda }^{}=\mathrm{\Lambda }(\lambda _1+\mathrm{\Lambda })(\lambda _2+\mathrm{\Lambda })\mathrm{}(\lambda _{|Z(\mathrm{\Lambda })|1}+\mathrm{\Lambda }),$$
$`(6.2)`$
where $`\lambda _\beta `$ is a representative element of the $`\beta `$’th coset and $`\lambda _0`$ is understood to vanish. Choosing $`X`$ to be the $`\alpha `$’th of these representatives so $`\mathrm{e}^{2\pi \mathrm{i}l^{}X}=\mathrm{e}^{2\pi \mathrm{i}\lambda _\alpha \lambda _\beta }`$ if $`l^{}\lambda _\beta +\mathrm{\Lambda }`$, we can rearrange the Poisson summation formula so that sums over $`\mathrm{\Lambda }`$ occur on both sides:
$$\underset{l\lambda _\alpha +\mathrm{\Lambda }}{}\mathrm{e}^{\pi \mathrm{i}l\widehat{\mathrm{\Omega }}l}=\frac{1}{\sqrt{det(\mathrm{i}\mathrm{\Omega })}}\underset{\beta =0}{\overset{|Z(\mathrm{\Lambda })|1}{}}\mathrm{e}^{2\pi \mathrm{i}\lambda _\alpha \lambda _\beta }\underset{l\lambda _\beta +\mathrm{\Lambda }}{}\mathrm{e}^{\pi \mathrm{i}l\widehat{\mathrm{\Omega }}^1l}.$$
Now consider $`\mathrm{\Omega }`$ to have the special structure (5.7), where, as before, $`Q_{ij}=e_ie_j`$, but is now only integral and not necessarily unimodular. Then by the properties of the reciprocal basis
$$\widehat{\mathrm{\Omega }}(\tau )=\tau _1I+i\tau _2\widehat{G}\text{where}\widehat{G}=f_iG_{ij}f_j^T\text{and}\widehat{G}^2=I$$
by virtue of equations (5.4) and (2.1). It follows that $`\widehat{\mathrm{\Omega }}^1(\tau )=\widehat{\mathrm{\Omega }}(1/\tau )`$. Hence, defining the $`|Z(\mathrm{\Lambda })|`$ theta functions
$$\mathrm{\Theta }_\alpha (\tau )=\underset{l\lambda _\alpha +\mathrm{\Lambda }}{}\mathrm{e}^{\pi \mathrm{i}l(\tau _1+\mathrm{i}\tau _2\widehat{G})l},\alpha =0,1\mathrm{}|Z(\mathrm{\Lambda })|1,$$
$`(6.3)`$
the Poisson summation formula now reads
$$\mathrm{\Theta }_\alpha (\tau )=\tau ^{b^+/2}(\tau ^{})^{b^{}/2}\frac{\mathrm{e}^{2\pi \mathrm{i}\eta /8}}{\sqrt{|Z(\mathrm{\Lambda })|}}\underset{\beta =0}{\overset{|Z(\mathrm{\Lambda })|1}{}}\mathrm{e}^{2\pi \mathrm{i}\lambda _\alpha \lambda _\beta }\mathrm{\Theta }_\beta (1/\tau )$$
$`(6.4)`$
where the determinant was evaluated by (5.11b), modified to take account of the extra factor $`|detQ|=|Z(\mathrm{\Lambda })|`$, by (2.10), arising because $`Q`$ is no longer necessarily unimodular.
This is the action of the $`S`$-transformation (2.14). The response to $`T`$, (5.13), is simple when $`\mathrm{\Lambda }`$ is even. Otherwise $`T^2`$ must be considered
$$\mathrm{\Theta }_\alpha (\tau +1)=\mathrm{e}^{\pi \mathrm{i}\lambda _\alpha ^2}\mathrm{\Theta }_\alpha (\tau )\text{if }\mathrm{\Lambda }\text{ is even}$$
$`(6.5a)`$
$$\mathrm{\Theta }_\alpha (\tau +2)=\mathrm{e}^{2\pi \mathrm{i}\lambda _\alpha ^2}\mathrm{\Theta }_\alpha (\tau )\text{otherwise}$$
$`(6.5b)`$
As mentioned earlier, the two matrices $`S`$, (2.14), and $`T`$, (5.13), generate $`Sp(2,𝐙)SL(2,𝐙)`$, a discrete subgroup of the original duality group, $`Sp(2,𝐑)`$, while $`S`$ and $`T^2`$ generate the Hecke subgroup, $`\mathrm{\Gamma }_\theta `$. $`S`$ and $`T`$ are not independent since they satisfy the relations
$$S^2=I_2=(ST)^3.$$
$`(6.6)`$
Given the responses (6.4) and (6.5) of the theta functions, the relations (6.6) will have remarkable consequences that we now develop. First define the following action of the general element $`B=\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)`$ of $`Sp(2,𝐑)`$ on a function $`f(\tau )`$ of $`\tau _1`$ and $`\tau _2`$, (or equivalently $`\tau `$ and $`\tau ^{}`$),
$$\widehat{A}_{k_1k_2}(B)f(\tau )=(c\tau +d)^{k_1}(c\tau ^{}+d)^{k_2}f\left(\frac{a\tau +b}{c\tau +d}\right).$$
$`(6.7)`$
This is meaningful if $`k_1`$ and $`k_2`$ are are a pair of integers as will be supposed. It is also relevant for $`k_1`$ and $`k_2`$ to be half-integers. Treatment of this case requires the introduction of the metaplectic group to take account of the sign ambiguities that arise because of the square roots. In the interests of avoiding an over cumbersome notation this will not be done.
The virtue of the definition (6.7) is that this action satisfies the group property:
$$\widehat{A}_{k_1k_2}(B)\widehat{A}_{k_1k_2}(B^{})=\widehat{A}_{k_1k_2}(BB^{}).$$
$`(6.8)`$
Given a column vector of such functions, $`f_1(\tau ),f_2(\tau )\mathrm{}f_N(\tau )`$, we say, following \[V Kac 1990\], that they form a vector modular form if there exist matrices $`D_{\beta \alpha }(B)`$ such that
$$\widehat{A}_{k_1k_2}(B)f_\alpha (\tau )=\underset{\beta =1}{\overset{N}{}}f_\beta (\tau )D_{\beta \alpha }(B).$$
$`(6.9)`$
The integers $`k_1`$ and $`k_2`$ are then called the weights of the functions $`f`$ with respect to whatever discrete subgroup of $`Sp(2,𝐙)`$ is considered.
It follows from (6.8) and (6.9) that, if the $`N`$ functions $`f`$ are linearly independent, then the matrices $`D`$ represent the relevant group:
$$D(B)D(B^{})=D(BB^{}).$$
$`(6.10)`$
Furthermore
$$D_{\alpha \beta }(I_2)=\delta _{\alpha \beta },D_{\alpha \beta }(I_2)=(1)^{k_1k_2}\delta _{\alpha \beta }.$$
$`(6.11)`$
But the theta functions (6.3) are not necessarily linearly independent as they may be related by permutation matrices (which are real). For example,
$$\mathrm{\Theta }_\alpha (\tau )P_{\alpha \beta }=\mathrm{\Theta }_\beta (\tau )\text{where}P_{\alpha \beta }=\delta _{\lambda _\alpha +\lambda _\beta ,0},$$
$`(6.12)`$
and it is understood in the Kronecker delta function that the equality to zero is mod $`\mathrm{\Lambda }`$.
It follows from rearrangement of (6.4) and (6.5) that the $`|Z(\mathrm{\Lambda })|`$ theta functions (6.3) constitute a vector modular form with weights $`(k_1,k_2)=(b^+/2,b^{}/2)`$ and that, for them,
$$D_{\alpha \beta }(S)=\frac{\mathrm{e}^{2\pi \mathrm{i}\eta /8}}{\sqrt{|Z(\mathrm{\Lambda })|}}\mathrm{e}^{2\pi \mathrm{i}\lambda _\alpha .\lambda _\beta },$$
$`(6.13)`$
$$D_{\alpha \beta }(T)=\mathrm{e}^{\pi \mathrm{i}\lambda _\alpha ^2}\delta _{\alpha \beta }\text{if }\mathrm{\Lambda }\text{ is even,}$$
$`(6.14a)`$
$$D_{\alpha \beta }(T^2)=\mathrm{e}^{2\pi \mathrm{i}\lambda _\alpha ^2}\delta _{\alpha \beta }\text{otherwise.}$$
$`(6.14b)`$
It is important to notice that in view of the assumed properties of the lattice $`\mathrm{\Lambda }`$, these matrix elements are independent of the choice of representatives $`\lambda _\alpha `$ of the cosets (6.2). Furthermore they are independent of the parameters in the $`b^+b^{}`$-dimensional moduli space of the matrices $`G`$.
The matrices (6.13) and (6.14) are unitary, obviously so for $`D(T)`$ and $`D(T^2)`$. The unitarity of $`D(S)`$ follows from the orthogonality properties of the characters of the abelian group $`Z(\mathrm{\Lambda })`$ which appear as the phases $`\mathrm{exp}(2\pi \mathrm{i}\lambda _\alpha \lambda _\beta )`$. Actually the unitarity of $`D(S)`$ was used in the derivation of (6.13) from (6.4).
By inspection,
$$D(S)^T=D(S),D^{}(S)=(1)^{\eta /2}PD(S)=(1)^{\eta /2}D(S)P$$
$`(6.15)`$
where $`P`$ is defined by (6.12). Hence
$$D(S)^2=D(S)D(S)^{}(1)^{\eta /2}P=(1)^{\eta /2}P=D(I_2)P$$
using unitarity and (6.11). As $`P`$ equals unity on theta functions, (6.12), the first of the relations (6.6) is therefore checked on the theta functions (6.3). The other identity is more interesting. By (6.4) and (6.5),
$$\left(\left(D(S)D(T)\right)^2\right)_{\alpha \beta }=\frac{\mathrm{e}^{4\pi \mathrm{i}\eta /8}}{|Z(\mathrm{\Lambda })|}\underset{\gamma }{}=0^{|Z(\mathrm{\Lambda })|1}\mathrm{e}^{\pi \mathrm{i}(2\lambda _\alpha \lambda _\gamma +\lambda _\gamma ^22\lambda _\gamma \lambda _\beta +\lambda _\beta ^2)}.$$
Rearranging the sum over $`\gamma `$, this can be written in the form
$$\left(\left(D(S)D(T)\right)^2\right)_{\alpha \beta }=\frac{\mathrm{\Psi }(\mathrm{\Lambda })\mathrm{e}^{2\pi \mathrm{i}\eta (\mathrm{\Lambda })/8}}{\sqrt{|Z(\mathrm{\Lambda })|}}\mathrm{e}^{\pi \mathrm{i}(2\lambda _\alpha \lambda _\beta +\lambda _\alpha ^2)},$$
$`(6.16)`$
where
$$\mathrm{\Psi }(\mathrm{\Lambda })\frac{\mathrm{e}^{2\pi \mathrm{i}\eta (\mathrm{\Lambda })/8}}{\sqrt{|Z(\mathrm{\Lambda })|}}\underset{\gamma =0}{\overset{|Z(\mathrm{\Lambda })|1}{}}\mathrm{e}^{\pi \mathrm{i}(\lambda _\gamma \lambda _\alpha \lambda _\beta )^2}=\frac{\mathrm{e}^{2\pi \mathrm{i}\eta (\mathrm{\Lambda })/8}}{\sqrt{|Z(\mathrm{\Lambda })|}}\underset{\gamma =0}{\overset{|Z(\mathrm{\Lambda })|1}{}}\mathrm{e}^{\pi \mathrm{i}(\lambda _\gamma )^2}.$$
$`(6.17)`$
So the quantity $`\mathrm{\Psi }(\mathrm{\Lambda })`$ is independent of $`\alpha `$ and $`\beta `$ by virtue of the way addition of cosets in (6.2) realises the group $`Z(\mathrm{\Lambda })`$. Now we observe that
$$\left(\left(D(S)D(T)\right)^{}\right)_{\alpha \beta }=\left(D^{}(S)D^{}(T)\right)_{\beta \alpha }=\frac{1}{\sqrt{|Z(\mathrm{\Lambda })|}}\mathrm{e}^{\frac{1}{8}2\pi \mathrm{i}\eta }\mathrm{e}^{\pi \mathrm{i}(2\lambda _\alpha \lambda _\beta \lambda _\alpha ^2)}.$$
This can be made proportional to $`(D(S)D(T))^2`$ by means of the matrix $`P`$, which changes the sign of the vector $`\lambda _\beta `$:
$$\left(\left(D(S)D(T)\right)^{}P\right)_{\alpha \beta }=\left(D^{}(S)D^{}(T)P\right)_{\beta \alpha }=\frac{1}{\sqrt{|Z(\mathrm{\Lambda })|}}\mathrm{e}^{\frac{1}{8}2\pi \mathrm{i}\eta }\mathrm{e}^{\pi \mathrm{i}(2\lambda _\alpha \lambda _\beta +\lambda _\alpha ^2)}.$$
Comparing the last equation with (6.16) we find that
$$\left(D(S)D(T)\right)^2=\mathrm{e}^{4\pi \mathrm{i}\eta /8}\mathrm{\Psi }\left(D(S)D(T)\right)^{}P.$$
The matrix $`D(S)D(T)`$ is unitary, and this implies that
$$\left(D(S)D(T)\right)^3=\mathrm{e}^{4\pi \mathrm{i}\eta /8}\mathrm{\Psi }P.$$
But, by (6.11), $`D(I_2)=\mathrm{exp}(4\pi \mathrm{i}\eta /8)`$ and therefore the result we are looking for is
$$D(I_2)\left(D(S)D(T)\right)^3=\mathrm{\Psi }(\mathrm{\Lambda })P.$$
This has to be equal to the identity in order to satisfy the relation $`(ST)^3=I_2`$, (6.6). When acting on theta functions the matrix $`P`$ is the identity, and this implies that $`\mathrm{\Psi }(\mathrm{\Lambda })`$, (6.17), must have a particular value, namely unity. Hence
$$\frac{1}{\sqrt{|Z(\mathrm{\Lambda })|}}\underset{\gamma =0}{\overset{|Z(\mathrm{\Lambda })|1}{}}\mathrm{e}^{\pi \mathrm{i}\lambda _\gamma ^2}=\mathrm{e}^{\frac{1}{8}2\pi \mathrm{i}\eta (\mathrm{\Lambda })}.$$
$`(6.18)`$
This is Milgram’s formula, relating the signature of any even lattice to its structure, but so far only proven when $`b^+`$ and $`b^{}`$ and hence their difference, $`\eta `$, the signature of $`\mathrm{\Lambda }`$, are all even. But it is easy to check (6.18) explicitly for the even lattice $`\mathrm{\Lambda }`$ given by $`\sqrt{2}𝐙`$, so $`Z(\mathrm{\Lambda })=Z_2`$, for either sign in the natural scalar product. Furthermore, if $`\mathrm{\Lambda }_1\mathrm{\Lambda }_2`$ is the even lattice which is the orthogonal sum of the two even lattices $`\mathrm{\Lambda }_1`$ and $`\mathrm{\Lambda }_2`$ it is easy to check that $`\mathrm{\Psi }(\mathrm{\Lambda }_1\mathrm{\Lambda }_2)=\mathrm{\Psi }(\mathrm{\Lambda }_1)\mathrm{\Psi }(\mathrm{\Lambda }_2)`$. Thus if two of the three lattices satisfies Milgram’s formula so does the third. Using this (6.18) can then be deduced for any even lattice.
In particular, Milgram’s formula, (6.18), implies that the signature of an even unimodular lattice is a multiple of eight as stated just after equation (4.8).
## 7. Theta functions and odd unimodular lattices
If the lattice $`\mathrm{\Lambda }`$ is unimodular, the group $`Z(\mathrm{\Lambda })`$ consists just of the unit element. Hence the previous construction (6.3) yields just one theta function and the matrices $`D`$ are simply phase factors, by unitarity, usually called multipliers.
If, in addition, $`\mathrm{\Lambda }`$ is even, the signature is a multiple of eight, by Milgram’s formula (6.18). Hence the multipliers (6.13) and (6.14a) are trivial:
$$D(S)=D(T)=1.$$
If, instead, the unimodular $`\mathrm{\Lambda }`$ is odd there is no restriction on its signature. Again there is just one theta function supporting transformations with respect to the Hecke subgroup, $`\mathrm{\Gamma }_\theta `$, generated by $`S`$ and $`T^2`$ but with non-trivial multipliers
$$D(S)=\mathrm{e}^{2\pi \mathrm{i}\eta /8}\text{ and }D(T^2)=1.$$
$`(7.1)`$
However, associated with an odd unimodular lattice $`\mathrm{\Lambda }`$ is the even lattice $`\mathrm{\Lambda }_{EVEN}`$ for which the group $`Z(\mathrm{\Lambda }_{EVEN})`$ consists of four elements, (4.11). The results of the previous section suggest that there are three more theta functions which, altogether, support an action under the full group $`SL(2,𝐙)`$ which we now investigate.
It is this construction that will be precisely relevant to the Maxwell partition functions associated with four-manifolds of Type II.
Since
$$SL(2,𝐙)=\mathrm{\Gamma }_\theta T\mathrm{\Gamma }_\theta ST\mathrm{\Gamma }_\theta ,$$
$`(7.2)`$
$`\mathrm{\Gamma }_\theta `$ is a subgroup of index three in $`SL(2,𝐙)`$ and so it is natural to expect that only three theta functions are needed for the complete action. Indeed this is true and the fourth theta function can be chosen to be orthogonal to the three whilst supporting an $`SL(2,𝐙)`$ action on its own (if it does not vanish), as we shall see.
The following vectors are chosen as representatives of the cosets in the decomposition (4.12) of $`\mathrm{\Lambda }_{TOTAL}=\mathrm{\Lambda }_{EVEN}^{}`$:
$$0,\lambda _v,\lambda _s=c/2,\lambda _t=\lambda _v+c/2,$$
$`(7.3)`$
where $`\lambda _v`$ is simply any element of $`\mathrm{\Lambda }_{ODD}`$. Then the defining properties of the characteristic vector $`c`$ lead to the conclusion that
$$\frac{1}{2}\underset{\gamma =0,v,s,t}{}\mathrm{e}^{\pi \mathrm{i}\lambda _\gamma ^2}=\frac{(11+2\mathrm{e}^{\pi \mathrm{i}c^2/4})}{2}=\mathrm{e}^{2\pi \mathrm{i}c^2/8}$$
Using this, Milgram’s formula (6.18) reduces to (4.8) which is thereby proven.
The notation in (7.3) is the same as one customarily used for the decomposition of the weight lattice of the Lie algebra $`so(2r)`$. Corresponding to this $`\mathrm{\Lambda }`$ is the hypercubic lattice $`𝐙^r`$. This is indeed a special case and many of the properties familiar in that case will turn out to be true in much greater generality.
This is despite the fact that the theta functions associated with the odd unimodular Euclidean lattice $`𝐙^r`$ take the form
$$\left(\underset{n𝐙}{}\mathrm{e}^{\pi \mathrm{i}\tau n^2}\right)^r$$
and no such factorisation occurs for more complicated odd unimodular lattices such as $`\mathrm{\Gamma }_{8n+4}`$ given by the construction (4.13), starting from the hypercubic lattice $`\mathrm{\Lambda }=𝐙^{8n+4}`$, or the theta functions associated with $`𝐙^r`$ with indefinite scalar product.
Corresponding to the decomposition (7.3) there are four theta functions, denoted, respectively $`\mathrm{\Theta }_0(\tau ),\mathrm{\Theta }_v(\tau ),\mathrm{\Theta }_s(\tau )`$ and $`\mathrm{\Theta }_t(\tau )`$. It is easy to evaluate the four-by-four matrices $`D(S)`$ and $`D(T)`$ given by (6.13), (6.14a) and (7.3) in terms of $`\omega =\mathrm{e}^{2\pi \mathrm{i}\eta (\mathrm{\Lambda })/8}`$ as
$$D_{\alpha \beta }(S)=\frac{\omega }{2}e^{2\pi \mathrm{i}\lambda _\alpha \lambda _\beta }=\frac{\omega }{2}\left(\begin{array}{cccc}1& 1& 1& 1\\ 1& 1& 1& 1\\ 1& 1& \omega ^2& \omega ^2\\ 1& 1& \omega ^2& \omega ^2\end{array}\right)$$
$`(7.4)`$
$$D(T)=\text{diag}(1,1,\omega ^1,\omega ^1).$$
$`(7.5)`$
It immediately follows from the definition (6.9) with (7.4) and (7.5) that
$$\widehat{A}(S)\left(\mathrm{\Theta }_s\mathrm{\Theta }_t\right)=\omega ^3\left(\mathrm{\Theta }_s\mathrm{\Theta }_t\right)\widehat{A}(T)\left(\mathrm{\Theta }_s\mathrm{\Theta }_t\right)=\omega ^1\left(\mathrm{\Theta }_s\mathrm{\Theta }_t\right)$$
Hence the linear combination $`\mathrm{\Theta }_s\mathrm{\Theta }_t`$ supports an action of $`SL(2,𝐙)`$ with the multipliers indicated. Quite often $`\mathrm{\Theta }_s\mathrm{\Theta }_t`$ vanishes, for example if the signature $`\eta (\mathrm{\Lambda })`$ is odd or if $`\mathrm{\Lambda }`$ is Euclidean and hypercubic.
The sub-space orthogonal to $`\mathrm{\Theta }_s\mathrm{\Theta }_t`$ is three dimensional and spanned by $`\mathrm{\Theta }_0`$, $`\mathrm{\Theta }_v`$ and $`\mathrm{\Theta }_s+\mathrm{\Theta }_t`$. Of course the theta function associated with the original odd unimodular lattice $`\mathrm{\Lambda }`$ is
$$\mathrm{\Theta }=\mathrm{\Theta }_0+\mathrm{\Theta }_v$$
and it is this which supports the action of the Hecke subgroup $`\mathrm{\Gamma }_\theta `$. It is convenient to define
$$\mathrm{\Theta }^T(\tau )\widehat{A}(T)\mathrm{\Theta }=\mathrm{\Theta }_0(\tau )\mathrm{\Theta }_v(\tau )$$
using (7.5), and
$$\mathrm{\Theta }^{ST}(\tau )\widehat{A}(S)\mathrm{\Theta }^T=\widehat{A}(ST)\mathrm{\Theta }=\omega (\mathrm{\Theta }_s(\tau )+\mathrm{\Theta }_t(\tau )),$$
using (7.4). Then $`\mathrm{\Theta }`$, $`\mathrm{\Theta }^T`$ and $`\mathrm{\Theta }^{ST}`$ form an alternative orthogonal basis for the three dimensional subspace in which the action of $`S`$ is now represented by permutation matrices.
It is this basis that is relevant to the distinct Maxwell partition functions arising when the space-time four-manifold is of Type II so that the unimodular lattice $`F_2(_4,𝐙)`$, (3.2), is unimodular and odd. The partition function for backgrounds supporting complex scalar wave functions is
$$Z_{SCALAR}(\tau )=\mathrm{\Delta }(\tau )\mathrm{\Theta }(\tau )$$
while the partition function for backgrounds supporting complex spinor wave functions is
$$Z_{SPINOR}(\tau )=\mathrm{\Delta }(\tau )\mathrm{\Theta }_{w/2}(\tau )=\mathrm{\Delta }(\tau )(\mathrm{\Theta }_s(\tau )+\mathrm{\Theta }_t(\tau ))=\omega ^1\mathrm{\Delta }(\tau )\mathrm{\Theta }^{ST}(\tau )$$
using (5.14) and assuming that the respective electric charges $`q_B`$ and $`q_F`$ are equal.
Of course $`\mathrm{\Delta }(\tau )\mathrm{\Theta }^T(\tau )`$ is simply $`Z_{SCALAR}`$ with the angle $`\theta `$, (2.3), (2.4) increased by $`2\pi `$.
We have seen that the four dimensional space of theta functions associated with the even lattice, $`\mathrm{\Lambda }_{EVEN}`$ constructed from any odd unimodular lattice via (4.10) usually decomposes under the $`SL(2,𝐙)`$ action as $`4=3+1`$, though sometimes, for example when the signature $`\eta (\mathrm{\Lambda })`$ is odd, it is $`4=3+0`$ as $`\mathrm{\Theta }_s\mathrm{\Theta }_t`$ vanishes. In fact when $`\eta (\mathrm{\Lambda })`$ is a multiple of four there is a further decomposition $`3=2+1`$ whose details depend on whether $`\eta (\mathrm{\Lambda })/4`$ is odd or even. This is because of the occurrence of the two new unimodular lattices (4.13) which are odd or even as $`\eta (\mathrm{\Lambda })/4`$ is. Associated with these are the theta functions $`\mathrm{\Theta }_0+\mathrm{\Theta }_s`$ and $`\mathrm{\Theta }_0+\mathrm{\Theta }_t`$ obeying
$$\widehat{A}(S)\left(\mathrm{\Theta }_0+\mathrm{\Theta }_s\right)=\omega (\mathrm{\Theta }_0+\mathrm{\Theta }_s),$$
$$\widehat{A}(S)\left(\mathrm{\Theta }_0+\mathrm{\Theta }_t\right)=\omega (\mathrm{\Theta }_0+\mathrm{\Theta }_t),$$
where, now, $`\omega =(1)^{\eta (\mathrm{\Lambda })/4}`$. Also
$$\widehat{A}(T)(\mathrm{\Theta }_0+\mathrm{\Theta }_s)=\mathrm{\Theta }_0+\omega \mathrm{\Theta }_s,$$
$$\widehat{A}(T)(\mathrm{\Theta }_0+\mathrm{\Theta }_t)=\mathrm{\Theta }_0+\omega \mathrm{\Theta }_t.$$
So, when $`\eta (\mathrm{\Lambda })`$ is a multiple of eight, so $`\omega =1`$, the three dimensional subspace contains the modular function
$$2\mathrm{\Theta }_0+\mathrm{\Theta }_s+\mathrm{\Theta }_t=\mathrm{\Theta }+\mathrm{\Theta }^T+\mathrm{\Theta }^{ST},$$
with multipliers $`D(S)=1=D(T)`$. On the other hand when $`\eta (\mathrm{\Lambda })8𝐙+4`$, so $`\omega =1`$ it is
$$2\mathrm{\Theta }_v\mathrm{\Theta }_s\mathrm{\Theta }_t=\mathrm{\Theta }\mathrm{\Theta }^T+\mathrm{\Theta }^{ST}$$
which has multipliers $`D(S)=D(T)=1`$.
## 8. Discussion
We have seen that the response of the extended partition function $`Z`$ associated with a four-manifold $`_4`$ depends on the topological properties of $`_4`$. For example, if $`_4`$ is of type I or III (according to the terminology of section 3), so that $`F_2(_4,𝐙)`$ is an even unimodular lattice, there is just one partition function and it transforms under the full electromagnetic duality group $`SL(2,𝐙)`$ acting on the dimensionless electromagnetic coupling $`\tau `$, (2.4), by fractional linear transformations. If $`_4`$ is of type II so that $`F_2(_4,𝐙)`$ is an odd unimodular lattice there are three possible partition functions, two appropriate to fluxes supporting complex scalar wave functions and the third, complex spinor wave functions. Individually these transform under three distinct but conjugate subgroups of $`SL(2,𝐙)`$ of index three (one of which is the Hecke subgroup). The full electromagnetic duality $`SL(2,𝐙)`$ is realised by permutations of these three, as explained above.
In talking of transformations we allow for the possiblity of non-zero “modular weights”. We saw after equation (6.12) that the modular weights of the theta functions were given by $`b^\pm (_4)/2`$. Witten argued that the prefactors due to the Van Vleck determinants (5.1) had modular weights both given by $`(1b_1(_4))/2`$. Hence the total modular weights of $`Z`$ are
$$\frac{1b_1(_4)+b^\pm (_4)}{2}=\frac{\chi (_4)\pm \eta (_4)}{4},$$
$`(8.1)`$
using (2.13) and (4.2). These numbers possess several interesting properties. First they are integers (or maybe half-integers in which case we should properly talk of metaplectic versions of $`SL(2,𝐙)`$). Secondly they are rather special topological numbers possessing the property of “locality” in the sense that the can be expressed as integrals over $`_4`$ of local densities.
Suppose space-time $`_4`$ exhibits a discrete $`Z_2`$ symmetry with no fixed points. Then the quotient $`_4/Z_2`$ obtained by identifying related points is also a four-manifold and it follows from the locality properties that
$$\chi \left(_4/Z_2\right)=\frac{\chi (_4)}{2}.$$
$`(8.2)`$
If the quotient $`_4/Z_2`$ is Poincaré dual, the same relation (8.2) holds for the Hirzebruch signature $`\eta `$ and hence for the modular weights (8.1).
If the partition function considered is a “strict” one, that is it is indeed a trace, (2.9), so that $`_4`$ has the form $`S^1\times _3`$, then, as $`S^1S^1/Z_2`$, $`_4_4/Z_2`$, where the $`Z_2`$ relates diametrically opposite points on the circle. Consequently, by (8.2), $`\chi `$, $`\eta `$ and the modular weights (8.1) all vanish and the partition function is actually invariant, in agreement with initial expectation. It is easy to check by calculation that $`F_2(S^1\times _3,𝐙)`$ is an even lattice so that type II is ruled out.
In general, the extended partition functions described in section 2 as being associated with more general four-manifolds are only modular covariant as the weights (8.1) need not vanish. This result is clearly of interest, but it has to be admitted that it is unclear what its physical meaning is. One reason is that a choice has to be made in selecting the Euclidean metric on $`_4`$. Fortunately the dependence on that choice is not too strong. For example, if the metric is altered by a Weyl rescaling,
$$g_{\mu \nu }(x)\lambda (x)^2g_{\mu \nu }(x),$$
$`(8.3)`$
then $`F`$ is unaltered if $`F`$ itself is. Consequently the matrix $`G`$ defined by (5.3) is unchanged as is the complete theta function. Perhaps this is related to the fact that when $`b^+`$ or $`b^{}`$ vanishes, then $`G=\pm Q`$ and so all dependence on the metric disappears from the theta function (which is then holomorphic or anti-holomorphic in $`\tau `$). If $`b^+`$ and $`b^{}`$ both vanish the action (2.1) is independent of $`\theta `$, as mentioned earlier. Now we see that the theta function is then simply a constant independent of $`\tau `$.
Of course, if $`_4`$ supports a Minkowski metric, its Euler number, $`\chi `$, vanishes and the modular weights (8.1) become equal and opposite.
We should like to emphasise that the breakdown of $`SL(2,𝐑)`$ to its discrete subgroup $`SL(2,𝐙)`$ was a consequence simply of the Dirac quantisation condition without recourse to the Zwanziger-Schwinger condition, which, accordingly, seems less fundamental.
In this paper we have endeavoured to elucidate the conceptual and mathematical structure in what seems to be the simplest possible context for electromagnetic duality. It would be physically interesting to extend the work in many different ways, for example: 1) to consider a larger number of abelian gauge fields on $`_4`$.
2) to consider space-times $`_{4k}`$ with $`2k`$-form field strengths with potential couplings to $`(2k2)`$-branes, possibly spinning.
3) to consider space-times $`_{4k+2}`$ with several $`2k+1`$-form field strengths.
4) to consider open space-times with boundary.
5) to consider the effect of space-time topology on conventional superstring theories. Undoubtedly yet more mathematical tools would have to brought into play.
Acknowledgements
DIO wishes to thank Gary Gibbons, Stephen Howes, Nadim Mahassen, Boris Pioline for discussions. He also wishes to thank the Mittag-Leffler Institute, UNESP Institute for Theoretical Physics (São Paulo) and the Instituut voor Theoretische Fysica, Utrecht for hospitality and their members for congenial and enlightening discussions. MA’s research has been supported by EPSRC and subsequently by PPARC through its Special Programme Grant PPA/G/S/1998/00613. We both wish to thank TMR grant FMRX-CT96-0012 for assistance. References M Alvarez and DI Olive; “The Dirac quantisation condition for fluxes on Four-manifolds”, hep-th/9906093, to appear in Commun. Math. Phys. O Alvarez; “Topological Quantisation and Cohomology”, Commun. Math. Phys. 100 (1985) 279-309 E Cremmer and B Julia; “The $`SO(8)`$ Supergravity”, Nucl. Phys. B159, (1979) 141-212 A D’Adda, R Horsley and P Di Vecchia; “Supersymmetric Monopoles and Dyons”, Phys. Lett. B76 (1978) 298-302. D Bohm and Y Aharonov; “Significance of electromagnetic potentials in the quantum theory”, Phys. Rev. 115 (1959), 485 PAM Dirac; “Quantised singularities in the electromagnetic field”, Proc. Roy. Soc. A33 (1931) 60-72 PAM Dirac; “Gauge invariant formulation of Quantum Electrodynamics”, Canadian Journal of Physics 33 (1955), 650-660 RP Feynman and AR Hibbs; Quantum Mechanics and Path Integrals, Chapter 10, McGraw-Hill, 1965 RP Feynman; Statistical Mechanics, Chapter 3, Benjamin/Cummings, 1972 MK Gaillard and B Zumino; “Duality Rotations for Interacting Fields”, Nucl. Phys. B193 (1981), 221-244 M Green, J Schwarz and E Witten; Superstring Theory, Vol. 2, Appendix 9B, Cambridge University Press, 1987 M Kalb and P Ramond; “Classical direct interstring action”, Phys. Rev. D9 (1974), 2273-2284 V Kac; Infinite dimensional Lie algebras, Chapter 13, Cambridge University Press, 1990 HA Kramers and GH Wannier; “Statistics of the Two-Dimensional Ferromagnet I”, Phys. Rev.60 (1941), 252-276 HB Lawson and M-L Michelsohn; Spin Geometry, Princeton Mathematical Series 38, Princeton, 1989 J Milnor, D. Husemoller; Symmetric bilinear forms, Ergebnisse der Mathematik und ihrer Grenzgebiete, Band 73, Springer-Verlag, 1973. C Montonen and DI Olive; “Magnetic monopoles as gauge fields?”, Phys. Lett. B72 (1977), 117-120. H Osborn; “Topological Charges for $`N=4`$ Supersymmetric Gauge Theories and Monopoles of Spin 1”, Phys. Lett. B83 (1979), 321-326. A Schwarz; Topology for Physicists, Springer, 1994. A Sen; “Dyon-monopole bound states, self-dual harmonic forms on the multi-monopole moduli space, and $`SL(2,ZZ)`$ invariance in string theory”, Phys. Lett. B329 (1994), 217-221 E Verlinde; “Global aspects of electric-magnetic duality”, Nucl. Phys. B455 (1995), 211-228. K Wilson; “Confinement of quarks”, Phys. Rev. D10 (1974), 2445-2459 E Witten and DI Olive; “Supersymmetry Algebras that Include Topological Charges”, Phys. Lett. B78 (1978), 97-101. E Witten; “On S-duality in abelian gauge theory”, Selecta Math (NS) 1 (1995), 383-410, hep-th/9505186 TT Wu and CN Yang; “Concept of non-integrable phase factors and global formulation of gauge fields”, Phys. Rev. D12 (1975), 3845-3857.
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# 1. Introduction
## 1. Introduction
The generally accepted ‘triviality’ of $`\lambda \mathrm{\Phi }^4`$ theories in four space-time dimensions is usually interpreted within leading-order perturbation theory in a very intuitive way. One starts with the perturbative one-loop $`\beta `$function
$$\beta _{\mathrm{pert}}(\lambda )=\frac{3\lambda ^2}{16\pi ^2}+𝒪(\lambda ^3)$$
(1.1)
and integrates the differential equation
$$\frac{dx}{x}=\frac{d\lambda }{\beta (\lambda )}$$
(1.2)
between a fixed energy scale $`x=\mu `$ and the ‘Landau pole’ $`x=\mathrm{\Lambda }`$ where $`\lambda (x)=+\mathrm{}`$. In this way, at energy scales $`𝒪(\mu )`$, the theory is governed by a 1PI 4-point function
$$\lambda \lambda (\mu )=\frac{32\pi ^2}{3\mathrm{ln}\frac{\mathrm{\Lambda }^2}{\mu ^2}}$$
(1.3)
that would vanish in the continuum limit where the ultraviolet cutoff $`\mathrm{\Lambda }\mathrm{}`$ within all loop diagrams. In this picture, differently from the original renormalization-group approach where $`\lambda (\mu )`$ is changed along a given integral curve of Eq.(1.2) by changing $`\mu `$, one considers all possible integral curves at the same time. After that, one changes $`\lambda `$, at fixed $`\mu `$, depending on the magnitude of the Landau pole associated with the various integral curves. In this way, the ‘bare’ theory is always defined in the infinite-coupling limit, i.e. as in an Ising model. However, at any finite scale $`\mu `$, the continuum limit corresponds to a vanishingly small interaction strength so that perturbation theory in the small parameter $`\lambda (\mu )`$ should provide very accurate predictions for low-energy physical observables.
By adopting this view of ‘triviality’ in the spontaneously broken phase of $`(\lambda \mathrm{\Phi }^4)_4`$, the euclidean propagator for the Higgs particle should approach the form
$$\stackrel{~}{G}(q)\frac{Z_{\mathrm{prop}}}{q^2+M_h^2}$$
(1.4)
with a residue
$$Z_{\mathrm{prop}}=1|𝒪(\lambda )|1$$
(1.5)
consistent with the Källen-Lehmann decomposition that dictates the spectral function $`\rho _h(s)`$ to approach $`\delta (sM_h^2)`$ in the continuum limit of a ‘trivial’ theory.
This simple picture neglects, however, that the origin of spontaneous symmetry breaking is not necessarily of perturbative nature. Indeed, one may be faced with non-analytic $`1/\lambda `$ effects so that, even with an infinitesimal two-body interaction strength, there may be non-perturbative effects. To understand the possible implications, one should remember the case of superconductivity. This is due to the basic instability of a normal Fermi system in the presence of an infinitesimally small attractive interaction between electrons. Due to the presence of a very large density of quantum states near the Fermi surface, superconductivity is a non-perturbative phenomenon.
In the case of spontaneous symmetry breaking with an elementary scalar field the delicate issue concerns the limit $`q0`$ that requires some care in the case of a macroscopic occupation of the same quantum state, i.e. of Bose-Einstein condensation. Just for this reason, and contrary to the most naive expectations, the approach to the continuum limit in the spontaneously broken phase of $`(\lambda \mathrm{\Phi }^4)_4`$ theories may contain unexpected features. These are discovered whenever one takes seriously ‘triviality’ as a technical statement controlling the approach to the continuum theory and, thus, supporting the idea of a trivially free fluctuation field with a gaussian quantum measure for $`\mathrm{\Lambda }\mathrm{}`$. In fact, as pointed out in refs., assuming gaussian (and post-gaussian) ansatz for the ground state wave functional, one finds
$$\frac{\mathrm{\Gamma }_2(0)}{M_h^2}=𝒪(1/\mathrm{ln}\mathrm{\Lambda })0$$
(1.6)
once the zero-momentum two-point function is computed through the effective potential
$$\mathrm{\Gamma }_2(0)\frac{d^2V_{\mathrm{eff}}}{d\varphi ^2}|_{\varphi =v}$$
(1.7)
As pointed out in refs., Eq. (1.6) requires a non-trivial re-scaling
$$ZZ_\varphi =𝒪(\mathrm{ln}\mathrm{\Lambda })\mathrm{}$$
(1.8)
of the vacuum field $`\varphi `$
$$\varphi _R^2\frac{\varphi ^2}{Z_\varphi }$$
(1.9)
in order to match the quadratic shape of the effective potential with the physical mass $`M_h`$ defined from the $`q0`$ behaviour of the propagator
$$\frac{d^2V_{\mathrm{eff}}}{d\varphi _R^2}|_{\varphi _R=v_R}M_h^2.$$
(1.10)
As such, $`Z=Z_\varphi `$ is quite distinct from the ‘trivial’ re-scaling $`Z=Z_{\mathrm{prop}}`$ in Eq.(1.5), and one may obtain a continuum limit where, although $`M_h`$ vanishes in units of the bare $`v`$, both $`M_h`$ and $`v_R`$ are finite quantities (with potentially important implications for the commonly quoted upper bounds on the Higgs mass from ‘triviality’ ).
The result in Eq.(1.6) is also striking for the following reason. The euclidean value $`q=0`$ corresponds, indeed, to the single point $`(q_o,𝐪)=0`$ in the continuum theory. However, in the cutoff theory, Eq.(1.6) implies that there is a region of 3-momenta say $`𝐪^2𝒪(1/\mathrm{ln}\mathrm{\Lambda })M_h^2`$ where the energy spectrum is not $`\stackrel{~}{E}(𝐪)=\sqrt{𝐪^2+M_h^2}`$. This region, although infinitesimal in units of $`M_h`$, can have a physical meaning (for $`M_h=𝒪(10^2)`$ GeV, think of the values $`|𝐪|10^5`$ eV/c corresponding to wavelengths much larger than 1 cm).
The previous result was obtained in the formalism of the gaussian ( and post-gaussian) approximations to the effective potential. In the next section, we shall outline a simple semi-perturbative argument that leads to the same conclusions and can help to understand in a more intuitive way the singular nature of the limit $`q0`$ when approaching the continuum theory from the broken phase.
## 2. A simple semi-perturbative calculation
Let us consider a one-component $`\lambda \mathrm{\Phi }^4`$ theory
$$=\frac{1}{2}(\mathrm{\Phi })^2U(\mathrm{\Phi })$$
(2.1)
with a classical potential ($`\lambda >0`$)
$$U(\mathrm{\Phi })=\frac{1}{2}m_B^2\mathrm{\Phi }^2+\frac{\lambda }{4!}\mathrm{\Phi }^4$$
(2.2)
Classically, non-vanishing constant field configurations $`\varphi =\pm v`$ occur where $`U^{}(\pm v)=0`$. At these values, one gets a quadratic shape
$$U^{\prime \prime }(\pm v)=\frac{\lambda v^2}{3}M_h^2$$
(2.3)
that represents the well known classical result for the Higgs mass.
In the quantum theory, the question of vacuum stability is more subtle and one has to replace $`U(\varphi )`$ with the quantum effective potential $`V_{\mathrm{eff}}(\varphi )`$. However, the basic expectation is that the excitation spectrum of the broken phase will maintain the Lorentz-covariant form $`\stackrel{~}{E}(𝐪)=\sqrt{𝐪^2+M_h^2}`$ down to $`𝐪=0`$ so that $`M_h`$ should coincide with $`\stackrel{~}{E}(0)`$ that represents the energy-gap of the broken phase.
To check this prediction, let us consider the one-loop structure of the gap-equation for the euclidean two-point function
$$\stackrel{~}{G}^1(q)\mathrm{\Gamma }_2(q)=q^2+m_B^2+\frac{\lambda \varphi ^2}{2}+\frac{\lambda }{2}\frac{d^4k}{(2\pi )^4}\stackrel{~}{G}(k)\frac{\lambda ^2\varphi ^2}{2}\frac{d^4k}{(2\pi )^4}\stackrel{~}{G}(k)\stackrel{~}{G}(k+q)$$
(2.4)
together with the vanishing of the one-loop tadpoles, i.e.
$$T(\varphi )m_B^2+\frac{\lambda \varphi ^2}{6}+\frac{\lambda }{2}\frac{d^4k}{(2\pi )^4}\stackrel{~}{G}(k)=0$$
(2.5)
Eqs.(2.4) and (2.5) can be understood, for instance, as coupled minimization equations of the effective potential for composite operators introduced by Cornwall, Jackiw and Tomboulis . As such, they are non-perturbative, being equivalent to the all-order resummation of one-loop graphs with the tree-level propagator in the external background field $`\varphi `$
$$\stackrel{~}{G}(q)_{\mathrm{tree}}=\frac{1}{q^2+m_B^2+\frac{\lambda \varphi ^2}{2}}$$
(2.6)
Notice that for $`\varphi 0`$, the gap-equation would also contain one-particle reducible terms (i.e. not contained in $`\mathrm{\Gamma }_2(q)`$) proportional to the zero-momentum limit of the shifted field propagator $`\stackrel{~}{G}(0)`$. In this sense, by considering the 1PI gap-equation, we are assuming a non-singular zero-momentum limit, even for $`\varphi =\pm v`$ where $`\pm v`$ are the solutions of (2.5) and represent, to this order, the minima of the effective potential.
The possibility to solve simultaneously Eqs.(2.4) and (2.5), in the continuum limit of the regularized theory by using Eqs. (1.3) - (1.5), amounts to describe spontaneous symmetry breaking as a quantum phenomenon of vacuum instability consistently with (the intuitive interpretation of) rigorous quantum field theoretical results on $`(\lambda \mathrm{\Phi }^4)_4`$ theories.
By using (2.5) in (2.4) we get
$$\stackrel{~}{G}^1(q)\mathrm{\Gamma }_2(q)=q^2+\frac{\lambda v^2}{3}A(q)$$
(2.7)
where
$$A(q)=1\frac{3\lambda }{2}\frac{d^4k}{(2\pi )^4}\stackrel{~}{G}(k)\stackrel{~}{G}(k+q)$$
(2.8)
Notice that a form of the propagator as in Eq.(1.4) can be a solution of Eq.(2.4) only if the coupling $`\lambda `$ is understood as an infinitesimally small quantity, i.e. as in Eq.(1.3) with $`\mu M_h`$. In this case, for $`G(q)\frac{1}{q^2}`$ at large euclidean $`q^2`$, one finds $`A(q)=A(0)+𝒪(1/\mathrm{ln}\mathrm{\Lambda })`$ and one indeed gets a form $`\mathrm{\Gamma }_2(q)=q^2+const`$, up to terms vanishing in the continuum limit where $`\lambda 0`$. In this sense, renormalized perturbation theory should be considered an external input whose overall consistency with the all-order Eqs.(2.4) and (2.5) can only be checked a posteriori.
By using Eq. (1.4) in Eqs.(2.7) and (2.8) we obtain the leading-order expression
$$Z_{\mathrm{prop}}=\frac{1}{1+\frac{\lambda v^2}{3M_h^2}\frac{\lambda Z_{\mathrm{prop}}^2}{64\pi ^2}}$$
(2.9)
and the estimate of the Higgs mass
$$M_h^2=\frac{\lambda v^2Z_{\mathrm{prop}}}{3}(1\frac{3\lambda Z_{\mathrm{prop}}^2}{32\pi ^2}\mathrm{ln}\frac{\mathrm{\Lambda }^2}{M_h^2})$$
(2.10)
At this point we find a strong contradiction. Indeed, by using Eq. (1.3) and assuming (1.5) we obtain from (2.10)
$$M_h^2=𝒪(\lambda ^2v^2)$$
(2.11)
that when inserted in Eq.(2.9), does not produce a $`Z_{\mathrm{prop}}1`$ when $`\mathrm{\Lambda }\mathrm{}`$.
It is true that in a more conventional fixed-order perturbative calculation one would tend to regard (2.10) as the first two terms in the expansion of the renormalized coupling constant at a scale $`M_h`$ (i.e. the renormalization-group improved version of the classical result (2.3)). However, this cannot work in our case for two reasons. First, as anticipated, Eqs.(2.4) and (2.5) resum one-loop graphs with the tree propagator (2.6) to all orders, so that our results cannot be considered fixed order calculations. Second, using the simple form of the propagator (1.4) to solve Eq.(2.4) requires an infinitesimal $`\lambda `$, as in (1.3). However, if one tries to improve on (2.4) and (2.5), by introducing genuine two-loop terms, the intuitive interpretation of ‘triviality’ based on Eqs.(1.1)-(1.5) is destroyed. This is due to the presence of a (spurious) ultraviolet fixed point at finite coupling in the two-loop perturbative $`\beta `$ function so that, beyond a leading-order calculation, there are no reasons for the 1PI four point function $`\lambda (\mu )`$ to vanish in the continuum limit as in Eq.(1.3) and for Eqs. (1.4) and (1.5) to be valid.
The previous results suggest that, contrary to our assumption, either ‘triviality’ cannot be understood within leading-order perturbation theory, or one is faced with a singular $`\stackrel{~}{G}(0)`$ when $`\varphi =\pm v`$. Namely, a result $`M_h^2=𝒪(\lambda v^2)`$ from Eqs. (2.5) and (2.4), in order Eq.(2.9) to agree with (1.5), can only be obtained if there are one-particle reducible contributions to the propagator. These, at a generic value of $`\varphi `$, are proportional to the combination
$$R=T(\varphi )\stackrel{~}{G}(0)$$
(2.12)
and the case $`\varphi =\pm v`$ should be understood as a limiting procedure due to a possible divergence in $`\stackrel{~}{G}(0)`$. To this end, let us introduce $`\varphi ^2v^2(1+\delta )`$ with $`|\delta |1`$, and define
$$g\stackrel{~}{G}(0)\frac{\lambda v^2}{3}$$
(2.13)
Now a non vanishing $`R`$ requires
$$g\frac{1}{\delta }$$
(2.14)
implying that, at $`\varphi =\pm v`$ where $`\delta =0`$, there is a mode whose energy vanishes for $`𝐪0`$. For this reason, the energy-gap of the broken phase, obtained from $`\stackrel{~}{E}(𝐪)`$ for $`𝐪0`$, cannot be $`M_h`$.
## 3. The zero-momentum discontinuity
We believe that the discrepancy we have pointed out, is related to the subtleties of Bose-Einstein condensation in an almost ideal gas where there is a macroscopic occupation of the same quantum state. This phenomenon cannot be fully understood in a purely perturbative manner. However, as anticipated in the Introduction, the peculiarity of the zero-momentum limit is found in other approaches such as:
i) variational evaluations of the effective potential
ii) lattice computations of the propagator and of the zero-momentum susceptibility $`\chi ^1\mathrm{\Gamma }_2(0)`$ in the broken phase . These show that Eq.(1.4), although valid at higher momenta, has non-perturbative corrections for $`q0`$. Indeed, these become larger and larger by approaching the continuum limit and therefore cannot represent perturbative $`𝒪(\lambda )`$ effects
Due to this qualitative agreement with other approaches, the semi-perturbative calculation of sect.2 can be considered a reasonably self-contained treatment of the basic zero-momentum discontinuity.
The discrepancy persists in the case of an O(N) continuous-symmetry $`\lambda \mathrm{\Phi }^4`$ theory. To this end, one can check with a non-perturbative Gaussian effective-action approach . In this case, the minimization conditions of the gaussian effective potential provide suitable relations that replace the ‘triviality pattern’ Eqs.(1.3) and (1.5) and where the equivalent of $`T(\varphi )`$ plays the role of a mass term for the Goldstone bosons. By neglecting one-particle reducible contributions in the $`\sigma `$field propagator, one finds (N-1) massless fields with propagator $`D_\pi (q)=1/q^2`$ and a $`\sigma `$field two-point function of the type as in Eqs.(2.7) and (2.8). By following the steps of ref., it is not difficult to check that one gets in this way the same discrepancy as in sect.2 for the mass of the $`\sigma `$field. This can also be understood since, for the Goldstone bosons, the non-perturbative wave functional of ref. reproduces ‘triviality’ and is exact. In fact, it yields (N-1) non-interacting fields that decouple from each other and from the $`\sigma `$field. As a consequence, one gets effectively the same type of $`\sigma \sigma `$ interactions as in our discrete-symmetry case.
On the other hand, in the case of an O(N) continuous symmetry, the singular nature of the zero-momentum limit of the singlet-Higgs propagator is well known. It was first pointed out by Symanzik for the linear $`\sigma `$-model , and later on by Patashinsky and Pokrowsky and Anishetty et al. .
Symanzik’s analysis for the $`\sigma `$ field, although purely perturbative, displays the essential features of the phenomenon, i.e. the perturbative contradiction between finite 1PI vertex diagrams and finite Green’s functions that introduces the zero-momentum discontinuity. Just for this reason, he introduced two different notations, namely $`\mathrm{\Gamma }_\sigma (0)M^2`$ and $`\mathrm{\Gamma }_\sigma (q^2)(q^2+\overline{M}^2)`$, to emphasize that the limit $`q0`$ is not defined.
From ref., on the other hand, one can get a better feeling of what is actually going on. In fact, in the case of a spontaneously broken O(N) symmetry, the longitudinal susceptibility is found
$$\chi _{||}(𝐪)\frac{1}{|𝐪|}arctg\frac{|𝐪|}{2\kappa }$$
(3.1)
where
$$\kappa ^2(\varphi v)^2$$
(3.2)
Now, for any $`\varphi v`$, the limit $`𝐪0`$ yields a finite result. However, just in the case $`\varphi =v`$, $`\chi _{||}(𝐪)`$ becomes singular when $`𝐪0`$ as in our case.
Our results show no qualitative difference with respect to the continuous-symmetry case of ref. and, therefore, the agreement is not surprising. Indeed, the zero-momentum discontinuity does not depend on the existence of a continuous symmetry of the Lagrangian. Rather, its physical origin has to be searched in the presence of the scalar condensate, i.e. in the phenomenon of Bose-Einstein condensation that leads to a gap-less mode and to a long-range $`1/r`$ potential .
This can be understood as follows. As discussed in refs., variational approximations to $`V_{\mathrm{eff}}`$ describe spontaneous symmetry breaking as an infinitesimally weak first-order transition. This occurs when the mass-gap at $`\varphi =0`$, say $`0m^2m_c^2`$, is still positive, but in a ‘hierarchical’ relation
$$\frac{m^2}{M_h^2}=𝒪(\frac{1}{\mathrm{ln}\frac{\mathrm{\Lambda }^2}{\mu ^2}})$$
(3.3)
with the mass scale of the broken phase .
As discussed in ref., the gap-less mode with $`\stackrel{~}{E}(𝐪)=const.|𝐪|`$ is associated with the infinitesimal region of momenta
$$𝐪^2\frac{M_h^2}{\mathrm{ln}\frac{\mathrm{\Lambda }^2}{\mu ^2}}$$
(3.4)
that goes, indeed, into the single point $`(q_o,𝐪)=0`$ in the continuum limit $`\mathrm{\Lambda }\mathrm{}`$. However, in the cutoff theory this region defines the non-relativistic limit $`|𝐪|m`$ where $`m`$ is the mass of the quanta in the condensate (see Eq. (3.3)).
In this regime, any scalar condensate, whatever its origin may be, is a highly correlated structure with long-range order due to the coherence effects associated with the phase of the non-relativistic condensate wave-function . Therefore, for $`𝐪0`$, the deviations from a Lorentz-covariant energy spectrum $`\stackrel{~}{E}(𝐪)=\sqrt{𝐪^2+M_h^2}`$ are not surprising.
Now, for $`M_h=𝒪(10^2)`$ GeV, it is a matter of taste to decide whether, for instance, values $`|𝐪|10^5`$ eV/c (corresponding to wavelengths much larger than 1 cm) may be considered infinitesimal or not. In the case of a positive answer, these deviations from exact Lorentz-covariance on such scales should be taken seriously. Indeed, as discussed in , the associated extremely weak $`1/r`$ potential does not disappear when coupling the scalar fields to gauge bosons.
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# I Introduction
## I Introduction
The light current quark masses are baisc input of quantum chromodynamics(QCD). The inequality of the light quark masses, especially, $`m_um_d`$, breaks the isospin symmetry or charge symmetry. This breaking of isospin symmetry induces various measurable physics processes such as $`\pi ^0\eta `$, $`\mathrm{\Lambda }\mathrm{\Sigma }^0`$ mixing and $`\omega \pi ^+\pi ^{}`$ decay etc. In this paper, we will focus on the $`\omega \pi ^+\pi ^{}`$ decay, which is considered as the important source of charge symmetry breaking in nuclear physics.
In ref.(we will quote it as $`𝐈`$ hereafter) we have shown that the chiral expansion at vector meson energy region converges slowly. Therefore, a well-defined effective field theory describing the physics in this energy region must be available for calculation on high order terms of the chiral expansion and meson loops. Obviously, method of chiral perturbative theory(ChPT) is impractical to cupature the high order term contribution because the number of free parameter increases very rapidly as perturbative order raising. In $`𝐈`$, following the spirit of Manohar-Georgi(MG) model, we have constructed chiral constituent quark model(ChCQM) including lowest vector meson resonances. The advatanges of this approach are that high order contribution of the chiral expansion and $`N_c^1`$ expansion can be calculated consistently, and only fewer free parameters are required. Low energy limit and unitarity of the model are also examined successfully. In particular, it is a attractive property that although the leading order theoretical prediction does not match with experimental data, larger contribution from high order of the chiral expansion and pseudoscalar meson one-loop corrects theoretical prediction close to data very much. It is just the characteristic of the chiral expansion in this energy region. Therefore, in this the present paper, we also need to calculate $`\omega \pi ^+\pi ^{}`$ decay to include all high order contribution and pseudoscalar meson one-loop correction.
This research is also motivated by the following reasons:
i) In the most recent references, the $`\omega \pi ^+\pi ^{}`$ decay was treated as being dominant via $`\rho `$-resonance exchange, and the direct $`\omega \pi ^+\pi ^{}`$ coupling is neglected. It has been pointed out in ref. that the neglect of $`\omega `$ “direct” coupling to $`\pi ^+\pi ^{}`$ is not valid. It can be naturally understood since $`\pi \pi `$ can make up of vector-isovector system, whose quantum numbers are same to $`\rho `$ meson. Thus in an effective lagrangian based on chiral symmetry, every $`\rho `$ field can be replaced by $`\pi \pi `$ and does not conflict with symmetry. Although authors of ref. also pointed out that the present experimental data still can not be used to separate “direct” $`\omega \pi \pi `$ coupling from $`\omega \rho `$ mixing contribution in model-independent way, it is very interesting to perform a theorectical investigation on “direct” $`\omega \pi \pi `$ coupling contribution. We will show that the contribution from interfernce of “direct” $`\omega \pi \pi `$ coupling and $`\omega \rho `$ is about 15$`\%`$. Thus “direct” $`\omega \pi \pi `$ coupling can not be neglected indeed.
ii) The present study involves the investigation of $`\rho ^0\omega `$ mixing, which has been an active subject\[5–14\]. The mixing amplitude for on-mass-shell vector mesons has been observed directly in the measurement of the pion form-factor in the time-like region from the process $`e^+e^{}\pi ^+\pi ^{}`$. For roughly twenty years, $`\rho ^0\omega `$ mixing amplitude was assumed constant or momentum independt even if $`\rho `$ and $`\omega `$ have the space-like momenta, far from the on-shell point. Several years ago, this assumption was firstly questioned by Goldman et. al., and the mixing amplitude was found to be significantly momentum dependent within a simple quark loop model. Subsequently, various authors have argued such momentum dependence of the $`\rho ^0\omega `$ mixing amplitude by using various approaches. In particular, the authors of ref. has pointed out that $`\rho ^0\omega `$ mixing amplitude must vanish at $`q^2=0`$(where $`q^2`$ denotes the four-momentum square of the vector mesons) within a broad class of model. This point will be also exmined in ChCQM.
iii) It has been known that $`\omega \pi ^+\pi ^{}`$ decay amplitude receive the contribution from two sources: isospin symmetry breaking due to $`ud`$ quark mass difference and electromagnetic interaction. In I we have shown that VMD in meson physics is natural consequence of the present formlism instead of input. The vector $`e^+e^{}`$ decays are also predicted successfully. Therefore, the dynamics of electromagnetic interactions of mesons has been well established, and the calculation for $`\omega \pi ^+\pi ^{}`$ decay from the transition $`\omega \gamma \rho \pi \pi `$ and “direct” $`\omega \gamma \pi \pi `$ is straightforward. In this the present paper, we will pay our attention to isospin breaking due to $`m_um_d`$. It is another purpose of this paper to determine isospin breaking parameter $`\delta m_qm_dm_u`$ via $`\omega \pi ^+\pi ^{}`$ decay. This parameter is urgently wanted by determination of light quark mass ratios.
The contents of the paper are organized as follows. In sect. 2 we review the basic notations of the chiral constituent quark model with the lowest vector meson resonances. In sect. 3, the tree level effective lagrangian, which including all order contribution of the chiral expansion, is obtained. The pseudoscalar meson one-loop corrections are calculated in sect. 5. In sect. 6, the formulas and numerical results of $`\omega \rho ^0`$ mixing amplitude and $`\omega \pi ^+\pi ^{}`$ are given. The sect. 7 is devoted to a brief summary.
## II Chiral Constituent Quark Model with Vector Meson
The simplest version of chiral quark model which was originated by Weinberg, and developed by Manohar and Georgi provides a QCD-inspired description on the simple constituent quark model. In view of this model, in the energy region between the chiral symmetry spontaneously broken (CSSB) scale and the confinement scale ($`\mathrm{\Lambda }_{QCD}0.20.3GeV`$), the dynamical field degrees of freedom are constituent quarks(quasi-particle of quarks), gluons and Goldstone bosons associated with CSSB(these Goldstone bosons correspond to lowest pseudoscalar octet). In this quasiparticle description, the effective coupling between gluon and quarks is small and the important interaction is the coupling between quarks and Goldstone bosons. In I we have further included the lowest vector meson resonances into this formlism. At chiral limit, this model is parameterized by the following chiral constituent quark lagrangian
$`_\chi `$ $`=`$ $`i\overline{q}(/+/\mathrm{\Gamma }+g__A/\mathrm{\Delta }\gamma _5i/V)qm\overline{q}q+{\displaystyle \frac{F^2}{16}}<_\mu U^\mu U^{}>+{\displaystyle \frac{1}{4}}m_0^2<V_\mu V^\mu >.`$ (1)
Here $`<\mathrm{}>`$ denotes trace in SU(3) flavour space, $`\overline{q}=(\overline{q}_u,\overline{q}_d,\overline{q}_s)`$ are constituent quark fields. $`V_\mu `$ denotes vector meson octet and singlet, or more convenience, due to OZI rule, they are combined into a singlet “nonet” matrix
$$V_\mu (x)=\lambda 𝐕_\mu =\sqrt{2}\left(\begin{array}{ccc}\frac{\rho _\mu ^0}{\sqrt{2}}+\frac{\omega _\mu }{\sqrt{2}}& \rho _\mu ^+& K_\mu ^+\\ \rho _\mu ^{}& \frac{\rho _\mu ^0}{\sqrt{2}}+\frac{\omega _\mu }{\sqrt{2}}& K_\mu ^0\\ K_\mu ^{}& \overline{K}_\mu ^0& \varphi _\mu \end{array}\right).$$
(2)
The $`\mathrm{\Delta }_\mu `$ and $`\mathrm{\Gamma }_\mu `$ are defined as follows,
$`\mathrm{\Delta }_\mu `$ $`=`$ $`{\displaystyle \frac{1}{2}}\{\xi ^{}(_\mu ir_\mu )\xi \xi (_\mu il_\mu )\xi ^{}\},`$ (3)
$`\mathrm{\Gamma }_\mu `$ $`=`$ $`{\displaystyle \frac{1}{2}}\{\xi ^{}(_\mu ir_\mu )\xi +\xi (_\mu il_\mu )\xi ^{}\},`$ (4)
and covariant derivative are defined as follows
$`_\mu U`$ $`=`$ $`_\mu Uir_\mu U+iUl_\mu =2\xi \mathrm{\Delta }_\mu \xi ,`$ (5)
$`_\mu U^{}`$ $`=`$ $`_\mu U^{}il_\mu U^{}+iU^{}r_\mu =2\xi ^{}\mathrm{\Delta }\xi ^{},`$ (6)
where $`l_\mu =v_\mu +a_\mu `$ and $`r_\mu =v_\mu a_\mu `$ are linear combinations of external vector field $`v_\mu `$ and axial-vector field $`a_\mu `$, $`\xi `$ associates with non-linear realization of spontanoeusly broken global chiral symmetry introduced by Weinberg. This realization is obtained by specifying the action of global chiral group $`G=SU(3)_L\times SU(3)_R`$ on element $`\xi (\mathrm{\Phi })`$ of the coset space $`G/SU(3)__V`$:
$$\xi (\mathrm{\Phi })g_R\xi (\mathrm{\Phi })h^{}(\mathrm{\Phi })=h(\mathrm{\Phi })\xi (\mathrm{\Phi })g_L^{},g_L,g_RG,h(\mathrm{\Phi })H=SU(3)__V.$$
(7)
Explicit form of $`\xi (\mathrm{\Phi })`$ is usual taken
$$\xi (\mathrm{\Phi })=\mathrm{exp}\{i\lambda ^a\mathrm{\Phi }^a(x)/2\},U(\mathrm{\Phi })=\xi ^2(\mathrm{\Phi }),$$
(8)
where the Goldstone boson $`\mathrm{\Phi }^a`$ are treated as pseudoscalar meson octet. The compensating $`SU(3)__V`$ transformation $`h(\mathrm{\Phi })`$ defined by eq.( 3) is th wanted ingredent for a non-linear realization of G. In practice, we shall be interested in transformations of $`V_\mu ,\mathrm{\Delta }_\mu ,\mathrm{\Gamma }_\mu `$ and constituent quark fields under $`SU(3)__V`$. The $`q,\overline{q}`$ transform as matter fields of SU(3)$`__V`$,
$$qh(\mathrm{\Phi })q,\overline{q}\overline{q}h^{}(\mathrm{\Phi }).$$
(9)
The vector meson fields transform homogeneously under SU(3)$`__V`$
$$V_\mu h(\mathrm{\Phi })V_\mu h^{}(\mathrm{\Phi }),$$
(10)
which was suggested by Weinberg and developed further by Callan, Coleman et. al.. Since under local G, the expilcit transformations of external vector and axial-vector fields are
$`l_\mu v_\mu a_\mu g_L(x)l_\mu g_L^{}(x)+ig_L(x)_\mu g_L^{}(x),`$ (11)
$`r_\mu v_\mu +a_\mu g_R(x)r_\mu g_R^{}(x)+ig_R(x)_\mu g_R^{}(x),`$ (12)
$`\mathrm{\Delta }_\mu `$ is SU(3)$`__V`$ is invariant field gradients and $`\mathrm{\Gamma }_\mu `$ transforms as field connection of SU(3)$`__V`$
$$\mathrm{\Delta }_\mu h(\mathrm{\Phi })\mathrm{\Delta }_\mu h^{}(\mathrm{\Phi }),\mathrm{\Gamma }_\mu h(\mathrm{\Phi })\mathrm{\Gamma }_\mu h^{}(\mathrm{\Phi })+h(\mathrm{\Phi })_\mu h^{}(\mathrm{\Phi }).$$
(13)
Thus the lagrangian( 1) is invariant under $`G_{\mathrm{global}}\times G_{\mathrm{local}}`$.
The several remarks are need here. 1) Note that there is no kinetic term for vector meson fields in $`_\chi `$. Therefore, in this formlism the vector mesons are treated as composited fields of constituent equarks instead of fundamental fields. The dynamics of vector meson resonances will be generated via loop effects of constituent quarks. 2) Note that there is kinetic term of pseudoscalar mesons in $`_\chi `$. This is different from some other chiral quark models, in which there is no such term. Existing of this kinetic term is consistent with basic assumption of our model, because in this energy region, the dynamical field degrees of freedom are both constituent quarks and Goldstone bosons associated with CSSB. 3) In $`_\chi `$ the parameter $`g_A0.75`$ is determined by $`\beta `$ decay of neutron. It has been pointed out in I that this value has included effects of intermediate axial-vector resonances exchanges at low energy. In addition, the constituent quark mass parameter $`m480`$MeV has been fitted in I via low energy limit of the model. Such large value is requireed by convergence of chiral expansion at vector meson energy scale.
In this paper, we must go beyond chiral limit for obtaining isospin breaking results. The light current quark matrix $`=\mathrm{diag}\{m_u,m_d,m_s\}`$ can be usually included into external scalar fields, i.e., $`\stackrel{~}{\chi }=s+ip`$, where $`s=s_{ext}+`$, $`s_{\mathrm{ext}}`$ and $`p`$ are scalar and pseudoscalar external fields respectively. The chiral transformation for $`\stackrel{~}{\chi }`$ is
$$\stackrel{~}{\chi }g_R\stackrel{~}{\chi }g_L^{}.$$
(14)
Thus together with $`\xi `$ and $`\xi ^{}`$, $`\stackrel{~}{\chi }`$ and $`\stackrel{~}{\chi }^{}`$ can form SU(3)$`__V`$ invariant scalar source $`\xi ^{}\stackrel{~}{\chi }\xi ^{}+\xi \stackrel{~}{\chi }^{}\xi `$ and pseudoscalar source $`(\xi ^{}\stackrel{~}{\chi }\xi ^{}\xi \stackrel{~}{\chi }^{}\xi )\gamma _5`$. Then current quark mass dependent lagrangian is written
$$\frac{1}{2}\overline{q}(\xi ^{}\stackrel{~}{\chi }\xi ^{}+\xi \stackrel{~}{\chi }^{}\xi )q\frac{\kappa }{2}\overline{q}(\xi ^{}\stackrel{~}{\chi }\xi ^{}\xi \stackrel{~}{\chi }^{}\xi )\gamma _5q.$$
(15)
The above lagrangian will return to QCD lagrangian $`\overline{\psi }\psi `$ in absence of pseudoscalar mesons at high energy. So that there is a free parameter $`\kappa `$ which can not be determined by symmetry alonely. From viewpoint of phenomenology, this ambiguity is similar to Kaplan-Manohar ambiguity in ChPT. It will be studied at elsewhere so that we do not further discuss it here. In next section, rigorous calculation will show that our results in this paper is independent of $`\kappa `$.
To conclude this section, the ChCQM lagrangian including the lowest vector meson resonances and light current quark masses is
$`_\chi `$ $`=`$ $`i\overline{q}(/+/\mathrm{\Gamma }+g__A/\mathrm{\Delta }\gamma _5i/V)qm\overline{q}q{\displaystyle \frac{1}{2}}\overline{q}(\xi ^{}\stackrel{~}{\chi }\xi ^{}+\xi \stackrel{~}{\chi }^{}\xi )q{\displaystyle \frac{\kappa }{2}}\overline{q}(\xi ^{}\stackrel{~}{\chi }\xi ^{}\xi \stackrel{~}{\chi }^{}\xi )\gamma _5q`$ (17)
$`+{\displaystyle \frac{F^2}{16}}<_\mu U^\mu U^{}>+{\displaystyle \frac{1}{4}}m_0^2<V_\mu V^\mu >.`$
The effective lagrangian describing interaction of vector meson resonances will be generated via loop effects of constituent quarks.
## III Leading Order Effective Lagrangian in $`N_c^1`$ Expansion
In this section, we will derive relevant effective lagrangian via calculating one-loop diagrams of constituent quarks, which is at the leading order in $`N_c^1`$ expansion. In our calculation, the light current quark masses will be treated as perturbation and be expanded to the leading order. Pseudoscalar mesons which are localed in external line should satisfy soft pion theorem, i.e., $`k^20`$(where $`k^2`$ denotes the four-momentum square of pseudoscalar mesons). However, $`q^2`$ is the four-momentum square of vector mesons, and obviously it is not very small comparing with chiral symmetry spantoneously broken scale. Therefore, the higher order terms of $`q^2`$ in the chiral expansion, have significant contribution which can not be neglected. Or in the other words, the chiral expansion at vector meson energy scale converge slowly. We will rigorously calculate all $`q^2`$-dependent term contribution.
We start with constituent quark lagrangian( 1), and define vector external source $`\overline{V}_\mu ^a(a=0,1,\mathrm{},8)`$, axial-vector external source $`\mathrm{\Delta }_\mu ^a`$, scalar external source $`S^a`$ and pseudoscalar external source $`P^a`$ as follows
$`\overline{V}_\mu ^\alpha `$ $`=`$ $`{\displaystyle \frac{1}{2}}<\lambda ^\alpha (V_\mu +i\mathrm{\Gamma }_\mu )>,\mathrm{\Delta }_\mu ^a={\displaystyle \frac{1}{2}}<\lambda ^a\mathrm{\Delta }_\mu >,`$ (18)
$`S^a`$ $`=`$ $`{\displaystyle \frac{1}{4}}<\lambda ^a(\xi ^{}\stackrel{~}{\chi }\xi ^{}+\xi \stackrel{~}{\chi }^{}\xi )>,P^a={\displaystyle \frac{\kappa }{4}}<\lambda ^a(\xi ^{}\stackrel{~}{\chi }\xi ^{}\xi \stackrel{~}{\chi }^{}\xi )>`$ (19)
(where $`\lambda ^1,\mathrm{},\lambda ^8`$ are SU(3) Gell-Mann matrices and $`\lambda ^0=\sqrt{\frac{2}{3}}`$). Then in lagrangian( 1), the terms associating with constituent quark fields can be rewritten as follow
$`_\chi ^q=\overline{q}(i/m)q+\overline{V}_\mu ^a\overline{q}\lambda ^a\gamma ^\mu q+ig_A\mathrm{\Delta }_\mu ^a\overline{q}\lambda ^a\gamma ^\mu \gamma _5qS^a\overline{q}\lambda ^aqP^a\overline{q}\lambda ^a\gamma _5q.`$ (20)
The effective action describing meson interaction can be obtained via integrating over degrees of freedom of fermions
$$e^{iS_{\mathrm{eff}}}𝒟\overline{q}𝒟qe^{i{\scriptscriptstyle d^4x_\chi (x)}}=<vac,out|in,vac>_{\overline{V},\mathrm{\Delta },S,P},$$
(21)
where $`<vac,out|in,vac>_{\overline{V},\mathrm{\Delta },S,P}`$ is vacuum expectation value in presence external sources. The above path integral can be performed explicitly, and heat kernal method has been used to regulate the result. However, this method is extremely difficult to compute very high order contributions in practice. In I we have provided an equivalent and efficient method to evaluate the effective action via calculating one-loop diagrams of constituent quarks directly. This method can capture all high order contributions of the chiral expansion.
In interaction picture, the equation( 21) is rewritten as follow
$`e^{iS_{\mathrm{eff}}}`$ $`=`$ $`<0|𝒯_qe^{i{\scriptscriptstyle d^4x_\chi ^\mathrm{I}(x)}}|0>`$ (22)
$`=`$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}i{\displaystyle d^4p_1\frac{d^4p_2}{(2\pi )^4}\mathrm{}\frac{d^4p_n}{(2\pi )^4}\stackrel{~}{\mathrm{\Pi }}_n(p_1,\mathrm{},p_n)\delta ^4(p_1p_2\mathrm{}p_n)}`$ (23)
$``$ $`i\mathrm{\Pi }_1(0)+{\displaystyle \underset{n=2}{\overset{\mathrm{}}{}}}i{\displaystyle \frac{d^4p_1}{(2\pi )^4}\mathrm{}\frac{d^4p_{n1}}{(2\pi )^4}\mathrm{\Pi }_n(p_1,\mathrm{},p_{n1})},`$ (24)
where $`𝒯_q`$ is time-order product of constituent quark fields, $`_\chi ^\mathrm{I}`$ is interaction part of lagrangian( 20), $`\stackrel{~}{\mathrm{\Pi }}_n(p_1,\mathrm{},p_n)`$ is one-loop effects of constituent quarks with $`n`$ external sources, $`p_1,p_2,\mathrm{},p_n`$ are four-momentas of $`n`$ external sources respectively and
$$\mathrm{\Pi }_n(p_1,\mathrm{},p_{n1})=d^4p_n\stackrel{~}{\mathrm{\Pi }}_n(p_1,\mathrm{},p_n)\delta ^4(p_1p_2\mathrm{}p_n).$$
(25)
To get rid of all disconnected diagrams, we have
$`S_{\mathrm{eff}}`$ $`=`$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}S_n,`$ (26)
$`S_1`$ $`=`$ $`\mathrm{\Pi }_1(0),`$ (27)
$`S_n`$ $`=`$ $`{\displaystyle \frac{d^4p_1}{(2\pi )^4}\mathrm{}\frac{d^4p_{n1}}{(2\pi )^4}\mathrm{\Pi }_n(p_1,\mathrm{},p_{n1})},(n2).`$ (28)
Hereafter we will call $`S_n`$ as $`n`$-point effective action.
The $`S_1`$ is tadpole-loop contribution of constituent quarks, which is independent of the purpose this paper. The two-point effective action $`S_2`$ has been evaluated in I,
$`S_2`$ $`=`$ $`{\displaystyle \frac{F_0^2}{16}}{\displaystyle d^4x<_\mu U^\mu U^{}>}+{\displaystyle \frac{d^4q}{(2\pi )^4}\mathrm{\Pi }_2^V(q)},`$ (29)
$`\mathrm{\Pi }_2^V(q)`$ $`=`$ $`{\displaystyle \frac{1}{2}}A(q^2)(q^2\delta _{\mu \nu }q_\mu q_\nu )<\overline{V}^\mu (q)\overline{V}^\nu (q)>,`$ (30)
where kinetic term of pseudoscalar mesons has been renormalized, and $`A(q^2)`$ is defined as follow
$`A(q^2)=g^2{\displaystyle \frac{N_c}{\pi ^2}}{\displaystyle _0^1}𝑑tt(1t)\mathrm{ln}(1{\displaystyle \frac{t(1t)q^2}{m^2}}).`$ (31)
Here a universal constant of the model, $`g`$, is defined to absorbe logarithmic divergence from quark loop integral
$`g^2={\displaystyle \frac{8}{3}}{\displaystyle \frac{N_c}{(4\pi )^{D/2}}}({\displaystyle \frac{\mu ^2}{m^2}})^{ϵ/2}\mathrm{\Gamma }(2{\displaystyle \frac{D}{2}}).`$ (32)
In I we have fitted $`g\pi ^1\sqrt{N_c/3}`$ which satisfy the first KSRF sum rule rigorously.
### A Three-point effective action
Up to the leading order of light current quark masses, there are three kinds of three-point effective action. They are made up of by external sources $`\overline{V}\mathrm{\Delta }\mathrm{\Delta }`$, $`\overline{V}\overline{V}S`$ and $`\overline{V}\mathrm{\Delta }P`$ respectively. The effective action with external source $`\overline{V}\mathrm{\Delta }\mathrm{\Delta }`$ has been obtained in I,
$`S_3^{(1)}={\displaystyle \frac{g_A^2}{2}}{\displaystyle d^4x\frac{d^4q}{(2\pi )^4}e^{iqx}B(q^2)q^\mu <\overline{V}^\nu (q)[\mathrm{\Delta }_\mu (x),\mathrm{\Delta }_\nu (x)]>},`$ (33)
where
$`B(q^2)=g^2`$ $`+`$ $`{\displaystyle \frac{N_c}{2\pi ^2}}{\displaystyle _0^1}dt_1t_1{\displaystyle _0^1}dt_2(1t_1t_2)[1+{\displaystyle \frac{m^2}{m^2t_1(1t_1)(1t_2)q^2}}`$ (34)
$`+`$ $`\mathrm{ln}(1{\displaystyle \frac{t_1(1t_1)(1t_2)q^2}{m^2}})].`$ (35)
Let us now calculate three-point effective action with external source $`\overline{V}\overline{V}S`$. Note that since $`\overline{V}_\mu =V_\mu +ie𝒬A_\mu +i\pi _\mu \pi +\mathrm{}`$, for $`\omega \rho ^0`$ mixing or $`\omega \pi \pi `$ coupling vertices, scalar external source $`S`$ will reduce to constant matrix.
$`iS_3^{(2)}`$ $`=`$ $`{\displaystyle \frac{i}{2}}{\displaystyle d^4xd^4yd^4z\overline{V}_\mu ^a(x)\overline{V}_\nu ^b(y)S^c<0|T\{\overline{q}(x)\gamma ^\mu \lambda ^aq(x)\overline{q}(y)\gamma ^\nu \lambda ^bq(y)\overline{q}(z)\lambda ^cq(z)\}|0>}`$ (36)
$`=`$ $`i{\displaystyle \frac{d^4q}{(2\pi )^4}\frac{d^4l}{(2\pi )^4}<\overline{V}_\mu (q)\overline{V}_\nu (q)S>Tr_{(c,L)}\{S_F(l)\gamma ^\nu S_F^2(l+q)\gamma ^\mu \}},`$ (37)
where $`Tr_{(c,L)}`$ denotes trace taking over color and Lorentz space, $`S_F(k)=i(/km+iϵ)^1`$ is propagator of constituent quark fields in momentum space. The direct calculation will give
$`S_3^{(2)}`$ $`=`$ $`{\displaystyle \frac{N_c}{12\pi ^2m}}{\displaystyle d^4x\frac{d^4q}{(2\pi )^4}e^{iqx}h_0(q^2)(q^2\delta _{\mu \nu }q_\mu q_\nu )<\overline{V}^\mu (q)\overline{V}^\nu (x)S(x)>}`$ (38)
$`=`$ $`{\displaystyle \frac{N_c}{24\pi ^2m}}{\displaystyle d^4x\frac{d^4q}{(2\pi )^4}e^{iqx}h_0(q^2)(q^2\delta _{\mu \nu }q_\mu q_\nu )<\overline{V}^\mu (q)\overline{V}^\nu (x)(\xi \stackrel{~}{\chi }^{}\xi +\xi ^{}\stackrel{~}{\chi }\xi ^{})>},`$ (39)
where
$`h_0(q^2)={\displaystyle _0^1}𝑑t{\displaystyle \frac{6t(1t)}{1t(1t)q^2/m^2}}.`$ (40)
Next, we calculate three-point effective action with external source $`\overline{V}\mathrm{\Delta }P`$.
$`iS_3^{(3)}`$ $`=`$ $`g_A{\displaystyle d^4xd^4yd^4z<0|T\{\overline{q}(x)\gamma ^\mu \lambda ^aq(x)\overline{q}(y)\gamma ^\nu \gamma _5\lambda ^bq(y)\overline{q}(z)\gamma _5\lambda ^cq(z)\}|0>\overline{V}_\mu ^a(x)\mathrm{\Delta }_\nu ^b(y)P^c(z)}`$ (41)
$`=`$ $`g_A{\displaystyle }{\displaystyle \frac{d^4q}{(2\pi )^4}}{\displaystyle \frac{d^4k}{(2\pi )^4}}{\displaystyle \frac{d^4l}{(2\pi )^4}}\{<\overline{V}_\mu (q)\mathrm{\Delta }_\nu (k)P(qk)>Tr_{(c,L)}[S_F(l)\gamma ^\nu \gamma _5S_F(lk)\gamma _5S_F(l+q)\gamma ^\mu ]`$ (43)
$`+<\overline{V}_\mu (q)P(qk)\mathrm{\Delta }_\nu (k)>Tr_{(c,L)}[S_F(l)\gamma ^\mu S_F(lq)\gamma _5S_F(l+k)\gamma ^\mu \gamma _5]\}.`$
Due to $`\mathrm{\Delta }_\mu _\mu \pi +\mathrm{}`$ and $`P\pi +\mathrm{}`$, for purpose of this paper, the soft pion theroem tell us $`k^20`$ and $`(k+q)^20`$. In addition, we can find that $`k_\mu \mathrm{\Delta }^\mu (k)k^2\pi (k)0`$. Then performing the loop-integral in eq.( 41), and employing the above discussion in our calculation, we obtain
$`S_3^{(3)}`$ $`=`$ $`{\displaystyle \frac{N_cm}{4\pi ^2}}g_A{\displaystyle }{\displaystyle \frac{d^4q}{(2\pi )^4}}{\displaystyle \frac{d^4k}{(2\pi )^4}}\{\alpha _1(q^2)<[\overline{V}_\mu (q),\mathrm{\Delta }^\mu (k)]P(kq)>`$ (45)
$`+\alpha _2(q^2)(q^2\delta _{\mu \nu }q_\mu q_\nu )<[\overline{V}^\nu (q),\mathrm{\Delta }^\nu (k)]P(kq)>\},`$
where
$`\alpha _1(q^2)`$ $`=`$ $`({\displaystyle \frac{4\pi \mu ^2}{m^2}})^{ϵ/2}\mathrm{\Gamma }(2{\displaystyle \frac{D}{2}}){\displaystyle _0^1}dt_1{\displaystyle _0^1}dt_2\{2t_1\mathrm{ln}(1{\displaystyle \frac{t_1(1t_1)(1t_2)q^2}{m^2}})`$ (47)
$`+{\displaystyle \frac{t_1t_2(1t_1)q^2}{m^2t_1(1t_1)(1t_2)q^2}}\},`$
$`\alpha _2(q^2)`$ $`=`$ $`{\displaystyle _0^1}𝑑t_1{\displaystyle _0^1}𝑑t_2{\displaystyle \frac{2t_1(1t_1)}{m^2t_1(1t_1)(1t_2)q^2}}.`$ (48)
The third three-point effective action $`S_3^{(3)}`$ is $`O(m_q)`$ and free parameter $`\kappa `$-dependent. However, if $`\overline{V}_\mu =\omega _\mu `$, $`S_3^{(3)}`$ vanish, and if $`\overline{V}_\mu =\rho _\mu ^0\lambda ^3`$, $`S_3^{(3)}`$ provide an isospin conservation $`\rho \pi \pi `$ vertex which is order $`m_u+m_d`$ and much smaller that leading order vertex. Thus the contribution from $`S_3^{(3)}`$ will be omitted in this paper.
### B Four-point effective action
There is only one four-point effective action relating to $`\omega \pi ^+\pi ^{}`$ decay. It is made up of by four external source $`\overline{V}\mathrm{\Delta }\mathrm{\Delta }S`$. As shown in the above subsection, here $`S`$ reduces to a constant matrix.
$`iS_4`$ $`=`$ $`{\displaystyle \frac{g_A^2}{2}}{\displaystyle d^4xd^4yd^4zd^4w\overline{V}_\mu ^a(x)\mathrm{\Delta }_\nu ^b(y)\mathrm{\Delta }_\sigma ^c(z)S^d(w)}`$ (50)
$`\times <0|T\{\overline{q}(x)\gamma ^\mu \lambda ^aq(x)\overline{q}(y)\gamma ^\nu \gamma _5\lambda ^bq(y)\overline{q}(z)\gamma ^\sigma \gamma _5\lambda ^cq(z)\overline{q}(w)\lambda ^dq(w)\}|0>`$
$`=`$ $`g_A^2{\displaystyle \frac{d^4q}{(2\pi )^4}\frac{d^4k}{(2\pi )^4}\frac{d^4l}{(2\pi )^4}}`$ (54)
$`\times \{<\overline{V}_\mu (q)\mathrm{\Delta }_\nu (k)\mathrm{\Delta }_\sigma (kq)S>Tr_{(c.L)}[S_F^2(l)\gamma ^\mu S_F(lq)\gamma ^\nu \gamma _5S_F(lqk)\gamma ^\sigma \gamma _5]`$
$`+<\overline{V}_\mu (q)\mathrm{\Delta }_\nu (k)S\mathrm{\Delta }_\sigma (kq)>Tr_{(c.L)}[S_F(l+q+k)\gamma ^\mu S_F(l+k)\gamma ^\nu \gamma _5S_F^2(l)\gamma ^\sigma \gamma _5]`$
$`+<\overline{V}_\mu (q)S\mathrm{\Delta }_\nu (k)\mathrm{\Delta }_\sigma (kq)>Tr_{(c.L)}[S_F(l+q)\gamma ^\mu S_F^2(l)\gamma ^\nu \gamma _5S_F(lk)\gamma ^\sigma \gamma _5]\}.`$
To perform integral of four-momenta $`l`$ in the above equation and employ indentities in Appendix to simplify result, we can obtain
$`S_4`$ $`=`$ $`{\displaystyle \frac{N_c}{8\pi ^2m}}g_A^2{\displaystyle d^4x\frac{d^4q}{(2\pi )^4}e^{iqx}(\delta _{\mu \nu }q_\sigma \delta _{\mu \sigma }q_\nu )}`$ (56)
$`\times \{h_1(q^2)<\{\overline{V}^\mu (q),S(x)\}\mathrm{\Delta }^\nu (x)\mathrm{\Delta }^\sigma (x)>+h_2(q^2)<\overline{V}^\mu (q)\mathrm{\Delta }^\nu (x)S(x)\mathrm{\Delta }^\sigma (x)>\},`$
where
$`h_1(q^2)`$ $`=`$ $`{\displaystyle _0^1}𝑑t_1t_1^2{\displaystyle _0^1}𝑑t_2(1t_2){\displaystyle \frac{32t_1^2t_2(1+2t_1)(1t_2)q^2/m^2}{[1t_1^2t_2(1t_2)q^2/m^2]^2}},`$ (57)
$`h_2(q^2)`$ $`=`$ $`{\displaystyle _0^1}𝑑t_1t_1^2{\displaystyle _0^1}𝑑t_2(1t_2){\displaystyle \frac{4(1t_1)[34t_1^2t_2(1t_2)q^2/m^2]}{[1t_1^2t_2(1t_2)q^2/m^2]^2}}.`$ (58)
### C Relevant effective vertices at tree level
In the following we will give all relevant effective vertice at tree level. Since one-loop correction of pseudoscalar mesons will be calculated, we also need to include four-pseudoscalar meson vertices, which will be derived from $`O(p^4)`$ effective lagrangian of this formlism. The effective vertices involving electromagnetic interaction have been calculated up to one-loop level in I. We will quote them in sect. 5 directly. All effective vertices can be divided into two part: one is isospin conservation and anthor is isospin broken. In addition, we should point out that, so far, all meson fields are still non-physical. The physical meson fields can be obtained via the following field rescaling which make kinetic terms of pseudoscalar mesons and vector mesons into standard form
$`\rho _\mu ^0{\displaystyle \frac{1}{g}}\rho _\mu ^0,\omega _\mu {\displaystyle \frac{1}{g}}\omega _\mu ,`$ (59)
$`\pi {\displaystyle \frac{2}{f_\pi }}\pi ,K{\displaystyle \frac{2}{f_\pi }}K.`$ (60)
Since in this paper $`K`$-mesons only appear as intermediate states in one-loop diagrams, for sake of convenience, we neglect the difference between $`f_\pi `$ and $`f_K`$(since the results yielded by this difference are twofold suppressed by light current quark mass expansion and $`N_c^1`$ expansion).
The $`\rho ^0\omega `$ mixing vertex, which breaks isospin symmetry, is included in eq.( 38)
$`_{\omega \rho }`$ $`=`$ $`{\displaystyle \frac{N_c}{12\pi ^2g^2m}}{\displaystyle \frac{d^4q}{(2\pi )^4}e^{iqx}h_0(q^2)(q^2\delta _{\mu \nu }q_\mu q_\nu )\omega ^\mu (q)\rho ^{0\nu }(x)<\lambda ^3(\xi \xi +\xi ^{}\xi ^{})>}`$ (61)
$`=`$ $`{\displaystyle \frac{N_c}{6\pi ^2g^2}}{\displaystyle \frac{m_um_d}{m}}{\displaystyle \frac{d^4q}{(2\pi )^4}e^{iqx}h_0(q^2)(q^2\delta _{\mu \nu }q_\mu q_\nu )\omega ^\mu (q)\rho ^{0\nu }(x)}+\mathrm{}.`$ (62)
The isospin symmetry unbroken vector$`\phi \phi `$ vertex is included in eqs.( 29) and ( 33)
$`_{V\phi \phi }^{(\mathrm{\Delta }I=0)}=2{\displaystyle \frac{d^4q}{(2\pi )^4}e^{iqx}b(q^2)q_\mu <V_\nu (q)[\mathrm{\Delta }^\mu ,\mathrm{\Delta }^\nu ]>},`$ (63)
where
$$b(q^2)=\frac{1}{gf_\pi ^2}[A(q^2)+g_A^2B(q^2)].$$
(64)
The isospin symmetry broken vector$`\phi \phi `$ vertex is include in eqs.( 38) and ( 56)
$`_{V\phi \phi }^{(\mathrm{\Delta }I=1)}`$ $`=`$ $`{\displaystyle \frac{N_c}{12\pi ^2mg}}{\displaystyle \frac{d^4q}{(2\pi )^4}e^{iqx}(\delta _{\mu \nu }q_\sigma \delta _{\mu \sigma }q_\nu )}`$ (67)
$`\times \{[h_0(q^2)+{\displaystyle \frac{3}{4}}g_A^2h_1(q^2)]<\{\overline{V}^\mu (q),\xi \xi +\xi ^{}\xi ^{}\}\mathrm{\Delta }^\nu (x)\mathrm{\Delta }^\sigma (x)>`$
$`{\displaystyle \frac{3}{4}}g_A^2h_2(q^2)<\overline{V}^\mu (q)\mathrm{\Delta }^\nu (x)(\xi \xi +\xi ^{}\xi ^{})\mathrm{\Delta }^\sigma (x)>\}.`$
In particular, define
$$s(q^2)=\frac{4}{gf_\pi ^2}[h_0(q^2)+\frac{3}{4}g_A^2(h_1(q^2)\frac{h_2(q^2)}{2})],$$
(68)
we have
$`_{\rho ^0\pi \pi }`$ $`=`$ $`i{\displaystyle \frac{d^4q}{(2\pi )^4}e^{iqx}b(q^2)(q^2\delta _{\mu \nu }q_\mu q_\nu )\rho ^{0\mu }(q)[\pi ^+(x)^\nu \pi ^{}(x)^\nu \pi ^+(x)\pi ^{}(x)]}`$ (69)
$`_{\omega \pi \pi }`$ $`=`$ $`{\displaystyle \frac{iN_c}{12\pi ^2}}{\displaystyle \frac{m_um_d}{m}}{\displaystyle \frac{d^4q}{(2\pi )^4}e^{iqx}s(q^2)(q^2\delta _{\mu \nu }q_\mu q_\nu )\omega ^\mu (q)}`$ (71)
$`\times [\pi ^+(x)^\nu \pi ^{}(x)^\nu \pi ^+(x)\pi ^{}(x)]`$
$`_{\rho ^0KK}^{(\mathrm{\Delta }I=0)}`$ $`=`$ $`{\displaystyle \frac{i}{2}}{\displaystyle \frac{d^4q}{(2\pi )^4}e^{iqx}b(q^2)(q^2\delta _{\mu \nu }q_\mu q_\nu )\rho ^{0\mu }(q)}`$ (72)
$`\times `$ $`\{[K^+(x)^\nu K^{}(x)^\nu K^+(x)K^{}(x)][K^0(x)^\nu \overline{K}^0(x)^\nu K^0(x)\overline{K}^0(x)]\}`$ (73)
$`_{\rho ^0KK}^{(\mathrm{\Delta }I=1)}`$ $`=`$ $`{\displaystyle \frac{iN_c}{24\pi ^2}}{\displaystyle \frac{m_um_d}{m}}{\displaystyle \frac{d^4q}{(2\pi )^4}e^{iqx}s(q^2)(q^2\delta _{\mu \nu }q_\mu q_\nu )\rho ^{0\mu }(q)}`$ (74)
$`\times `$ $`\{[K^+(x)^\nu K^{}(x)^\nu K^+(x)K^{}(x)]+[K^0(x)^\nu \overline{K}^0(x)^\nu K^0(x)\overline{K}^0(x)]\}`$ (75)
$`_{\omega KK}^{(\mathrm{\Delta }I=0)}`$ $`=`$ $`{\displaystyle \frac{i}{2}}{\displaystyle \frac{d^4q}{(2\pi )^4}e^{iqx}b(q^2)(q^2\delta _{\mu \nu }q_\mu q_\nu )\omega ^\mu (q)}`$ (76)
$`\times `$ $`\{[K^+(x)^\nu K^{}(x)^\nu K^+(x)K^{}(x)]+[K^0(x)^\nu \overline{K}^0(x)^\nu K^0(x)\overline{K}^0(x)]\}`$ (77)
$`_{\omega KK}^{(\mathrm{\Delta }I=1)}`$ $`=`$ $`{\displaystyle \frac{iN_c}{24\pi ^2}}{\displaystyle \frac{m_um_d}{m}}{\displaystyle \frac{d^4q}{(2\pi )^4}e^{iqx}s(q^2)(q^2\delta _{\mu \nu }q_\mu q_\nu )\omega ^\mu (q)}`$ (78)
$`\times `$ $`\{[K^+(x)^\nu K^{}(x)^\nu K^+(x)K^{}(x)][K^0(x)^\nu \overline{K}^0(x)^\nu K^0(x)\overline{K}^0(x)]\}`$ (79)
Up to $`O(p^4)`$, the tree level four-pseudoscalar meson effective lagrangian has been derived in I. The isospin symmetry unbroken four-pseudoscalar vertex is included in the following lagrangian
$`_{4\phi }^{(\mathrm{\Delta }I=0)}`$ $`=`$ $`{\displaystyle \frac{f_\pi ^2}{16}}<_\mu U^\mu U^{}>+{\displaystyle \frac{N_c}{12(4\pi )^2}}<_\mu U_\nu U^{}_\mu U_\nu U^{}>`$ (81)
$`{\displaystyle \frac{N_c}{12(4\pi )^2}}(1g_A^4)<_\mu U^\mu U^{}_\nu U^\nu U^{}>,`$
where we have used $`g^2=N_c/(3\pi ^2)`$. The isospin symmetry broken four-pseudoscalar vertex is proportional to $`m_dm_u`$, which is included in the following lagrangian
$`_{4\phi }^{(\mathrm{\Delta }I=1)}={\displaystyle \frac{f_\pi ^2}{8}}B_0<(U+U^{})>+{\displaystyle \frac{N_cm}{(4\pi )^2}}g_A^2<_\mu U^\mu U^{}(U^{}+U)>.`$ (82)
I can be found that the eqs.( 61)-( 82) are free parameter $`\kappa `$ independent. Moreover, we can see that every vector$`\phi \phi `$ in eq.( 69) includes an antisymmetry factor $`(q^2\delta _{\mu \nu }q_\mu q_\nu )`$(where $`q_\mu `$ denotes four-momenta of vector mesons). Thus the first term of eq.( 82) does not contribute to $`\omega \pi ^+\pi ^{}`$ decay via pseudoscalar meson loops. This antisymmetry factor also constrains that the vertices with one of factors $`K\overline{K}`$, $`_\mu K^\mu \overline{K}`$ and $`K^2\overline{K}`$ do not contribute to $`\omega \pi ^+\pi ^{}`$ decay via pseudoscalar meson loops. Then the relevant four-pseudoscalar vertices can explicitly read as follows,
$`_{4\pi }`$ $`=`$ $`{\displaystyle \frac{2}{f_\pi ^2}}(\pi ^+_\mu \pi ^{})(\pi ^+^\mu \pi ^{})+{\displaystyle \frac{2N_c}{3\pi ^2f_\pi ^4}}_\mu \pi ^+_\nu \pi ^{}[^\mu \pi ^+^\nu \pi ^{}(1g_A^4)^\mu \pi ^{}^\nu \pi ^+],`$ (83)
$`_{KK\pi \pi }^{(\mathrm{\Delta }I=0)}`$ $`=`$ $`{\displaystyle \frac{4}{f_\pi ^2}}(K^+_\mu K^{}+_\mu K^0\overline{K}^0)\pi ^+^\mu \pi ^{}`$ (85)
$`+{\displaystyle \frac{2N_c}{3\pi ^2f_\pi ^4}}(_\mu K^+_\nu K^{}+_\nu K^0_\mu \overline{K}^0)[2^\mu \pi ^+^\nu \pi ^{}(1g_A^4)^\mu \pi ^{}^\nu \pi ^+],`$
$`_{KK\pi \pi }^{(\mathrm{\Delta }I=1)}`$ $`=`$ $`{\displaystyle \frac{16N_c}{\pi ^2f_\pi ^4}}g_A^2m(m_uK^+_\mu K^{}m_dK^0_\mu \overline{K}^0)\pi ^+^\mu \pi ^{}.`$ (86)
## IV One-loop Corrections of Pseudoscalar Mesons
In this section we calculate one-loop correction of mesons. Because of $`m__V^2>m__K^2>>m_\pi ^2`$, it can be expected that the dominant contribution comes from one-loop digrams of pseudoscalar mesons. In addition, we can treat pion as massless particle but must take $`m__K^20`$. This difference is very important, since $`\pi `$-loop yields imaginary part of $`𝒯`$-matrix but $`K`$-loop does not at $`m_\omega `$ scale. In our calculation, the mass difference between $`K^\pm `$ and $`K^0`$ is also neglected.
There are three kinds of one-loop diagrams correcting to “direct” $`\omega \pi \pi `$ couping and $`\omega \rho ^0`$ mixing(fig. 1 and fig. 2).
It must be pointed out that, in $`𝒯`$-matrices yielded by figure 1-(b) and figure 2-(b), the contribution of imaginary part is dominant. We have shown in I that it can not ensure unitarity of $`S`$-matrix if we only calculate figure 1-(b) and figure 2-(b). The unitarity can be ensured through summing over all diagrams in chain approximation(fig. 3 and fig. 4).
### A Tadpole diagram
Since pion is treated as massless particle, the nonzero tadpole diagram contribution is yielded by $`K`$ or $`\eta _8`$ mesons(fig. 1-(a) and fig. 2-(a)). For sake of convenience, here we assume $`m_{\eta _8}=m__K`$.
Correction to $`\omega \rho ^0`$ mixing
The tree level $`\omega \rho ^0\phi \phi `$ is contained in eq.(35),
$`_{\omega \rho ^0\phi \phi }={\displaystyle \frac{N_c}{12\pi ^2g^2f_\pi ^2m}}{\displaystyle \frac{d^4q}{(2\pi )^4}e^{iqx}h_0(q^2)(q^2\delta _{\mu \nu }q_\mu q_\nu )\omega ^\mu (q)\rho ^{0\nu }(x)<\lambda ^3(\phi ^2+\phi ^2+\phi \phi )>}.`$ (87)
In momentum space, the calculation on fig. 2-(a) is straighforward
$`_{\omega \rho ^0}^{(tad)}`$ $`=`$ $`{\displaystyle \frac{N_c}{12\pi ^2g^2f_\pi ^2m}}{\displaystyle \frac{d^4q}{(2\pi )^4}e^{iqx}h_0(q^2)(q^2\delta _{\mu \nu }q_\mu q_\nu )\omega ^\mu (q)\rho ^{0\nu }(x)}`$ (89)
$`\times {\displaystyle }{\displaystyle \frac{d^4k}{(2\pi )^4}}{\displaystyle \frac{i}{k^2m__K^2+iϵ}}{\displaystyle \underset{a=4}{\overset{8}{}}}<\lambda ^3(\lambda ^a\lambda ^a+\lambda ^a\lambda ^a+\lambda ^a\lambda ^a)>.`$
The generators $`\lambda ^a`$ of SU(N) obey the completeness relations
$`{\displaystyle \underset{a=1}{\overset{N^21}{}}}<\lambda ^aA\lambda ^aB>`$ $`=`$ $`{\displaystyle \frac{2}{N}}<AB>+2<A><B>,`$ (90)
$`{\displaystyle \underset{a=1}{\overset{N^21}{}}}<\lambda ^aA><\lambda ^aB>`$ $`=`$ $`\mathrm{\hspace{0.33em}2}<AB>{\displaystyle \frac{2}{N}}<A><B>.`$ (91)
Then we have
$`{\displaystyle \underset{a=4}{\overset{8}{}}}<\lambda ^3(\lambda ^a\lambda ^a+\lambda ^a\lambda ^a+\lambda ^a\lambda ^a)>={\displaystyle \frac{16}{3}}<\lambda ^3>={\displaystyle \frac{16}{3}}(m_um_d).`$ (92)
Substituting eq.( 92) into eq.( 89) and performing loop integral, we obtain
$`_{\omega \rho ^0}^{(tad)}`$ $`=`$ $`{\displaystyle \frac{4}{3}}\zeta {\displaystyle \frac{N_c}{6\pi ^2g^2}}{\displaystyle \frac{m_um_d}{m}}{\displaystyle \frac{d^4q}{(2\pi )^4}e^{iqx}h_0(q^2)(q^2\delta _{\mu \nu }q_\mu q_\nu )\omega ^\mu (q)\rho ^{0\nu }(x)},`$ (93)
where
$$\zeta =\frac{2\lambda }{(4\pi )^2}\frac{m__K^2}{f_\pi ^2},\lambda =(\frac{4\pi \mu ^2}{m__K^2})^{ϵ/2}\mathrm{\Gamma }(1\frac{D}{2}).$$
(94)
Here we define a constant $`\lambda `$ to absorbe quadratic divergence from meson loop integral. Its value has been determined as $`\lambda =2/3`$ in I by OZI rule.
Correction to “direct” $`\omega \pi \pi `$ mixing
The isospin symmetry broken $`\omega `$-4$`\phi `$ couping vertex is included in eq.( 67). Expanding eq.( 67) to contain four pseudoscalar meson fields, we can obtain
$`_{\omega \pi \pi }^{(tad)}`$ $`=`$ $`{\displaystyle \frac{N_c}{12\pi ^2f_\pi ^2m}}{\displaystyle \frac{d^4q}{(2\pi )^4}e^{iqx}s(q^2)(\delta _{\mu \nu }q_\sigma \delta _{\mu \sigma }q_\nu )\omega ^\mu (q)\frac{d^4k}{(2\pi )^4}\frac{i}{k^2m__K^2+iϵ}}`$ (97)
$`\times {\displaystyle \underset{a=4}{\overset{8}{}}}\{<I_2(^\nu \pi \lambda ^a[\lambda ^a,^\sigma \pi ]+\lambda ^a[\lambda ^a,^\nu \pi ]^\sigma \pi )>`$
$`+{\displaystyle \frac{1}{2}}<I_2^\nu \pi ^\sigma \pi (\lambda ^a\lambda ^a+\lambda ^a\lambda ^a+\lambda ^a\lambda ^a)>\}`$
$`=`$ $`{\displaystyle \frac{10}{3}}\zeta {\displaystyle \frac{iN_c}{12\pi ^2}}{\displaystyle \frac{m_um_d}{m}}{\displaystyle \frac{d^4q}{(2\pi )^4}e^{iqx}s(q^2)(q^2\delta _{\mu \nu }q_\mu q_\nu )\omega ^\mu (q)}`$ (99)
$`\times [\pi ^+(x)^\nu \pi ^{}(x)^\nu \pi ^+(x)\pi ^{}(x)],`$
where $`I_2=\mathrm{diag}\{1,1,0\}`$, $`s(q^2)`$, to see eq.( 68) and eq.( 90) has been used.
### B $`K`$-loop contribution
Here $`K`$-loop denotes that one-loop diagrams in fig.1-(c) and fig.2-(c). In this subsection, $`T`$ will denote time-order product of $`K`$-meson field.
Correction to $`\omega \rho ^0`$ mixing
The $`\omega \rho ^0`$ effective action yielded by $`K`$-loop is follow
$`iS_{\omega \rho }^{(Kloop)}`$ $`=`$ $`{\displaystyle d^4xd^4y<0|T\{_{\rho KK}^{(\mathrm{\Delta }I=0)}(x)_{\omega KK}^{(\mathrm{\Delta }I=1)}(y)+_{\rho KK}^{(\mathrm{\Delta }I=1)}(x)_{\omega KK}^{(\mathrm{\Delta }I=0)}(y)\}|0>}`$ (100)
$`=`$ $`{\displaystyle \frac{N_c}{3\pi ^2}}{\displaystyle \frac{m_um_d}{m}}{\displaystyle \frac{d^4q}{(2\pi )^4}\frac{d^4l}{(2\pi )^4}b(q^2)s(q^2)(q^2\delta _{\mu \nu }q_\mu q_\nu )(q^2\delta _{\alpha \beta }q_\alpha q_\beta )}`$ (102)
$`\times \omega ^\mu (q)\rho ^{0\alpha }(q)(l+q)^\nu l^\beta \mathrm{\Delta }_K(l)\mathrm{\Delta }_K(l+q),`$
where $`\mathrm{\Delta }_K(l)=i(l^2m__K^2+iϵ)^1`$ is propagator of $`K`$-meson. Integrating over $`l_\mu `$ in the above equation and defining
$`\mathrm{\Sigma }_K(q^2)={\displaystyle \frac{1}{(4\pi )^2}}\{\lambda (m__K^2{\displaystyle \frac{q^2}{6}})+{\displaystyle _0^1}𝑑t[m__K^2t(1t)q^2]\mathrm{ln}(1{\displaystyle \frac{t(1t)q^2}{m__K^2}})\},`$ (103)
we have
$`S_{\omega \rho }^{(Kloop)}={\displaystyle \frac{N_c}{6\pi ^2}}{\displaystyle \frac{m_um_d}{m}}{\displaystyle \frac{d^4q}{(2\pi )^4}b(q^2)s(q^2)\mathrm{\Sigma }_K(q^2)q^2(q^2\delta _{\mu \nu }q_\mu q_\nu )\omega ^\mu (q)\rho ^{0\nu }(q)}.`$ (104)
The corresponding effective lagrangian read
$`_{\omega \rho }^{(Kloop)}={\displaystyle \frac{N_c}{6\pi ^2}}{\displaystyle \frac{m_um_d}{m}}{\displaystyle \frac{d^4q}{(2\pi )^4}e^{iqx}b(q^2)s(q^2)\mathrm{\Sigma }_K(q^2)q^2(q^2\delta _{\mu \nu }q_\mu q_\nu )\omega ^\mu (q)\rho ^{0\nu }(x)}.`$ (105)
Correction to “direct” $`\omega \pi \pi `$ couping
The “direct” $`\omega \pi \pi `$ couping effective action yielded by $`K`$-loop can be evaluated as follow
$`iS_{\omega \pi \pi }^{(Kloop)}`$ $`=`$ $`{\displaystyle d^4xd^4y<0|T\{_{\rho KK}^{(\mathrm{\Delta }I=0)}(x)_{\pi \pi KK}^{(\mathrm{\Delta }I=1)}(y)+_{\rho KK}^{(\mathrm{\Delta }I=1)}(x)_{\pi \pi KK}^{(\mathrm{\Delta }I=0)}(y)\}|0>}`$ (106)
$`=`$ $`{\displaystyle \frac{N_c}{\pi ^2f_\pi ^2}}{\displaystyle \frac{m_um_d}{m}}{\displaystyle \frac{d^4q}{(2\pi )^4}\frac{d^4k}{(2\pi )^4}\frac{d^4l}{(2\pi )^4}[16m^2f_\pi ^2b(q^2)\frac{2}{3}s(q^2)\frac{(3g_A^4)N_c}{18\pi ^2f_\pi ^2}q^2s(q^2)]}`$ (108)
$`\times (q^2\delta _{\mu \nu }q_\mu q_\nu )(kl)l^\nu \omega ^\mu (q)\pi ^+(qk)\pi ^{}(k)\mathrm{\Delta }_K(l)\mathrm{\Delta }_K(l+q)`$
$``$ $`{\displaystyle \frac{N_c}{\pi ^2f_\pi ^2}}{\displaystyle \frac{m_um_d}{m}}i{\displaystyle \frac{d^4q}{(2\pi )^4}\frac{d^4k}{(2\pi )^4}[8m^2f_\pi ^2b(q^2)\frac{1}{3}s(q^2)(1+\frac{q^2N_c}{4\pi ^2f_\pi ^2})]\mathrm{\Sigma }_K(q^2)}`$ (110)
$`\times (q^2\delta _{\mu \nu }q_\mu q_\nu )k^\nu \omega ^\mu (q)\pi ^+(qk)\pi ^{}(k),`$
where we have taken soft pion limit and $`3g_A^43`$ due to $`g_A^40.3`$. The corresponding effective lagrangian reads
$`_{\omega \pi \pi }^{(Kloop)}`$ $`=`$ $`{\displaystyle \frac{iN_c}{2\pi ^2f_\pi ^2}}{\displaystyle \frac{m_um_d}{m}}{\displaystyle \frac{d^4q}{(2\pi )^4}e^{iqx}[8m^2f_\pi ^2b(q^2)\frac{1}{3}s(q^2)(1+\frac{q^2N_c}{4\pi ^2f_\pi ^2})]\mathrm{\Sigma }_K(q^2)}`$ (112)
$`\times (q^2\delta _{\mu \nu }q_\mu q_\nu )\omega ^\mu (q)[\pi ^+(x)^\nu \pi ^{}(x)^\nu \pi ^+(x)\pi ^{}(x)].`$
### C Chain contribution of $`\pi `$-loop
Finally, we calculate chain approximation corrections of $`\pi `$-loop in fig. 3 and fig. 4.
Correction to “direct” $`\omega \pi \pi `$ coupling
The effective action yielded by $`\pi `$-loop in fig. 1-(b) is evaluated as follow
$`iS_{\omega \pi \pi }^{(\pi loop)}`$ $`=`$ $`{\displaystyle d^4xd^4y<0|T\{_{\omega \pi \pi }(x)_{4\pi }(y)\}|0>}`$ (113)
$`=`$ $`{\displaystyle \frac{4N_c}{3\pi ^2f_\pi ^2}}{\displaystyle \frac{m_um_d}{m}}{\displaystyle \frac{d^4q}{(2\pi )^4}\frac{d^4k}{(2\pi )^4}\frac{d^4l}{(2\pi )^4}s(q^2)(1+\frac{q^2N_c}{4\pi ^2f_\pi ^2})(q^2\delta _{\mu \nu }q_\mu q_\nu )}`$ (115)
$`\times (lk)l^\nu \omega ^\mu (q)\pi ^+(qk)\pi ^{}(k)\mathrm{\Delta }_\pi (l)\mathrm{\Delta }_\pi (l+k),`$
where we have employed soft pion limit and $`g_A^43`$, and $`\mathrm{\Delta }_\pi (l)=i(l^2+iϵ)^1`$ is propagator of pion. Integrating over $`l`$ in the above equation and defining
$`\mathrm{\Sigma }_\pi (q^2)={\displaystyle \frac{q^2}{(4\pi )^2}}\{{\displaystyle \frac{\lambda }{6}}+{\displaystyle _0^1}𝑑tt(1t)\mathrm{ln}{\displaystyle \frac{t(1t)q^2}{m__K^2}}+{\displaystyle \frac{i}{6}}Arg(1)\theta (q^24m_\pi ^2)\},`$ (116)
we have
$`S_{\omega \pi \pi }^{(\pi loop)}`$ $`=`$ $`{\displaystyle \frac{2N_c}{3\pi ^2f_\pi ^2}}{\displaystyle \frac{m_um_d}{m}}{\displaystyle \frac{d^4q}{(2\pi )^4}\frac{d^4k}{(2\pi )^4}s(q^2)(1+\frac{q^2N_c}{4\pi ^2f_\pi ^2})\mathrm{\Sigma }_\pi (q^2)}`$ (118)
$`\times (q^2\delta _{\mu \nu }q_\mu q_\nu )k^\nu \omega ^\mu (q)\pi ^+(qk)\pi ^{}(k).`$
In eq.( 116), $`Arg(1)=\pi `$ has been fitted in I due to requirement of unitarity. The corresponding effective lagrangian reads
$`_{\omega \pi \pi }^{(\pi loop)}`$ $`=`$ $`{\displaystyle \frac{iN_c}{3\pi ^2f_\pi ^2}}{\displaystyle \frac{m_um_d}{m}}{\displaystyle \frac{d^4q}{(2\pi )^4}e^{iqx}s(q^2)(1+\frac{q^2N_c}{4\pi ^2f_\pi ^2})\mathrm{\Sigma }_\pi (q^2)}`$ (120)
$`\times (q^2\delta _{\mu \nu }q_\mu q_\nu )\omega ^\mu (q)[\pi ^+(x)^\nu \pi ^{}(x)^\nu \pi ^+(x)\pi ^{}(x)].`$
Comparing eq.( 120) and tree level vertex( 69), we can find that every one-loop in fig.3 contributes a factor
$`\mathrm{\Xi }(q^2)=4f_\pi ^2(1+{\displaystyle \frac{q^2N_c}{4\pi ^2f_\pi ^2}})\mathrm{\Sigma }_\pi (q^2).`$ (121)
Thus summing over all diagrams in fig. 3, we obtain
$`_{\omega \pi \pi }^{}={\displaystyle \frac{iN_c}{12\pi ^2}}{\displaystyle \frac{m_um_d}{m}}{\displaystyle \frac{d^4q}{(2\pi )^4}e^{iqx}\frac{s(q^2)}{1+\mathrm{\Xi }(q^2)}(q^2\delta _{\mu \nu }q_\mu q_\nu )\omega ^\mu (q)[\pi ^+(x)^\nu \pi ^{}(x)^\nu \pi ^+(x)\pi ^{}(x)]}.`$ (122)
Correction to $`\omega \rho ^0`$ mixing
If in fig. 2-(b), tree level vertex $`_{\omega \pi \pi }`$ is replaced by $`_{\omega \pi \pi }^{}`$ which contains all diagram contribution in fig. 3, then summing tree diagram and fig. 2-(b) is just chain approximation correction to $`\omega \rho ^0`$ mixing. The effective action yielded by $`\pi `$-loop in fig. 2-(b) is evaluated as follow
$`iS_{\omega \rho }^{(\pi loop)}`$ $`=`$ $`{\displaystyle d^4xd^4y<0|T\{_{\omega \pi \pi }^{}(x)_{\rho \pi \pi }(y)\}|0>}`$ (123)
$`=`$ $`{\displaystyle \frac{N_c}{6\pi ^2}}{\displaystyle \frac{m_um_d}{m}}i{\displaystyle \frac{d^4q}{(2\pi )^4}b(q^2)s(q^2)\mathrm{\Sigma }_\pi (q^2)\frac{q^2(q^2\delta _{\mu \nu }q_\mu q_\nu )}{1+\mathrm{\Xi }(q^2)}\omega ^\mu (q)\rho ^{0\nu }(q)}.`$ (124)
Thus chain approximation of fig. 4 yield effective lagrangian as follow
$`_{\omega \rho }^{}={\displaystyle \frac{N_c}{6\pi ^2}}{\displaystyle \frac{m_um_d}{m}}{\displaystyle \frac{d^4q}{(2\pi )^4}e^{iqx}\{g^2h_0(q^2)\frac{q^2}{(1+\mathrm{\Xi }(q^2))}b(q^2)s(q^2)\mathrm{\Sigma }_\pi (q^2)\}(q^2\delta _{\mu \nu }q_\mu q_\nu )\omega ^\mu (q^2)\rho ^{0\nu }(x)}.`$ (125)
## V $`\omega \rho ^0`$ Mixing and $`\omega \pi ^+\pi ^{}`$ Decay
Due to VMD, the eleectromagnetic interaction contributes to $`\omega \pi ^+\pi ^{}`$ decay through $`\omega \gamma \rho ^0\pi \pi `$ and $`\omega \gamma \pi \pi `$. In I we have evaluated $`\rho \pi \pi `$ vertex, $`\rho ^0\gamma `$ mixing vertex and “direct” $`\gamma \pi \pi `$ vertex up to one-loop level. The “direct” $`\gamma \pi \pi `$ vertex reads
$`_{\gamma \pi \pi }^c={\displaystyle \frac{d^4q}{(2\pi )^4}e^{iqx}\overline{F}_\pi (q^2)A_\mu (q)[\pi ^+(x)^\mu \pi ^{}(x)^\mu \pi ^+(x)\pi ^{}(x)]},`$ (126)
where $`A_\mu `$ is photon field, $`\overline{F}_\pi (q^2)`$ is nonresonant background part of pion form factor. Explicitly, $`\overline{F}_\pi (q^2)`$ reads
$`\overline{F}_\pi (q^2)=1+{\displaystyle \frac{q^2b_\gamma (q^2)}{1+\mathrm{\Sigma }(q^2)}},`$ (127)
where
$`b_\gamma (q^2)`$ $`=`$ $`{\displaystyle \frac{b(q^2)}{2(1+3\zeta )}}D(q^2){\displaystyle \frac{C(q^2)\mathrm{\Sigma }_0(q^2)}{1+11\zeta /3}},`$ (128)
$`C(q^2)`$ $`=`$ $`{\displaystyle \frac{1}{2f_\pi ^2}}[A(q^2)+2g_A^2B(q^2)],`$ (129)
$`\mathrm{\Sigma }_0(q^2)`$ $`=`$ $`{\displaystyle \frac{2}{f_\pi ^2}}[2\mathrm{\Sigma }_\pi (q^2)\mathrm{\Sigma }_K(q^2)],`$ (130)
$`\mathrm{\Sigma }(q^2)`$ $`=`$ $`[1+{\displaystyle \frac{q^2C(q^2)}{1+11\zeta /3}}]\mathrm{\Sigma }_0(q^2),`$ (131)
$`D(p^2)`$ $`=`$ $`{\displaystyle \frac{1}{16\pi ^2f_\pi ^2}}\{\lambda +{\displaystyle _0^1}dxx(1x)\mathrm{ln}[(1{\displaystyle \frac{x(1x)p^2}{m__K^2}})({\displaystyle \frac{x(1x)p^2}{m__K^2}})^2]`$ (133)
$`{\displaystyle \frac{2}{3}}i\pi \theta (p^24m_\pi ^2)\}.`$
The complete $`\rho \pi \pi `$ vertex reads
$$_{\rho \pi \pi }^c=\frac{d^4q}{(2\pi )^4}e^{iqx}g_{\rho \pi \pi }(q^2)(q^2\delta _{\mu \nu }q_\mu q_\nu )\rho ^{0\mu }(q)[\pi ^+(x)^\mu \pi ^{}(x)^\mu \pi ^+(x)\pi ^{}(x)],$$
(134)
with
$$g_{\rho \pi \pi }(q^2)=\frac{b(q^2)}{(1+2\zeta )(1+\mathrm{\Sigma }(q^2))}.$$
(135)
Moreover, the complete $`\rho \gamma `$ mixing vertex reads
$$_{\rho \gamma }^c=\frac{1}{2}e\frac{d^4q}{(2\pi )^4}e^{iqx}b_{\rho \gamma }(q^2)(q^2\delta _{\mu \nu }q_\mu q_\nu )\rho ^{0\mu }(q)A^\nu (x),$$
(136)
where
$$b_{\rho \gamma }(q^2)=\frac{A(q^2)}{g(1+\zeta )}f_\pi ^2b(q^2)\frac{\mathrm{\Sigma }_0(q^2)}{1+2\zeta }[1+\frac{q^2b_\gamma (q^2)}{1+\mathrm{\Sigma }(q^2)}].$$
(137)
Thus due to VMD, the complete $`\omega \gamma `$ mixing vertex can be obtained directly
$$_{\omega \gamma }^{}=\frac{1}{6}e\frac{d^4q}{(2\pi )^4}e^{iqx}b_{\rho \gamma }(q^2)(q^2\delta _{\mu \nu }q_\mu q_\nu )\rho ^{0\mu }(q)A^\nu (x).$$
(138)
Eqs.( 136) and ( 138) will lead to $`\omega \rho ^0`$ mixing at the order of $`\alpha _{\mathrm{e}.\mathrm{m}.}`$ through the transition process $`\omega \gamma \rho ^0`$, which is
$$_{\omega \rho }^{\mathrm{e}.\mathrm{m}.}=\frac{1}{12}e^2\frac{d^4q}{(2\pi )^4}e^{iqx}b_{\rho \gamma }^2(q^2)(q^2\delta _{\mu \nu }q_\mu q_\nu )\omega ^\mu (q)\rho ^{0\nu }(x).$$
(139)
In addition, eqs.( 126) and ( 138) also lead to “direct” $`\omega \pi \pi `$ couping at the order of $`\alpha _{\mathrm{e}.\mathrm{m}.}`$ through the transition process $`\omega \gamma \pi \pi `$, which is
$`_{\omega \pi \pi }^{\mathrm{e}.\mathrm{m}.}`$ $`=`$ $`{\displaystyle \frac{i}{6}}e^2{\displaystyle \frac{d^4q}{(2\pi )^4}e^{iqx}\overline{F}_\pi (q^2)b_{\rho \gamma }(q^2)(q^2\delta _{\mu \nu }q_\mu q_\nu )}`$ (141)
$`\times \omega ^\mu (q)[\pi ^+(x)^\nu \pi ^{}(x)^\nu \pi ^+(x)\pi ^{}(x)].`$
Eq.( 139) together with eqs.( 93), (105) and (125) give the complete $`\omega \rho ^0`$ mixing vertex as follow
$$_{\omega \rho }^c=\frac{d^4q}{(2\pi )^4}e^{iqx}\mathrm{\Theta }_{\omega \rho }(q^2)(q^2\delta _{\mu \nu }q_\mu q_\nu )\omega ^\mu (q)\rho ^{0\nu }(x),$$
(142)
where vector meson fields have been normalized to physical fields, and
$`\mathrm{\Theta }_{\omega \rho }(q^2)={\displaystyle \frac{N_c}{6\pi ^2}}{\displaystyle \frac{m_um_d}{m}}\{g^2h_0(q^2)(1{\displaystyle \frac{4}{3}}\zeta )+q^2b(q^2)s(q^2)[\mathrm{\Sigma }_K(q^2){\displaystyle \frac{\mathrm{\Sigma }_\pi (q^2)}{1+\mathrm{\Xi }(q^2)}}]\}+{\displaystyle \frac{\alpha \pi }{3}}b_{\rho \gamma }^2(q^2).`$ (143)
The complete “direct” $`\omega \pi \pi `$ vertex can be obtained via summing eqs.(97), (112), (122) and (141),
$`_{\omega \pi \pi }^c`$ $`=`$ $`i{\displaystyle \frac{d^4q}{(2\pi )^4}e^{iqx}g_{\omega \pi \pi }(q^2)(q^2\delta _{\mu \nu }q_\mu q_\nu )\omega ^\mu (q)[\pi ^+(x)^\nu \pi ^{}(x)^\nu \pi ^+(x)\pi ^{}(x)]}.`$ (144)
where all meson fields have been normalized to physical fields, and $`g_{\omega \pi \pi }(q^2)`$ is defined as follow
$`g_{\omega \pi \pi }(q^2)`$ $`=`$ $`{\displaystyle \frac{N_c}{12\pi ^2}}{\displaystyle \frac{m_um_d}{m}}\{s(q^2)({\displaystyle \frac{1}{1+\mathrm{\Xi }(q^2)}}{\displaystyle \frac{10}{3}}\zeta )`$ (146)
$`6f_\pi ^2\mathrm{\Sigma }_K(q^2)[8m^2f_\pi ^2b(q^2){\displaystyle \frac{s(q^2)}{3}}(1+{\displaystyle \frac{q^2N_c}{4\pi ^2f_\pi ^2}})]\}+{\displaystyle \frac{2\alpha \pi }{3}}\overline{F}_\pi (q^2)b_{\rho \gamma }(q^2).`$
Thus G-parity forbidden $`\omega \pi ^+\pi ^{}`$ includes a nonresonant background contribution, eq.( 144), and $`\rho `$ resonance exchange contribution(eqs.( 134) and (142)). The decay width on $`\omega `$ mass-shell is
$`\mathrm{\Gamma }(\omega \pi ^+\pi ^{})={\displaystyle \frac{m_\omega ^5}{48\pi }}|{\displaystyle \frac{m_\omega ^2\mathrm{\Theta }_{\omega \rho }(m_\omega ^2)g_{\rho \pi \pi }(m_\omega ^2)}{m_\omega ^2m_\rho ^2+im_\rho \mathrm{\Gamma }_\rho }}g_{\omega \pi \pi }(m_\omega ^2)|^2(1{\displaystyle \frac{4m_\pi ^2}{m_\omega ^2}})^{3/2}.`$ (147)
Using the experimental data $`B(\omega \pi ^+\pi ^{})=(2.21\pm 0.30)\%`$ together with eq.( 147), we obtain
$$m_um_d=(3.9\pm 0.22)\mathrm{MeV}$$
(148)
at energy scale $`\mu m_\omega `$. Here the error bar is from the uncertainty in branch ratio of the process $`\omega \pi ^+\pi ^{}`$. In the standard way, the $`\omega \rho ^0`$ mixing amplitude is
$`{\displaystyle \frac{d^4q}{(2\pi )^4}\mathrm{\Pi }_{\omega \rho }(q^2)}=<\omega |{\displaystyle d^4x_{\omega \rho }(x)}|\rho >\mathrm{\Pi }_{\omega \rho }(q^2)`$ $`=`$ $`q^2\mathrm{\Theta }_{\omega \rho }(q^2).`$ (149)
The off-shell $`\omega \rho ^0`$ mixing amplitude is obviously momentum dependent, and vanished at $`q^2=0`$. This is consistent with the arguement by O’Connell et. al. in ref. that this mixing amplitude must vanish at the transition from time-like to space-like four momentum within a broad class of models. In addition, the value of isospin broken parameter( 148) leads on $`\omega `$ mass-shell $`\omega \rho ^0`$ mixing amplitude as follow
$$\mathrm{Re}\mathrm{\Pi }_{\omega \rho }(m_\omega ^2)=(3956\pm 280)\mathrm{MeV}^2,\mathrm{Im}\mathrm{\Pi }_{\omega \rho }(m_\omega ^2)=(1697\pm 130)\mathrm{MeV}^2.$$
(150)
In ref., the on-shell mixing amplitude has extracted from the $`e^+e^{}\pi ^+\pi ^{}`$ experimental data in a model-dependent way. In eq.( 150), the real part of on-shell mixing amplitude agree with result of ref.. The imaginary part, however, is much larger than one in ref. which is around $`300`$MeV<sup>2</sup>. It must be pointed out that, in ref. the author’s analysis bases on a model without “direct” $`\omega \pi \pi `$ coupling. Therefore, it is insignificant to compare the value of on-shell mixing amplitude of this the present paper with one of ref.. Fortunately, the ratio between $`\omega \pi \pi `$ decay amplitude and $`\rho \pi \pi `$ decay amplitude should be model-indenpendent. This value can test whether a model is right or not. The on-shell mixing amplitude in ref. yields
$$R_{\omega \rho }^{\mathrm{exp}}=\frac{<\pi ^+\pi ^{}|\omega >}{<\pi ^+\pi ^{}|\rho >}=(0.0060\pm 0.0009)+(0.0322\pm 0.0050)i.$$
(151)
The present paper predicts
$$R_{\omega \rho }=\frac{m_\omega ^2\mathrm{\Theta }_{\omega \rho }(m_\omega ^2)}{m_\omega ^2m_\rho ^2+im_\rho \mathrm{\Gamma }_\rho }\frac{g_{\omega \pi \pi }(m_\omega ^2)}{g_{\rho \pi \pi }(m_\omega ^2)}=(0.0084\pm 0.0007)+(0.0331\pm 0.0021)i.$$
(152)
We can see that this theoretical prediction agree with experimental excellently.
Moreover, the follows have also been revealed in our studies of this paper:
i) If we take $`g_{\omega \pi \pi }(q^2)=0`$ in eq.( 147), we have $`B(\omega \pi ^+\pi ^{})=(2.56\pm 0.34)\%`$. So that the contribution from interference between “direct” $`\omega \pi \pi `$ coupling and $`\omega \rho ^0`$ mixing is about $`15\%`$. The dominant contribution are from $`\rho `$-resonance exchange. This conclusion indicates all pervious studies which without “direct” $`\omega \pi \pi `$ coupling are good approximation even though this neglect is an ad hoc assumption. However, in mechanism of $`\omega \pi ^+\pi ^{}`$ with “direct” $`omega\pi \pi `$ coupling, larger imagnary part of on-shell $`\omega \rho ^0`$ mixing amplitude is allowed, but it is not allowed in the mechnism without the direct couping.
ii) If we do not consider the contributions from one-loop diagrams of pseudoscalar mesons, i.e. setting $`\mathrm{\Sigma }_K(q^2)=\mathrm{\Sigma }_\pi (q^2)=0`$, we obtain $`B(\omega \pi ^+\pi ^{})=(2.86\pm 0.47)\%`$. Thus the contribution from one-loop of pseudoscalar mesons is about $`30\%`$ and can not be omitted. This conclusion is consistent with I. In addition, in this case, the on-shell $`\omega \rho ^0`$ mixing amplitude is about $`4700`$MeV<sup>2</sup>. So that we can see that the larger imagnary part of on-shell $`\omega \rho ^0`$ mixing amplitude is yielded by pseudoscalar meson loops. In I, we have shown that this larger imagnary part is required by the unitarity of this effective field theory.
## VI Summary
In this paper, we study G-parity forbidden $`\omega \pi ^+\pi ^{}`$ decay up to one-loop level of mesons. This process is yielded by isospin symmetry breaking due to $`m_um_d`$ and electromagnetic interaction of mesons. The decay amplitude contains two parts of contributions which are from “direct” $`\omega \pi \pi `$ couping and $`\omega \rho ^0`$ mixing respectively. In the previous studies, the “direct” $`\omega \pi \pi `$ couping is neglected. We show that the “direct” $`\omega \pi \pi `$ couping and its interference with $`\omega \rho ^0`$ mixing contribute to on-shell decay amplitude about 15$`\%`$ only. It also interprets why the previous studies are good approximations even without “direct” $`\omega \pi \pi `$ couping. We suggest that the decay amplitude ratio $`R_{\omega \rho }`$ should be model-independent, and our prediction agree with experimental data excellently.
The formula of $`\omega \rho ^0`$ mixing amplitude is also obtained. Since our calculation is beyond the chiral expansion(including all orders contribution of the chiral expansion) and one-loop contribution of pseudoscalar mesons is considered, the momentum-dependence of the off-shell mixing amplitude is very complicated. However, the mixing amplitude also vanishes at $`q^2=0`$. For case of on $`\omega `$ mass-shell, the mixing amplitude emerges larger imagnary which is from one-loop contribution of pion and is required by unitarity of this effective field theory.
In our calculation, all vertices are expanded to the leading order light current quark masses. At this order, the decay amplitude yielded by isospin broken is proportional to $`m_dm_u`$. The theorectical prediction of isospin breaking parameter is $`m_dm_u=(3.9\pm 0.22)`$MeV at energy scale $`\mu m_\omega `$. This value is important for determining light quark masses at vector meson energy scale.
Appendix
Here we provide some identities which are used in calculation on four-point effective action. In sect. 3.2 we have used $`q`$, $`k`$ and $`kq`$ to denote four-momentum square of external source $`\overline{V}_\mu `$, $`\mathrm{\Delta }_\nu `$ and $`\mathrm{\Delta }_\sigma `$. For the purpose of this paper, $`\mathrm{\Delta }_\nu (k)`$ reduces to $`k_\nu \pi (k)`$. So that due to soft pion theorem we have $`k^20`$, $`(q+k)^20`$ and $`k_\mu \mathrm{\Delta }^\mu (k)0`$. Moreover, due to space-like condition of vector meson fields, $`q^\mu V_\mu (q)=0`$, we have
$`(\delta _{\mu \sigma }q_\nu +\delta _{\mu \nu }q_\sigma )<\{V^\mu (q),S\}\mathrm{\Delta }^\nu (k)\mathrm{\Delta }^\sigma (kq)>=\mathrm{\hspace{0.33em}0},`$ (153)
$`(\delta _{\mu \sigma }q_\nu +\delta _{\mu \nu }q_\sigma )<V^\mu (q)\mathrm{\Delta }^\nu (k)S\mathrm{\Delta }^\sigma (kq)>=\mathrm{\hspace{0.33em}0}.`$ (154)
$`q_\nu q_\sigma k_\mu <\{V^\mu (q),S\}\mathrm{\Delta }^\nu (k)\mathrm{\Delta }^\sigma (kq)>{\displaystyle \frac{q^2}{2}}q_\sigma \delta _{\mu \nu }<\{V^\mu (q),S\}\mathrm{\Delta }^\nu (k)\mathrm{\Delta }^\sigma (kq)>,`$ (155)
$`q_\nu q_\sigma k_\mu <V^\mu (q)\mathrm{\Delta }^\nu (k)S\mathrm{\Delta }^\sigma (kq)>{\displaystyle \frac{q^2}{2}}q_\sigma \delta _{\mu \nu }<V^\mu (q)\mathrm{\Delta }^\nu (k)S\mathrm{\Delta }^\sigma (kq)>,`$ (156)
$`k_\mu \delta _{\nu \sigma }<\{V^\mu (q),S\}\mathrm{\Delta }^\nu (k)\mathrm{\Delta }^\sigma (kq)>q_\sigma \delta _{\mu \nu }<\{V^\mu (q),S\}\mathrm{\Delta }^\nu (k)\mathrm{\Delta }^\sigma (kq)>,`$ (157)
$`k_\mu \delta _{\nu \sigma }<V^\mu (q)\mathrm{\Delta }^\nu (k)S\mathrm{\Delta }^\sigma (kq)>q_\sigma \delta _{\mu \nu }<V^\mu (q)\mathrm{\Delta }^\nu (k)S\mathrm{\Delta }^\sigma (kq)>,`$ (158)
The following integral identities are also used in our calculation
$`{\displaystyle _0^1}𝑑xx^2{\displaystyle _0^1}𝑑y(1y){\displaystyle \frac{xy}{[l^2m^2+xy(1x)q^2]^4}}={\displaystyle _0^1}𝑑xx^2{\displaystyle _0^1}𝑑y(1y){\displaystyle \frac{(1x)}{[l^2m^2+xy(1x)q^2]^4}}`$ (159)
$`=`$ $`{\displaystyle _0^1}𝑑xx^2{\displaystyle _0^1}𝑑y(1y){\displaystyle \frac{x(12y)}{[l^2m^2+x^2y(1x)q^2]^4}}={\displaystyle _0^1}𝑑xx^2{\displaystyle _0^1}𝑑y(1y){\displaystyle \frac{(1x)}{[l^2m^2+x^2y(1x)q^2]^4}}.`$ (160)
ACKNOWLEDGMENTS
This work is partially supported by NSF of China through C. N. Yang and the Grant LWTZ-1298 of Chinese Academy of Science.
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# Reorientation of Anisotropy in a Square Well Quantum Hall Sample
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## Abstract
We have measured magnetotransport at half-filled high Landau levels in a quantum well with two occupied electric subbands. We find resistivities that are isotropic in perpendicular magnetic field but become strongly anisotropic at $`\nu `$ = 9/2 and 11/2 on tilting the field. The anisotropy appears at an in-plane field, $`B_{ip}2.5`$ T, with the easy-current direction parallel to $`B_{ip}`$ but rotates by 90 at $`B_{ip}10`$ T and points now in the same direction as in single-subband samples. This complex behavior is in quantitative agreement with theoretical calculations based on a unidirectional charge density wave state model.
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A two-dimensional (2D) electron gas is an attractive system for many-body physics studies. A particularly rich variety of phenomena associated with strong interactions among electrons appears in the regime of the fractional quantum Hall effect (FQHE) . During much of the last decade, studies of the FQHE have focused on even-denominator Landau level filling factors such as the compressible $`\nu =1/2`$ state and the $`\nu =5/2`$ incompressible quantum Hall fluid. Most recently, strongly anisotropic transport has been observed in high quality GaAs/Al<sub>x</sub>Ga<sub>1-x</sub>As single heterojunctions at filling factors $`\nu `$ = 9/2, 11/2, etc., and in 2D hole systems starting at $`\nu =5/2`$. In these experiments, the magnetoresistance shows a strong peak in one current direction and a deep minimum in the perpendicular current direction. Tilting the magnetic field away from the sample normal causes the high resistance direction to change from its original orientation to the in-plane magnetic field direction.
The origin of the magnetotransport anisotropy has not been firmly established yet. The most appealing interpretation suggests that the 2D electron gas spontaneously breaks the translational symmetry by forming a unidirectional charge density wave (UCDW), as predicted by Hartree-Fock theory . This idea has spurred much theoretical interest . Because of uncertainty about the reliability of this Hartree-Fock prediction, there has been a special emphasis placed on tests of its ability to explain experimental results on “stripe” orientation in tilted magnetic fields. In particular, Jungwirth et al. carried out detailed many-body RPA/Hartree-Fock calculations combined with a self-consistent local-spin-density-approximation (LSDA) description of one-particle states in experimental sample geometries. For the sample parameters of the traditional, single-interface specimens of Refs. with a single electric subband occupied, the theory gives stripes oriented perpendicular to the field, consistent with experiment.
A theoretical study of UCDWs in parabolic quantum wells that have two subbands occupied in zero magnetic field, has predicted much more complex behavior of the UCDW state, including stripe states induced by an in-plane field and rotation of stripe orientation at critical in-plane field strengths. A comparison between theory and experiment in a geometry for which this intricate behavior occurs, constitutes an excellent test of the UCDW explanation of anisotropic transport in higher Landau levels. Since parabolic quantum wells are experimentally difficult to realize and suffer from poor mobility we, instead, chose a square well structure which is expected to exhibit similarly complex behavior, provided that more than one electric subband is occupied in zero field.
Our sample, detailed in Figure 1(c), consists of a 350Å wide GaAs quantum well bracketed between thick Al<sub>0.24</sub>Ga<sub>0.76</sub>As layers grown on a (100) GaAs substrate by MBE. Two Si delta-doping layers are placed symmetrically above and below the quantum well at a distance of 800Å. The specimen has a size of 5mm $`\times `$ 5mm and is contacted via eight indium contacts, placed symmetrically around the perimeter. The electron density is established after illuminating the sample with a red light-emitting diode at $``$4.2K and we measure an electron mobility of $`\mu `$ = 7 $`\times `$ 10<sup>6</sup> cm<sup>2</sup>/V s. The total electron density, $`n=4.6\times 10^{11}`$ cm<sup>-2</sup>, is determined from low-field Hall data. The subband densities, $`n_1=3.3\times 10^{11}`$ cm<sup>-2</sup> and $`n_2=1.3\times 10^{11}`$ /cm<sup>-2</sup>, are obtained by Fourier analysis of the low-field Shubnikov-de Haas oscillation. Their values coincide with the results of our numerical self-consistent LSDA calculation. All angular-dependent measurements were carried out at T = 40 mK in a top-loading dilution refrigerator equipped with an in-situ rotator placed inside a 33T resistive magnet. A low-frequency ($``$ 7 Hz) lock-in technique at a current $`I`$ = 10 nA is used. We define the axis of rotation as the y-axis. Consequently, the in-plane field, $`B_{ip}`$, is along the x-axis when the sample is rotated. Therefore, $`R_{xx}`$ refers to “$`I`$ parallel to $`B_{ip}`$” and $`R_{yy}`$ refers to “$`I`$ perpendicular to $`B_{ip}`$.
Figure 1(a) shows an overview of magnetoresistance at zero-tilt. The shaded region highlights the transport features around $`\nu `$ = 9/2 and 11/2. The integer quantum Hall effect (IQHE) states at $`\nu =1,2,3,\mathrm{}`$ and the FQHE states at $`\nu `$ = 2/3, etc. are clearly visible. Figure 1(b) shows the results of self-consistent LSDA calculations of Landau levels (measured from the bottom of the quantum well) in perpendicular magnetic field.
Figure 2 shows the $`R_{xx}`$ and $`R_{yy}`$ data for $`4<\nu <6`$ at four different tilt angles $`\theta `$ = 0, 41.2, 67.9, and 76.2. The tilt angle is determined using the shift of prominent QHE states, which depend only on the perpendicular magnetic field, $`B_{perp}=B\times `$ cos$`\theta `$.
In the absence of $`B_{ip}`$ ($`\theta `$ = 0), $`R_{xx}`$ and $`R_{yy}`$ show a peak at $`\nu `$ = 9/2 and a slight dip at $`\nu `$ = 11/2 and negligible anisotropy. The small difference in magnitude between $`R_{xx}`$ and $`R_{yy}`$ is probably a result of the different contacts involved in both measurements. This practically isotropic behavior of $`R_{xx}`$ and $`R_{yy}`$ is distinctively different from results on single-subband, single-heterojunctions, where the states at $`\nu `$ = 9/2 and 11/2 are strongly anisotropic in the absence of $`B_{ip}`$. This lack of anisotropy in our sample has a simple interpretation. The diagram in Figure 1(b) indicates that the $`\nu `$ = 9/2 and 11/2 states are the $`\nu `$ = 3/2 state of the lowest Landau level (N=0) in the second quantum well subband (i=2). The $`\nu `$ = 3/2 state in single subband samples exhibits isotropic transport, which seems to carry over to the second subband. Yet, exceptional behavior develops on tilting the specimen.
At $`\theta =41.2^{}`$ the $`R_{xx}`$ and $`R_{yy}`$ traces are very different from those taken at zero field-tilt and different from each other. The $`\nu `$ = 9/2 and $`\nu `$ = 11/2 states are strongly anisotropic with the hard-axis perpendicular to $`B_{ip}`$ ($`R_{yy}`$) and the easy-axis parallel to $`B_{ip}`$ ($`R_{xx}`$). The direction of this tilt-induced anisotropy (TIA) is rotated by 90 as compared to the direction in traditional single-subband, single-heterojunction structures . As the tilt angle increases further, the $`R_{xx}`$ and $`R_{yy}`$ traces approach each other again at $`\theta 67.9^{}`$ rendering the transport nearly isotropic (Figure 2(c)). Beyond this angle the anisotropy reemerges but the hard-axis and easy-axis have traded places, as seen in Figure 2(d).
In Figures 3(a,b) we plot $`R_{xx}`$ and $`R_{yy}`$ at filling factors $`\nu `$ = 9/2 and 11/2 versus $`B_{ip}`$. Their general behavior is rather similar. Practically isotropic transport prevails in the range of $`0<B_{ip}<2`$ T, but there is a clear onset to anisotropy at $`B_{ip}2.5`$ T. The level of anisotropy rapidly increases, reaching its peak at $`B_{ip}5.0`$ T, whereupon the $`R_{xx}`$ and $`R_{yy}`$ values approach each other again and cross at $`B_{ip}10`$ T. For higher in-plane fields the transport is again anisotropic, but its direction has rotated by 90. Figures 3(c,d) show the anisotropy factor, defined as $`(R_{xx}R_{yy})/(R_{xx}+R_{yy})`$ and derived from the data of the panels above. They clearly depict the initially, practically isotropic behavior followed by a strong anisotropy that rotates direction by 90 at $`B_{ip}10`$ T. The direction of anisotropy in single-subband samples corresponds to the high $`B_{ip}`$ direction in our double-subband specimen.
We now turn to the analysis of correspondence between the measured TIA and the theory based on the UCDW picture. For an infinitely narrow electron layer the effective 2D Coulomb interaction, $`V(\stackrel{}{q})`$, reduces to $`e^{q^2\mathrm{}^2/2}/q(L_N(q^2/2))^2\mathrm{\hspace{0.17em}2}\pi e^2\mathrm{}/ϵ`$ where $`L_N(x)`$ is the Laguerre polynomial, $`\stackrel{}{q}`$ is the wavevector, $`\mathrm{}`$ is the magnetic length, and $`ϵ`$ is the dielectric function. Starting from $`N=1`$, zeros of $`L_N(q^2/2)`$ occur at finite $`q=q^{}`$, producing a zero in the repulsive Hartree interaction at wave vectors where the attractive exchange interaction is strong. For the half-filled valence Landau level the corresponding UCDW state consists of alternating occupied and empty stripes of electron guiding center states with a modulation period $`2\pi /q^{}`$.
In finite-thickness 2D systems subjected to tilted magnetic fields, the dependence of the effective interaction on wavevector magnitude $`q`$ and orientation $`\varphi `$ relative to the in-plane field direction can be accurately approximated by $`V(\stackrel{}{q})=V_0(q)+V_2(q)\mathrm{cos}(2\varphi )`$. At $`B_{ip}=0`$, the isotropic term $`V_0(q)`$ has a wavevector-dependence similar to that of the effective interaction in the infinitely narrow 2D layer. The corresponding curve for the valence Landau level at $`\nu =9/2`$, shown in the top inset of Figure 4, has no zeros at finite $`q`$-vectors because the half-filled valence Landau level is the $`N=0`$ state of the second subband (as shown in detail in Figure 1(b)) . Hence, the UCDW state is not expected to form, consistent with the isotropic transport measured in perpendicular field.
Because of the finite thickness of the 2D system in our 350Å wide quantum well, the orbital effect of the in-plane field causes Landau levels emanating from different electric subbands to coincide, depending on the strength of $`B_{ip}`$. The in-plane field mixes electric and magnetic levels so the subband and orbit radius indices are no longer good quantum numbers. However, the effect of $`B_{ip}`$ near the level (anti)crossing can sometimes be viewed approximately as a transfer of valence electrons from the lowest ($`N=0`$) Landau level of the second subband to a higher ($`N>0`$) Landau level of the first subband. For filling factor $`\nu =9/2`$, such a circumstance occurs in our sample at $`B_{ip}3`$ T, as seen from the top and bottom insets of Figure 4. Indeed, $`V_0(q)`$ is only slightly modified at low in-plane fields, while a clear minimum develops for $`B_{ip}>3`$ T. As discussed above for the case of perpendicular magnetic field it is the minimum of the interaction energy at finite wavevector that opens the possibility for the formation of the UCDW state. The theoretical and experimental critical in-plane fields corresponding to the onset of the UCDW and TIA, respectively, are remarkably close.
The non-zero anisotropy coefficient $`V_2(q)`$ of the effective interaction at $`B_{ip}>0`$ is responsible for the non-zero UCDW anisotropy energy $`E_A`$, defined as the total Hartree-Fock energy of stripes oriented parallel with $`B_{ip}`$ minus the total energy of stripes perpendicular to $`B_{ip}`$. The direction of the anisotropy results from a delicate competition between electrostatic and exchange contributions to $`E_A`$ and can be determined only by an accurate calculation which takes into account details of the experimental configuration. As shown in Figure 4, the stripes align parallel with $`B_{ip}`$ at low in-plane fields, consistent with the measured easy-current direction parallel with $`B_{ip}`$. The sign of the UCDW anisotropy energy changes at $`B_{ip}=10`$ T which coincides with the experimental critical field for the interchange of easy and hard current axes. This theoretical discussion of the $`\nu =9/2`$ state was found to apply for $`\nu =11/2`$ as well.
In conclusion, we have observed complex transport behavior in a two-subband QW at half-filled high Landau levels. Both the transition to an anisotropic transport state, at finite $`B_{ip}`$, and the rotation of the direction of anisotropy by 90 at higher $`B_{ip}`$ are explained quantitatively by the UCDW picture. The close agreement between complex experimental data and theoretical results leaves little doubt as to the origin of the observed transport anisotropies in high Landau levels.
We would like to thank E. Palm and T. Murphy for experimental assistance, and N. Bonesteel, R.R. Du, and K. Yang for useful discussion. A portion of this work was performed at the National High Magnetic Field Laboratory which is supported by NSF Cooperative Agreement No. DMR-9527035 and by the State of Florida. The work at Indiana University was supported by NSF grant DMR-9714055, and at the Institute of Physics ASCR by the Grant Agency of the Czech Republic under grant 202/98/0085. D.C.T. and W.P. are supported by the DOE and the NSF.
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# Realistic regularization of the QED Green functions
## I Introduction
Perturbative predictions about quantum-electrodynamic phenomena implied by a QED Lagrangian can be computed using the Feynman rules, a regularization method to circumvent ultraviolet divergencies, and a renormalization scheme. Regularization method results in regularized n-point Green functions; a suitable limiting procedure (a renormalization scheme) then leads to physically sensible predictions that are independent of the particular regularization method used. But no known regularized n-point Green functions can be regarded as being based on physically realistic premises about quantum-electrodynamic phenomena: the derivation of each is formalistic since it disregards some of the basic tenets of conventional physics (e.g., by lacking a Lagrangian, by not being Lorentz-invariant, by introducing particles with wrong metric or statistics…). So the perturbative predictions of QED presently cannot be directly derived from physically realistic premises; for a history of, and comments on this basic, conceptual inconsistency see, e.g.,. Dirac believed that removal of this conceptual inconsistency may lead to an important advance in field theories.
To show that one can remove this inconsistency already in four-dimensional space-time, we will introduce a new, physically motivated modification of the QED Lagrangian and consider it within the theoretical framework of ’t Hooft and Veltman that presents an alternative to the convential perturbative quantum field theory: They avoid canonical formalism and take diagrams as the basis from which everything must be derived; so they give a perturbative definition of the S-matrix directly in terms of diagrams corresponding to a given Lagrangian as specified by postulated Feynman rules. The question is: How do we modify the QED Lagrangian so that the resulting regularized S-matrix is derived from physically realistic premises?
We are using the adjectives formalistic and realistic in the sense of Pauli and Villars. Introducing their formalistic regularization method, they remarked: ”It seems very likely that the ’formalistic’ viewpoint used in this paper and by other workers can only be a transitional stage of the theory, and that the auxiliary masses will eventually be entirely eliminated, or the ’realistic’ standpoint will be so much improved that the theory will not contain any further accidental compensations.” Which we intend to do.
Gupta has shown already in 1952 that one can modify the QED Lagrangian so that the new Lagrangian results in the S-matrix of QED regularized by certain Pauli-Villars method. And twenty years later ’t Hooft and Veltman introduced the method of unitary regulators (HV-method) that (i) is a variant of Pauli-Villars methods for regularizing propagators, (ii) requires only an exceedingly simple modification of the initial Lagrangian, and (iii) is very suitable for proving the causality of the regularized $`n`$-point Green functions and the unitarity of the resulting S-matrix. Unfortunately both methods are formalistic since they introduce also unphysical, auxiliary particles with wrong metric or statistics. To get rid of this serious conceptual deficiency, we will generalize the HV-method to avoid auxiliary particles.
We will demonstrate the utility of the generalized HV-method by showing that there are finite perturbative $`n`$-point Green functions of quantum-electrodynamic phenomena derived from a realistic perturbative theory (a rp-theory, for short) such that:
* A rp-theory of quantum-electrodynamic phenomena is specified in a continuous, four-dimensional space-time by a local, Lorentz-invariant, physically motivated modification of a QED Lagrangian.
* The Feynman rules for this modified Lagrangian, defined as specified by ’t Hooft and Veltman, result in regularized Feynman diagrams that equal the diagrams of QED but with regularized propagators that have no additional singularities.
* All constants of a rp-theory are measurable in principle; there are no auxiliary parameters or particles.
* For certain values of these constants, the QED propagators are such low-energy approximations to their regularizations as acceptable for renormalization.
* The n-point Green functions of a rp-theory, defined as specified by ’t Hooft and Veltman in terms of regularized Feynman diagrams, are C-, P-, T- and Lorentz-invariant; causal; and charge and total four-momentum conserving.
Such a rp-theory of quantum-electrodynamic phenomena is not yet known; we cannot incorporate a finite-cutoff, Pauli-Villars, dimensional, or lattice regularization of QED in a rp-theory.
## II Lorentz-invariant regularization without additional singularities
As in the HV-method to each additional singularity of a regularized Feynman propagator corresponds an additional particle, we will first specify Lorentz-invariantly regularized Feynman propagators that have no additional singularities and have the Källén-Lehman representation used in proving causality and unitarity. Regarding metric and other conventions we follow Refs. ; in particular, a four-vector $`k=(\stackrel{}{k},ik^0)`$, and $`k^2\stackrel{}{k}\stackrel{}{k}(k^0)^2`$.
Consider a Lorentz-invariantly regularized spin 0 Feynman propagator, say, $`\mathrm{\Delta }_F(x)`$ whose space-time Fourier transform
$$(2\pi )^4i\stackrel{~}{\mathrm{\Delta }}_F(k)=\phi (k^2)(k^2+m^2iϵ)^1,\phi (m^2)=1,$$
(1)
where: (a) $`\phi (z)`$ is an analytic function of complex variable $`z`$ with a finite discontinuity somewhere across the segment $`zz_d<m^2`$ of the negative real axis; (b) $`|\phi (z)|<A|z|^r`$ with $`r3/2`$ as $`|z|\mathrm{}`$; (c) $`\phi (z)`$ is real on the positive real axis; (d) $`\phi (z)`$ depends on some constant $`\mathrm{\Lambda }`$ so that for any $`\mathrm{\Lambda }\mathrm{\Lambda }_0>0`$ it has properties (a) to (c) and
$$\underset{|z|<z_0}{sup}|\phi ^{(n)}(z)\delta _{n0}|0\text{as}\mathrm{\Lambda }\mathrm{}\text{for any}z_0>0,n=0,1,2,$$
(2)
and
$$\underset{z0,\mathrm{\Lambda }\mathrm{\Lambda }_0}{sup}|z^{(n+3)/2}\phi ^{(n)}(z)|<\mathrm{},n=0,1,2.$$
As a consequence, the spin 0 propagator provides a low-energy approximation to its regularization (1) which itself is acceptable for renormalization.
Using Cauchy’s integral formula we can conclude that the Lorentz-invariant regularization (1) of the spin 0 Feynman propagator admits the Källén-Lehman representation
$$(2\pi )^4i\stackrel{~}{\mathrm{\Delta }}_F(k)=_0^{\mathrm{}}\frac{\rho (s)}{k^2+siϵ}𝑑s$$
(3)
with
$$\rho (s)=\delta (sm^2)+(2\pi i)^1(m^2s)^1\underset{y0}{lim}[\phi (siy)\phi (s+iy)],$$
(4)
$`s,y>0`$. Note that $`\rho (s)`$ is real, $`\rho (s)=O(s^r)`$ as $`s\mathrm{}`$, and
$$_0^{\mathrm{}}s^m\rho (s)𝑑s=0\text{for}m=0,1,\mathrm{}<r1.$$
(5)
So we can decompose the regularized spin 0 propagator $`\mathrm{\Delta }_F(x)`$ into positive and negative energy parts: $`\mathrm{\Delta }_F(x)=\mathrm{\Theta }(x_0)\mathrm{\Delta }^+(x)+\mathrm{\Theta }(x_0)\mathrm{\Delta }^{}(x)`$.
The function $`i(2\pi )^4(\sqrt{\mathrm{\Lambda }^2m^2}+\mathrm{\Lambda })^n(k^2+m^2iϵ)^1(\sqrt{k^2+\mathrm{\Lambda }^2iϵ}+\mathrm{\Lambda })^n`$, $`\mathrm{\Lambda }>m`$, $`n=1,2,\mathrm{}`$, is an example of a Lorentz-invariantly regularized spin 0 Feynman propagator that satisfies the above conditions with $`r=n/2`$. Unfortunately, we cannot use such propagators for a realistic regularization of the QED Green functions since we do not know how to construct the corresponding local, Lorentz-invariant Lagrangians.
A propagator that satisfies conditions (a)-(c) is by (3) a generalization of the spin-0 propagator regularized by a Pauli-Villars regulator that has a continuous mass spectrum. Thus, to use such propagators to construct a rp-theory, we have to extend the ’t Hooft-Veltman construction of Lagrangians in HV-method to an infinite number of additional fields. To provide an example of how this can be done, we will present a local, Lorentz-invariant Lagrangian whose propagators for interacting fields can be taken as spin 1 and spin $`\frac{1}{2}`$ propagators regularized so that they acquire no additional singularities and have the Källén-Lehman representation.
## III An example of Lagrangian that regularizes QED propagators
Following Veltman , we will consider QED with massive photons in unitary gauge. Its Lagrangian reads
$$_{\mathrm{QED}}=\frac{1}{4}(_\mu A_\nu _\nu A_\mu )^2\frac{1}{2}\mu ^2A^2\overline{\psi }(\gamma ^\mu \begin{array}{c}\\ \end{array}_\mu +m)\psi +ie\overline{\psi }\gamma ^\mu \psi A_\mu +A_\mu J_\mu +\overline{J_e}\psi +\overline{\psi }J_e,$$
(6)
where $`J_\mu (x)`$, $`\overline{J_e}(x)`$, and $`J_e(x)`$ are four-vector and bispinor source fields, and $`\mu `$ is the non-vanishing photon mass—a physical constant $`<2\times 10^{16}`$ eV . The Feynman propagators for the four-vector field $`A_\mu (x)`$ and for the bispinor field $`\psi (x)`$ are:
$$i(2\pi )^4\frac{\delta _{\mu \nu }+\mu ^2k_\mu k_\nu }{k^2+\mu ^2iϵ},i(2\pi )^4\frac{i\gamma ^\mu k_\mu +m}{k^2+m^2iϵ}.$$
(7)
We could use $`_{\mathrm{QED}}`$ to define a rp-theory as specified in Section I, were the propagators (7) faster decreasing when $`k^2`$ tends to infinity.
However, one can modify the QED Lagrangian (6) so that the propagators for the fields $`A_\mu `$ and $`\psi `$ are such regularizations of propagators (7) that have no additional singularities when calculated according to the generalized ’t Hooft-Veltman method. Take, for example, the following real-valued, local, Lorentz-invariant Lagrangian:
$$_{\mathrm{TR}}=_1_{1/2}+ie\overline{\psi }\gamma ^\mu \psi A_\mu +A_\mu J_\mu +\overline{J_e}\psi +\overline{\psi }J_e$$
(9)
with
$`_1`$ $``$ $`q_1^1{\displaystyle d^4p\mathrm{\Psi }_\mu ^{}(x,p)[\mathrm{\Lambda }t(p^2)+p^\nu \begin{array}{c}\\ \end{array}_\nu ]\mathrm{\Psi }^\mu (x,p)}`$ (14)
$`+q_1^1s_1{\displaystyle }d^4pd^4p^{}f(p_{}^{}{}_{}{}^{2})f(p^2)[\mathrm{\Psi }_\mu ^{}(x,p^{})\mathrm{\Psi }^\mu (x,p)+p_\nu ^{}p^\nu \mathrm{\Psi }_\mu (x,p^{})\mathrm{\Psi }^\mu (x,p)`$
$`p^\mu \mathrm{\Psi }_\mu (x,p^{})p_{}^{}{}_{}{}^{\nu }\mathrm{\Psi }_\nu (x,p)],`$
$`_{1/2}`$ $``$ $`q_{1/2}^1{\displaystyle d^4p\overline{\mathrm{\Psi }}_{1/2}(x,p)[\mathrm{\Lambda }t(p^2)+p^\mu \begin{array}{c}\\ \end{array}_\mu ]\mathrm{\Psi }_{1/2}(x,p)}`$ (18)
$`q_{1/2}^1s_{1/2}{\displaystyle }d^4p^{}d^4pf(p_{}^{}{}_{}{}^{2})f(p^2)[\overline{\mathrm{\Psi }}_{1/2}(x,p^{})\gamma ^\mu \mathrm{\Psi }_{1/2}(x,p)p_\mu +\mathrm{c}.\mathrm{c}.],`$
$$A_\mu (x)d^4pf(p^2)\mathrm{\Psi }_\mu (x,p),\psi (x)d^4pf(p^2)\mathrm{\Psi }_{1/2}(x,p),$$
(19)
where: $`\mathrm{\Psi }_\mu (x,p)`$ and $`\mathrm{\Psi }_\mu ^{}(x,p)`$ are four-vector-valued functions of two four-vectors $`x`$ and $`p`$; $`\mathrm{\Psi }_{1/2}(x,p)`$ is a bispinor-valued function of $`x`$ and $`p`$; $`2a\begin{array}{c}\\ \end{array}_\mu ba(_\mu b)(_\mu a)b`$; $`\overline{\mathrm{\Psi }}_{1/2}\mathrm{\Psi }_{1/2}^{}\gamma ^4`$; $`t(p^2)`$ and $`f(p^2)`$ are real-valued functions of real $`p^2`$, $`d^4pf^2(p^2)=1`$; $`q_1`$, $`s_1`$, $`q_{1/2}`$, $`s_{1/2}`$, and $`\mathrm{\Lambda }`$ are real constants—not auxiliary parameters.
There are three kinds of reasons for the chosen form (III) of the Lagrangian $`_{\mathrm{TR}}`$:
* It is the purpose of this paper to show that there are Lagrangians that generalize the t’Hooft and Veltman method of unitary regulators to an infinite number of additional fields but do not introduce additional particles. So we constructed the Lagrangian $`_{\mathrm{TR}}`$ modifying $`_{\mathrm{QED}}`$ on the analogy of HV-method: (i) We introduced an infinite number of four-vector and bispinor fields of $`x`$ that have a continuous index $`p`$, namely $`\mathrm{\Psi }_\mu (x,p)`$, $`\mathrm{\Psi }_\mu ^{}(x,p)`$, and $`\mathrm{\Psi }_{1/2}(x,p)`$. (ii) We replaced the free part of $`_{\mathrm{QED}}`$ with the free Lagrangian of these fields, $`_1_{1/2}`$, which is of the first order in $``$ and has a nondiagonal mass matrix. (iii) In the interaction and source terms of $`_{\mathrm{QED}}`$, we replaced the fields $`A_\mu (x)`$ and $`\psi (x)`$ with weighted integrals (19) of $`\mathrm{\Psi }_\mu (x,p)`$ and $`\mathrm{\Psi }_{1/2}(x,p)`$ over the continuous index $`p`$.
* We tried to simplify the calculations of regularized propagators. We could do without the four-vector function $`\mathrm{\Psi }_\mu ^{}(x,p)`$ which we introduced solely to be able to use the same functions $`t(p^2)`$ and $`f(p^2)`$ in $`_1`$ and $`_{1/2}`$. We introduced $``$ $``$ to make $`_{\mathrm{TR}}`$ itself real-valued, not only its action real as required.
* The physical motivation for the type of Lagrangian we constructed, which we considered in detail in Ref. , is twofold: (i) The Euler-Lagrange equations of $`_{\mathrm{TR}}`$ resemble the Boltzmann integro-differential transport equation, which can better model rapidly varying, “ultra-high-energy”, macroscopic fluid phenomena than the differential equations of motion of fluid dynamics. (To this end it uses an infinite number of fields to take some account of the underlying microscopic behaviour.) So the Euler-Lagrange equations of $`_{\mathrm{TR}}`$ may be regarded as classical transport equations of motion for the one-particle distribution of some infinitesimal entities, such as X-ons surmised to underly all physical phenomena by Feynman. (ii) Ever since the EPR gedanken experiment, it is known that interpretations of certain quantum phenomena suggest the existence of causal faster-than-light effects. The Euler-Lagrange equations of $`_{\mathrm{TR}}`$ are the first Lorentz-invariant equations of motion that classicaly model such effects, because their retarded solutions have unbounded front velocities . Which is a major qualitative advantage of $`_{\mathrm{TR}}`$ over $`_{\mathrm{QED}}`$.
Using the Euler-Lagrange equations of $`_{\mathrm{TR}}`$ with $`e=0`$ and proceeding as in Ref. , we calculate the causal dependence of $`\mathrm{\Psi }_\mu (x,p)`$ and $`\mathrm{\Psi }_\mu ^{}(x,p)`$ on $`J_\mu (x)`$, and of $`\mathrm{\Psi }_{1/2}(x,p)`$ on $`J_e(x)`$. Thereby we can infer that the Feynman propagator for the four-vector field $`A_\mu (x)`$ defined by (19) equals
$$i(2\pi )^4\stackrel{~}{g}_1\frac{\delta _{\mu \nu }+\stackrel{~}{\mu }^2k_\mu k_\nu }{k^2+\stackrel{~}{\mu }^2},$$
(20)
$$\stackrel{~}{g}_1(k^2)q_1s_1^2I_{10}I_{20}^2,\stackrel{~}{\mu }(k^2)|s_1|^1I_{20}^1,$$
(21)
where $`I_{mn}(k^2)`$ is an analytic function of the complex variable $`k^2`$ such that
$$I_{mn}(k^2)=2\pi ^2\mathrm{\Lambda }^m_0^{\mathrm{}}y^{m+n}f^2(y)t^m(y)[\sqrt{1+\mathrm{\Lambda }^2k^2yt^2(y)}+1]^m𝑑y$$
(22)
for $`k^2>0`$; and the Feynman propagator for the bispinor field $`\psi (x)`$ defined by (19) equals
$$i(2\pi )^4\stackrel{~}{g}_{1/2}\frac{i\gamma ^\mu k_\mu +\stackrel{~}{m}}{k^2+\stackrel{~}{m}^2},$$
(23)
$$\stackrel{~}{g}_{1/2}(k^2)q_{1/2}s_{1/2}^1I_{10}I_{20}^1,\stackrel{~}{m}(k^2)s_{1/2}^1\{1s_{1/2}^2[I_{10}I_{11}+\frac{1}{4}k^2I_{20}^2]\}I_{20}^1;$$
(24)
where $`k^2`$ has to be replaced everywhere with $`k^2iϵ`$, by the Feynman prescription.
If functions $`t(p^2)`$ and $`f(p^2)`$ are such that
$$_0^{\mathrm{}}f^2(y)t(y)|\sqrt{y}/t(y)|^{l+1}𝑑y=0$$
(25)
for $`l=0,1,\mathrm{},n`$, then for complex values of $`k`$ as $`|k^2|\mathrm{}`$:
$$\left|\stackrel{~}{g}_1\frac{\delta _{\mu \nu }+\stackrel{~}{\mu }^2k_\mu k_\nu }{k^2+\stackrel{~}{\mu }^2}\right|=O(|k^2|^{(1n)/2}),\left|\stackrel{~}{g}_{1/2}\frac{\stackrel{~}{m}i\gamma ^\mu k_\mu }{k^2+\stackrel{~}{m}^2}\right|=O(|k^2|^{n/2}).$$
(26)
When the function $`y/t^2(y)`$ takes only a finite number of real values $`v_i`$, $`i=1,2,\mathrm{}`$, we can explicitly evalute integrals (22); we obtain
$$I_{mn}(k^2)=\mathrm{\Lambda }^m\underset{i}{}A_{mni}v_i^m[\sqrt{1+\mathrm{\Lambda }^2v_ik^2}+1]^m,$$
(27)
where $`A_{mni}`$ are real constants. Considering such a case, we can show that for any $`\mu ^2`$, $`m^2`$ and integer $`n`$, there exist functions $`f(p^2)`$ and $`t(p^2)`$, and constants $`s_1`$, $`s_{1/2}`$, $`q_1`$, $`q_{1/2}`$, and $`\mathrm{\Lambda }_0>0`$ such that the propagators (20) and (23) with $`\mathrm{\Lambda }>\mathrm{\Lambda }_0`$ are regularizations of spin 1 and spin $`\frac{1}{2}`$ propagators (7) such that: (i) they have properties analogous to those of propagator (1), and (ii) there is a positive constant $`k_0^2`$ such that for all $`k^2\mathrm{\Lambda }^2k_0^2`$ the functions $`I_{mn}(k^2)`$, $`\stackrel{~}{g}_{1/2}(k^2)`$, $`\stackrel{~}{\mu }(k^2)`$, $`\stackrel{~}{g}_1(k^2)`$, and $`\stackrel{~}{m}(k^2)`$ are real. In such a case: (i) The constants $`s_1`$, $`s_{1/2}`$, $`q_1`$, and $`q_{1/2}`$ are such that
$$\stackrel{~}{\mu }^2(\mu ^2)=\mu ^2,\stackrel{~}{m}^2(m^2)=m^2,$$
(28)
$$\stackrel{~}{g}_{1/2}(\mu ^2)=1+d\stackrel{~}{\mu }^2(k^2=\mu ^2)/dk^2,\stackrel{~}{g}_1(m^2)=1+d\stackrel{~}{m}^2(k^2=m^2)/dk^2.$$
(29)
So the propagators (20) and (23) have poles at $`k^2=\mu ^2`$ and $`k^2=m^2`$, where their behaviour is given by the spin 1 and spin $`\frac{1}{2}`$ propagators (7) with $`ϵ=0`$. (ii) The difference between spin 1 propagator and propagator (20) depends on the value of $`\mathrm{\Lambda }`$ so that it satisfies relations analogous to (2); and the same goes for spin $`\frac{1}{2}`$ propagators. (iii) The propagators (20) and (23) are analytic functions of $`k^2`$ that (a) are not continuous everywhere across the negative real axis, (b) have no additional singularities to those of spin 1 and spin $`\frac{1}{2}`$ propagators (7), and (c) satisfy relations (26). For any integer $`n3`$, their Källén-Lehman integral representations are superconvergent: in $`x`$-space we can decompose the Feynman propagators (20) and (23) into positive and negative energy parts without contact terms.
As a consequence of (i) and (ii) above, the classical, inhomogeneous Maxwell equations can be obtained from the Euler-Lagrange equations of $`_{\mathrm{TR}}`$ with $`J_\mu =0`$ and $`J_e=0`$ and the definitions (19) by limiting $`\mathrm{\Lambda }0`$.
## IV Realistic regularization of the QED Green functions
To obtain a perturbative S-matrix of quantum-electrodynamic phenomena based on the Lagrangian $`_{\mathrm{TR}}`$, say $`S_{\mathrm{TR}}`$, we use the ’t Hooft-Veltman definition of an S-matrix. In view of results of Sec.III, there are functions $`f(p^2)`$ and $`t(p^2)`$, and constants $`s_1`$, $`s_{1/2}`$, $`q_1`$, $`q_{1/2}`$, and $`\mathrm{\Lambda }`$ such that the n-point Green functions of $`_{\mathrm{TR}}`$ and the corresponding S-matrix $`S_{\mathrm{TR}}`$ have the following properties:
* As the Lagrangian $`_{\mathrm{TR}}`$ has the same interaction and source terms as the QED Lagrangian $`_{\mathrm{QED}}`$, they are expressed in terms of QED diagrams with the spin 1 and spin $`\frac{1}{2}`$ propagators (7) replaced with their regularizations (20) and (23), whereas the vertices are the same as in QED, i.e., $`(2\pi )^4\gamma _\mu `$; so all diagrams are finite.
* To any order in the fine structure constant the n-point Green functions are causal ; charge and total four-momentum conserving; Lorentz-invariant; and C-, P- and T-invariant up to a phase factor .
* If not only the propagators (20) and (23) but also the higher-order two-point Green functions of $`_{\mathrm{TR}}`$ have no additional singularities, then $`S_{\mathrm{TR}}`$ relates the same particles as the S-matrix of QED with massive photons in unitary gauge: electrons and positrons, each with two possible polarization vectors, and massive photons with three possible polarization vectors; none of them with wrong metric or statistics. As the propagators (20) and (23) admit the Källén-Lehman representation, this scattering matrix $`S_{\mathrm{TR}}`$ is unitary to any order in the fine structure constant.
* In the asymptote $`\mathrm{\Lambda }\mathrm{}`$, the propagators (20) and (23) behave as sufficies for renormalization.
So, the perturbative $`n`$-point Green functions of $`_{\mathrm{TR}}`$ are the result of a rp-theory as defined in Section I.
In view of (iv), we can compute by renormalization the renormalized n-point Green functions of QED with massless photons from the n-point Green functions of $`_{\mathrm{TR}}`$ by choosing an appropriate dependence of $`e`$, $`s_1`$, $`s_{1/2}`$, $`q_1`$, and $`q_{1/2}`$ on $`\mathrm{\Lambda }`$, and then limiting $`\mathrm{\Lambda }\mathrm{}`$ and the renormalized photon mass to zero .
## V Comments
Generalizing the ’t Hooft and Veltman method of unitary regulators we have shown, for the first time as far as we know, that one can regularize the QED Green functions in accordance with the basic tenets of theoretical physics by suitably modifying the free part of QED Lagrangian. As we mentioned in Sec.III, the physical motivation for such modification has been the Feynman surmise about X-ons, the Boltzmann improvement on fluid dynamics by the transport theory based on his equation, and interpretations of certain quantum-electrodynamic phenomena that suggest causal faster-than-light effects.
Within the framework of perturbative quantum field theory as defined by ’t Hooft and Veltman, the Lagrangian $`_{\mathrm{TR}}`$ is related to the physical world solely through the perturbatively defined scattering matrix $`S_{\mathrm{TR}}`$. We see no physical properties of $`S_{\mathrm{TR}}`$ that require the spectral function (4) and the Hamiltonians corresponding to free Lagrangians $`_1`$ and $`_{1/2}`$ (which are not free-particle Lagrangians) to be positive as they turn out to be within the framework of canonical formalism.
The need for a regularization of QED that would result in a realistic physical model was felt very strongly by the founders of QED, Dirac and Heisenberg, already some sixty-five years ago. But neither they nor their contemporaries succeded in getting rid of the ultraviolet divergencies by a physicaly motivated modification of the QED Lagrangian. In the late 1940s, however, Tomonaga, Schwinger and Feynman “solved” the problem of QED ultraviolet divergencies through renormalization—a solution which does not require the preceding regularization to be realistic, and removes all parameters characteristic of it. As they obtained spectacularly succesful formulas for quantum-electrodynamic phenomena, the problem of finding a realistic, Lagrangian-based regularization of the QED Green functions was not so urgent any more. As there had been no progress whatsoever towards a solution of this problem, it mainly came to be considered as practically unsolvable; those who hoped otherwise were often considered “irrational”, as Isham, Salam, and Strathdee complained twenty-five years later. Thus nowadays, as far as we know, no quantum-field theorist, excepting the string theorist, pays much attention to this problem, which many of the preceding generations—e.g., Dirac, Heisenberg, Landau, Pauli, and Salam, to mention some—still hoped to be solved somehow someday. But the string theorists abandon one of the basic premises of conventional physics, the four-dimensionality of space-time. We have shown, however, that such drastic steps may be avoided when modifying QED Lagrangian to get rid of ultraviolet divergencies. But the question remains which modification of the type considered is the most appropriate for better describing quantum-electrodynamic phenomena and their faster-than-light effects than the conventional QED, and what is the content of such a perturbative theory.
## VI acknowledgement
Authors greatly appreciate discussions with M. Poljšak and B. Bajc and their suggestions.
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# On Orientifolds, Discrete Torsion, Branes and M Theory
## 1 Introduction
Orientifold planes are objects in string theory which are defined perturbatively by gauging a discrete symmetry which involves reversing the sign of coordinates transverse to the plane while changing the orientation of the string . Before we delve into the details of their construction and their dynamics let us mention some motivation and interest in their study.
Orientifold planes are of interest in many aspects of string theory. They turn out to be useful in studying disconnected components in moduli space of various string compactifications. They are crucial in the brane construction of gauge theories with $`Sp/SO`$ gauge groups.<sup>1</sup><sup>1</sup>1See for a recent interesting application. In addition they give a simple realization of matter in symmetric and antisymmetric second rank representations in $`SU(N)`$ gauge theories.
As will be discussed in detail in the following sections, orientifold planes turn out to be characterized by discrete fluxes. Some of these fluxes originate from NS fields (See for a recent discussion) and some other fluxes come from RR fields. While the first (NS) are visible from the usual perturbative description of orientifolds, the latter (RR) cannot be treated in a perturbative formulation of the Type II string since, as of now, there is no formalism which includes RR backgrounds in string theory. For this reason we need to rely on other methods in order to study such objects, and once such a formalism is available, these questions should be revisited. The best analysis which is available at the moment will include essentially two tools \- various dualities, and discrete fluxes represented by intersections of branes with the orientifold plane.
Here we study two subjects related to orientifolds - their lift to M theory (See a recent related discussion in ), and the different variants they appear in. We begin with a review of the relevant features of orientifolds in section 2. We describe the 4 types of $`Op`$ planes which are known to exist for $`p5`$. We explain how these 4 types are classified by a pair of $`𝐙_2`$ parameters, one perturbative coming from the $`B_{NS}`$ form and the other from one of the RR forms. These parameters are “discrete torsion” parameters which arise whenever the field strength of a p-form potential has a corresponding non trivial discrete cohomology in space time (see also recent related work ). Equivalently we show that these variants arise as brane configurations in the spirit of (see recent related work ) from intersections of an O plane with a brane which divides it into two parts. We search for all the possible discrete torsion variants, and we find that for $`p1`$ there are some additional ones. Later in that section, in 2.2, we explain how T duality relates a wrapped orientifold with a pair of lower dimensional orientifolds and the mapping between the various discrete torsion variants.
In section 3 we systematically review the $`𝐙_2`$ orbifolds of M theory. There are only a few of these and upon compactification they give rise to the various $`𝐙_2`$ orientifolds and orbifolds of string theory. We review the M lifts of the O0, O2, O4 in terms of 11d O planes - the OM1, OM2 and OM5. In the case of the $`O2`$ we give a new M theory explanation for the difference between the two variants $`O2^+`$ and $`\stackrel{~}{O2}^+`$.
In section 4 we describe the M lift of the O1, O3 and O5 in terms of the same OM planes. In all cases we pay special attention to the $`SL(2,𝐙)`$symmetry (see also related work ). For the O3 we show how the M lift realizes the $`SL(2,𝐙)`$geometrically, as usual. For the O1 and O5 the S transformation of $`SL(2,𝐙)`$transforms the orientifold planes into orbifolds, which originate in the same OM planes.
In section 5 we discuss in detail the interrelations between $`Op`$ planes and orbifolds for $`p=1,5`$. We explore the possible discrete torsion variants and discuss new ones for the orbifold and orientifold lines and for the orientifold point. We compare the discrete torsion classification with a recent claim on a K theory classification (see for earlier and related work) and find groups of the same order but with a different group structure. We are able to find the tension of some of the variants, and it remains to do so for the other cases.
Finally we discuss some miscellaneous applications in section 6. We discuss the dyon spectrum for 4d $`𝒩=4`$ with $`SO/Sp`$ gauge groups, where the $`Sp`$ theory has two versions. Then the discussion on the monopole spectrum of such a theory is generalized to a field theory in an arbitrary dimension. Finally, we discuss the problems associated with a $`(p,q)`$web of orientifolds.
Let us mention some open questions
* We show how an intersection of an orientifold, or an OM plane, with a brane causes a tension jump, and so a fractional charge is deposited on the brane. This jump should be understood and probably required by the worldvolume theory on the brane. In particular consistency with Dirac quantization is required.
* We find the tension of some of the discrete torsion variants of the orbifold lines, and it would be interesting to know all of them.
Preliminary results of this work were presented in Santa Barbara and can be found in the following link: http://online.itp.ucsb.edu/online/susy99/hanany/.
## 2 <br>Introduction to Orientifolds
An orientifold plane in $`p+1`$ dimensions is defined as Type II string theory on $`R^{p,1}\times R^{9p}/(I\mathrm{\Omega }J)`$ where $`I`$ is the inversion of all coordinates in the transverse space $`R^{9p}`$, $`\mathrm{\Omega }`$ is the orientation reversal on the world sheet of the fundamental string and $`J`$ is the identity operator for $`p=0,1\mathrm{mod}4`$ and $`(1)^{F_L}`$, the left moving spacetime fermion number operator, for $`p=2,3\mathrm{mod}4`$ . We will denote an orientifold p-plane as an “$`Op`$ plane”.
By definition the orientifold acts on $`B_{NS}`$ (with components parallel to the $`Op`$ plane) with a sign reversal. The action on the other NS fields is trivial, so let us specify the action on all of the RR forms. The $`O9`$ which projects Type IIB to Type I keeps invariant the left-right symmetric part of the spectrum, which in the RR sector is the 2-form $`B_{RR}`$. By T dualizing we arrive at the following rule for the action of an $`Op`$ plane on the RR forms $`C_p^{}`$
$`C_p^{}C_p^{}p^{}=p+1\mathrm{mod}4`$
$`C_p^{}C_p^{}p^{}=p+3\mathrm{mod}4.`$ (1)
Note that the above sign comes in addition to a component dependent sign which comes from the tensor transformation rules. Forms which get reversed by the action of the group are termed “twisted”, and this notation should not be confused with forms from a twisted sector.
The transverse space to an $`Op`$ plane is an $`\mathrm{𝐑𝐏}^{\mathrm{𝟖}𝐩}`$. This space has some (discrete) torsion cohomologies which are summarized in appendix A. Whenever a space has discrete torsion, forms of an appropriate rank are topologically classified by it. Namely
$$[G_{p+1}]H^{p+1}(X,𝐙)$$
(2)
where $`H^{p+1}(X,Z)`$ is the integral cohomology of the space $`X`$ and $`G_{p+1}`$ is the RR field strength $`p+1`$-form. Twisted forms require the use of twisted cohomologies, $`\stackrel{~}{H}^{p+1}`$. In the case of orientifolds we will see that discrete torsion allows for the orientifold to appear in several variants.
For a general orientifolds ($`p5`$) there are at least four types of orientifold planes which are distinguished by two $`𝐙_2`$ charges. One $`𝐙_2`$ charge comes from the NS sector and the other comes from the RR sector. For this reason, the former charge is seen in perturbative string theory while the latter is not. The perturbative distinction was first discussed in modern language in . Later, the other $`𝐙_2`$ charge was discovered in <sup>2</sup><sup>2</sup>2In there is an earlier discussion of the four variants of 5 planes in an M theory language which is directly related to discrete torsions. for the case $`p=3`$ essentially by using field theory intuition which is based on Montonen Olive Duality or Type IIB self duality. The case $`p=4`$ was studied in detail by .
For an $`Op`$ plane, the two $`𝐙_2`$ charges are given by $`b`$, which is the class $`[dB_{NS}]=[H_{NS}]\stackrel{~}{H}^3(\mathrm{𝐑𝐏}^{\mathrm{𝟖}𝐩})`$ and by $`c`$, which is the class of an RR form $`[dC_{5p}]=[G_{6p}]𝐙_2`$ (the relevant cohomology, either twisted or untwisted, is found to be $`𝐙_2`$). In the next subsection we will give these a physical interpretation via brane-orientifold intersections. The above classes can be defined for $`p5`$, while for $`6p8`$ one can still define $`b`$ through the perturbative action on open strings. In general the existence of a discrete torsion implies the existence of variants, though variants could possibly exist without it, as in the case of the $`b`$ variant. The absence of a $`c`$ charge for $`p=6`$ suggests that, contrary to naive expectations, one can not have an $`\stackrel{~}{O6}^{}`$ plane, i.e. half a D6 brane stuck on an $`O6^{}`$ plane. That turns out to be the case as implied by the results of .
One can systematically check for any $`p`$ which additional forms can get discrete torsions. For low $`p`$ some extra variants exist (see section 5): for $`p1`$ one can define an additional charge $`c^{}`$ by $`[dC_{1p}]=[G_{2p}]𝐙_2`$, and for $`p=1`$ one has $`H_{NS}\stackrel{~}{H}^7(\mathrm{𝐑𝐏}^\mathrm{𝟕})=𝐙_\mathrm{𝟐}`$. In addition there are two other discrete torsion whose variants we ignore - $`H^0(\mathrm{𝐑𝐏}^𝐧)=𝐙`$ leads to a integral discrete variant for $`p=2\mathrm{mod}4`$ which we interpret as describing the possible massive theories of type IIA (this interpretation is based on the brane intersection picture, where this torsion is seen to be due to an intersection with a D8 plane). The other type which we ignore is $`H^{8p}=𝐙`$ for odd $`p`$ and $`\stackrel{~}{H}^{8p}=𝐙`$ for even $`p`$, which simply means that one can add to an $`Op`$ plane any integer number of $`Dp`$ branes.
The orientifold planes will be denoted according to their $`𝐙_2`$ charges. A trivial $`b`$ charge will be denoted by a <sup>-</sup> superscript and a non-trivial $`b`$ will be denoted by <sup>+</sup>. A non-trivial $`c`$ charge will be denoted by adding a $`\stackrel{~}{}`$. The charges of these planes are $`2^{p5}`$ for the $`Op^{}`$ plane, $`+2^{p5}`$ for both the $`Op^+`$ and the $`\stackrel{~}{Op^+}`$ and $`\frac{1}{2}2^{p5}`$ for the $`\stackrel{~}{Op^{}}`$ plane. The tensions of these objects are measured in units of the Dp brane tension and are identical to their charge. Note a change of notation for the various orientifold planes from the papers . The orientifold planes are denoted there $`Op^{},Op^+,\widehat{Op},`$ and $`\stackrel{~}{Op}`$, respectively. The new notation is based on the two $`𝐙_2`$ charges and makes it more simple to work with. $`n`$ physical $`Dp`$ branes stacked upon an $`Op`$ plane leads to a gauge group $`SO(2n),Sp(n),SO(2n+1)`$ and $`Sp(n)`$ for $`Op^{},Op^+,\stackrel{~}{Op^{}},\stackrel{~}{Op^+}`$, respectively. The two $`Sp(n)`$ theories differ, in some cases, by a theta angle and in their monopole spectrum. This will be discussed in detail in sections 6.1 and 6.2.
### 2.1 Brane Realization of Discrete Torsion
In this section we will use a basic fact about orientifold planes and NS five branes and will develop a set of relations between intersecting branes and orientifold planes. The basic observation, made in , can be described in figure 1. Consider a NS five brane which spans the 012345 coordinates and intersecting an $`Op`$ plane which spans the $`0,\mathrm{},p1`$, and 6 coordinates, $`p6`$. It splits the orientifold plane into two different parts, one to its left and one to its right and the orientifold plane changes its type as in figure 1.
The NS brane in this configuration has the special property that it is reflected onto itself by the orientifold. In a compact configuration such a brane would have half of the charge of a brane which does not intersect the orientifold and so we call it a $`\frac{1}{2}`$ NS brane. Note that a 1/2 brane which extends outside of the orientifold has the same charge density as a unit brane, only it covers half of the volume due to the projection. On the other hand, a half brane inside the orientifold has half the charge density. A 1/2 brane cannot move in directions transverse to the orientifold plane, because only integral branes can move away off the orientifold. In the following brane configurations one should bear in mind that all branes which intersect with an orientifold plane are 1/2 branes.
Let us emphasize this point in more detail. Consider a configuration in which a physical NS brane is located far away from the orientifold, as in figure 2a. Let this brane move slowly towards the orientifold plane. It moves together with its image under the orientifold projection. As the two images meet on the orientifold plane, they can split along it as in figure 2b. At this point the rule in figure 1 can be used and a new type of orientifold plane is generated in between the two half NS branes.
We can look at it from a different point of view starting from figure 2b and moving to figure 2a. As observed in , the type of the orientifold plane changes as one crosses the $`\frac{1}{2}`$ NS brane. This is determined by the $`b`$ charge of the orientifold which changes as one crosses the $`\frac{1}{2}`$ NS brane. On the other hand, if we add an additional $`\frac{1}{2}`$ NS brane, the type of the orientifold changes back to its original value, as in figure 2b. In this case we can have a dynamical process in which the two half NS branes combine together along the 6 coordinate and leave the orientifold along the 789 directions as in figure 2a. This is the inverse process to the one described in the last paragraph.
One can relate such configurations to the value of the two form flux $`b`$ by the following reasoning. The combined configuration of a $`Dp`$ brane and half an NS brane is located at a point in 789 directions. There is a corresponding $`\mathrm{𝐑𝐏}^\mathrm{𝟐}`$ which surrounds the configuration. The field which couples magnetically to the NS brane is the 2-form NS field. Consequently, the integral of the two form over the $`\mathrm{𝐑𝐏}^\mathrm{𝟐}`$ measures the number of 1/2 NS branes, mod 2, which are located within the $`\mathrm{𝐑𝐏}^\mathrm{𝟐}`$ <sup>3</sup><sup>3</sup>3A discussion in a similar spirit can be found in ..
$$\mathrm{exp}\left(i_{\mathrm{𝐑𝐏}^\mathrm{𝟐}}B_{NS}\right)=()^{\mathrm{\#}\frac{1}{2}\mathrm{NS}\mathrm{branes}}.$$
(3)
This holonomy (“Wilson loop”) $`H^p(X,U(1))`$ of a general p-form potential has the same discrete torsion as the field strength class $`H^{p+1}(X,Z)`$ which was mentioned before. This is seen via the exact sequences $`0𝐙𝐑U(1)0`$ and
$`\mathrm{}`$ $`H^p(X,𝐑)`$ $`H^p(X,U(1))`$
$`H^{p+1}(X,𝐙)`$ $`H^{p+1}(X,𝐑)`$ $`\mathrm{}`$ (4)
where the cohomologies with real coefficients do not have torsion parts.
We will assume that the configuration of figure 1 exists and will apply some S and T dualities to see what other results can be obtained from this basic configuration. First let us apply S duality to this configuration in the case $`p=3`$ (see also ). For this we need to know what are the S transformations of the various orientifold planes. This is easily done by looking at their $`𝐙_2`$ charges which are the NS 3-form and the RR 3-form field strengths. Let us recall that $`O3^{}`$ has charge $`(b,c)=(0,0)`$ and the $`\stackrel{~}{O3^+}`$ has charge (1,1). These orientifolds are self-dual under S duality. $`O3^+`$ has (1,0) and transforms under S duality to (0,1) which is the $`\stackrel{~}{O3^{}}`$.
Equipped with this data let us make an S duality on the configuration of figure 1. One gets that an $`O3^{}`$ transforms to an $`\stackrel{~}{O3^{}}`$ while crossing a half D5 brane. By crossing another half NS brane the $`b`$ charge of the orientifold jumps and one gets $`\stackrel{~}{O3^+}`$. This is expected if one assumes that the configuration of half D5 and half NS brane are self dual under S duality and that the order of the D5 and NS brane is not important. The summary is that when crossing a half NS brane the $`b`$ charge changes while when crossing a half D5 brane the $`c`$ charge of the orientifold changes. $`b`$ measures the number of half NS branes enclosed by an $`\mathrm{𝐑𝐏}^\mathrm{𝟐}`$ surrounding the intersection of O3 plane and NS brane while $`c`$ measures the number of half D5 branes enclosed by an $`\mathrm{𝐑𝐏}^\mathrm{𝟐}`$ surrounding the intersection of O3 plane and D5 brane.
Next we can perform T duality on these results. Consider dualizing this system in a supersymmetric fashion to a system of a NS brane along 012345 and a D6 brane along 0123789. The orientifold plane dualizes to an O4 plane along 01236. The two $`𝐙_2`$ charges are measured by $`b`$ and by $`c`$ which is a Wilson line of the RR one form of Type IIA. T duality then implies that when an O4 plane crosses a half NS brane its $`b`$ charge changes while when crossing a half D6 brane its $`c`$ charge changes (A Wilson line $`c`$ can be associated to our configuration by repeating the arguments for the O3 plane). The system is located at a point in the 45 directions. The object which couples magnetically to the D6 brane is the RR 1-form of Type IIA. Correspondingly $`c`$ measures the number of half D6 branes, mod 2, trapped inside the $`\mathrm{𝐑𝐏}^\mathrm{𝟏}`$.
Let us summarize the situation after further applying T duality to other directions. An $`Op`$ plane along 012,…,$`p1`$ and 6, a half NS brane along 012345 and a half D(p+2) brane along 012,…,$`p1`$, 789. When crossing a half NS brane the $`Op`$ plane changes its $`b`$ charge and when crossing a half D(p+2) brane the $`Op`$ plane changes its $`c`$ charge which is measured by a ‘Wilson loop integral’ of the RR $`(5p)`$ form potential. This configuration exists for any $`p`$ which is 5 or less. $`b`$ measures the number of half NS branes in an $`\mathrm{𝐑𝐏}^\mathrm{𝟐}`$ enclosing the configuration while $`c`$ measures the number of half $`D(p+2)`$ branes inside an $`\mathrm{𝐑𝐏}^{\mathrm{𝟓}𝐩}`$ which surrounds the configuration.
### 2.2 T Duality
Let us apply T-duality in a direction along the orientifold plane. An $`O(p+1)`$ plane wrapping a circle of radius $`L`$ turns after T duality to a pair of $`Op`$ planes on a circle of radius $`L^{}=1/L`$. Since we have (at least) 4 possible types for each $`Op`$ plane we expect 16 possible types for the wrapped $`O(p+1)`$ plane.
Although this may sound surprising at first, we will demonstrate the 16 possibilities by analyzing the possible discrete torsions. The transverse space is $`\mathrm{𝐑𝐏}^{\mathrm{𝟕}𝐩}`$ with cohomologies given in appendix A, and we should find the 9 dimensional fields that can have discrete torsions. Before wrapping the $`O(p+1)`$ plane the fields with $`𝐙_2`$ discrete torsions are $`B_{NS}`$ and $`C_{5(p+1)=4p}`$. After compactifying on $`L`$ there are two more such fields \- the reduction of the metric on $`L`$ (an untwisted 1 form), and the reduction of $`C_{6p}`$ on $`L`$. To find the relation with the T dual picture it is useful to list the various discrete torsions in both pictures
| $`O(p+1)`$ | $`B_{NS}`$ | $`C_{4p}`$ | $`g_{\mu \nu }/L`$ | $`C_{6p}/L`$ |
| --- | --- | --- | --- | --- |
| $`O(p)`$ | $`B_{NS}`$ | $`C_{5p}/L^{}`$ | $`B_{NS}/L^{}`$ | $`C_{5p}`$ |
(5)
where the notation $`C_q/L`$ means “the form $`C_q`$ reduced on the circle $`L`$”. Let us restate this mapping in terms of brane intersections
| $`O(p+1)`$ | NS5 | $`D(p+3)`$ | KK monopole | $`\widehat{D}(p+1)`$ |
| --- | --- | --- | --- | --- |
| $`O(p)`$ | NS5 | $`\widehat{D}(p+2)`$ | $`\widehat{NS}5`$ | $`D(p+2)`$ |
(6)
where the hat above a brane means that it does not wrap the circle.
Let us discuss some examples. An NS5 intersecting the $`O(p+1)`$ turns into an NS5 which intersects both of the $`Op`$’s while an intersection with a $`D(p+3)`$ turns into a $`D(p+2)`$ which intersects only one out of the pair. If we want to get an RR discrete torsion on both $`Op`$’s then we should intersect the $`O(p+1)`$ with a $`D(p+1)`$, one which does not wrap the circle (and still has 4 mixed directions relative to the O plane). If we take an $`Op^{},Op^+`$ pair, which has zero charge, then it transforms into an $`O(p+1)`$ intersecting a KK monopole on the circle $`L^{}`$ and must have zero charge as well.
## 3 Review of M-lifts of Orientifolds
In this section we review known lifts of orientifold planes to M theory. Such a lift requires lifting the $`𝐙_2`$ action to M theory. The objects we describe are lifts of both orientifolds and orbifolds and accordingly are denoted as $`OMp`$ planes. Since the worldsheet formulation is lost in the lift, orientation reversal is meaningless. Nevertheless we shall sometimes continue to call them “orientifolds”.
Lacking a fundamental definition of M theory, we are satisfied by specifying the $`𝐙_2`$ action on the 11d supergravity fields. An $`OMp`$ plane includes a transverse spatial reflection, so we look at M theory on $`𝐑^{p,1}\times 𝐑^{10p}/𝐙_2`$ where the first factor is the worldvolume of the $`OMp`$ plane, the $`𝐙_2`$ in the second part is the reflection, and often it will be more convenient to replace the last factor with $`𝐓^{10p}/𝐙_2`$. The reflection determines the action on the metric ($`gg`$). The action on the 3-form $`C`$ is determined by requiring invariance of the topological term in the action $`CGG`$, where $`G=dC`$, to be
$$C()^pC.$$
(7)
Supersymmetry provides another constraint. When acting on fermions, the inversion of $`10p`$ coordinates squares to the identity for $`10p=0,1\mathrm{mod}4`$ and to $`()^F`$ for $`10p=2,3\mathrm{mod}4`$ (this is a consequence of $`(\mathrm{\Gamma }_1\mathrm{\Gamma }_2\mathrm{}\mathrm{\Gamma }_n)^2=()^{n(n1)/2}`$). Thus in order to have a supersymmetric orbifold and a $`𝐙_2`$ action rather than a $`𝐙_4`$ action we require
$$p=1,2\mathrm{mod}4.$$
(8)
We see that the $`𝐙_2`$ objects intrinsic to M theory, are the OM1, OM2, OM5, OM6 and OM9. It is no surprise that the orientifolds intrinsic to M theory include a 2 plane and a 5 plane which we denote by OM2 and OM5. We shall review the definition of these orientifolds and the way they were used to find the M lift of the O4 and the O2. Then we look for intersections of branes and orientifolds in M theory, where it turns out that there is only one such configuration: a $`1/2`$ M5 stuck on an OM2. After that we review the lift of another Type IIA orientifold, the O0, through an M theory object which we may call OM1.
We summarize here the results. The OM5 and OM2 carry charges given by
| $`OM2^{}`$ | $`OM2^+`$ | OM5 |
| --- | --- | --- |
| -1/16 | +3/16 | -1/2 |
(9)
where the $`OM2^{},OM2^+`$ are two discrete torsion variants of the $`OM2`$ plane (which will be discussed later in subsection 3.2.1), and the charge is in units of bulk OM2 and OM5 branes. The other planes support a “twisted sector” matter rather than charge (they are neutral in 11d)
(12)
The OM6 and OM9 will not be discussed any further in this paper. The OM6 is M theory on $`𝐓^\mathrm{𝟒}/𝐙_\mathrm{𝟐}`$ and the OM9 is a Hořava - Witten plane .
### 3.1 The OM5 plane
The “orientifold” M on $`𝐓^5/𝐙_2`$ was studied in . The “untwisted” sector has a 6d gravitational anomaly that can be canceled by 16 tensor multiplets. Moreover, local anomaly cancellation would seem to require adding a twisted sector of 1/2 a tensor multiplet at each of the 32 fixed points. This problem is avoided by using a different method of canceling the anomaly - assigning to each one of the fixed points a charge
$$Q(OM5)=1/2,$$
(13)
in units of physical M5 branes. Note that while we can usually put a half brane on top of an orientifold plane, one cannot put half an M5 on an OM5 due to M theory flux quantization .
#### 3.1.1 M-lift of the O4
The OM5 can account for the O4 planes of Type IIA (see also ). When lifting the $`𝐙_2`$ action from Type IIA to M theory it is required to specify the action on the circle $`R_{11}`$. Topologically there are three possible actions: the identity, a reversal and a shift through half a circle. The identity is interpreted as an OM5 wrapped on $`R_{11}`$. Reversal is not allowed since it would accumulate to an inversion in 6 directions which would break supersymmetry (8). A shift is the same as a non-trivial discrete torsion for the Type IIA 1-form $`A`$. It is exactly the discrete torsion which is present in the $`\stackrel{~}{O4}^\pm `$ planes. So “untilded” ($`c=0`$) orientifolds M-lift to $`𝐑^5/𝐙_2`$ while the “tilded” ones ($`c=1`$) M-lift to $`(𝐑^5\times 𝐒^\mathrm{𝟏})/𝐙_\mathrm{𝟐}`$ with the $`𝐙_2`$ acting on the circle by a shift.
Let us look more at the shift orbifold $`(𝐑^5\times 𝐒^\mathrm{𝟏})/𝐙_\mathrm{𝟐}`$. It has two kinds of minimal 1-cycles, one of them wraps $`𝐒^\mathrm{𝟏}`$ (“the circle”) and the other is a straight line between points identified by the $`𝐙_2`$ action (“antipodal line”). For points away from the origin of $`𝐑^5`$ the circle is smaller, but as we approach the origin the antipodal line becomes smaller, and at the origin itself it is half the size of the circle.
We summarize the different O4 planes and their M theory origin
1. $`O4^{}`$. Charge $`Q=1/2`$. This is simply an OM5 wrapping $`R_{11}`$.
2. $`\stackrel{~}{O4^{}}`$. Charge $`Q=0`$. This is the smooth shift orbifold, and as such indeed does not carry charge. (Recall that a bound state of an OM5 and a half M5 is not allowed).
3. $`O4^+`$. $`Q=+1/2`$. This is an OM5 with a full stuck M5. (It is stuck by imposing as Wilson loop an element of $`O(2)`$ which is not connected to the identity ).
4. $`\stackrel{~}{O4^+}`$. Charge $`Q=+1/2`$. This is the shift orbifold with a stuck M5 at the origin on the circle of half radius.
The two types of $`\stackrel{~}{O4}`$ planes were conjectured to be related to the two elements of $`\pi _4(Sp(n))=𝐙_2`$.
### 3.2 The OM2 plane
The OM2 orientifold was studied in (see also ). The charge can be found by considering the interaction $`CI_8(R)`$ for M theory on $`X=𝐓^\mathrm{𝟖}/𝐙_\mathrm{𝟐}`$. The effective charge is $`_XI_8(R)=\chi /24`$. Although this space is singular, we can define its “resolved cohomologies” by adding to the invariant (untwisted) cohomologies an extra 256 cohomologies (in $`H^{2,2}`$) from the RR twisted sector. This totals the Euler characteristic to 384 and so the total charge to -16. Dividing by 256, the number of fixed points, we find $`Q=1/16`$ in units of M2 charge.
This orientifold allows a variant due to a discrete flux. The transverse space is $`\mathrm{𝐑𝐏}^7`$ and the only field strength form in M theory is the 4-form $`G`$, so we are interested in the cohomology $`H^4(\mathrm{𝐑𝐏}^7,𝐙)=𝐙_2`$ (Appendix A). The discrete torsion adds a charge of a $`+1/4=\frac{1}{2}_{\mathrm{𝐑𝐏}^7}\frac{C}{2\pi }\frac{G}{2\pi }`$ in units of M2. We denote the OM2 with trivial discrete torsion by $`OM2^{}`$ and the one with non-trivial torsion by $`OM2^+`$. Their charges are summarized by
$`Q(OM2^{})=1/16`$
$`Q(OM2^+)=Q(OM2^{})+1/4=+3/16.`$ (14)
#### 3.2.1 The brane - orientifold intersection in M theory
The $`OM2^+`$can be realized by branes. To do that we put a half M5 brane on an $`OM2^{}`$plane. This is an analogue (the only one) of the Type II configurations in section (2).
Note that like the case of the OM5, one cannot attach a half M2 brane on top of an OM2 due to M theory flux quantization .
Once we compactify M theory we get (at least) two more possibilities for brane intersections. One possible configuration comes from lifting OF1 and D6 with charge jump of 1/16. It lifts to a wound OM2 intersecting a KK6 - a Kaluza Klein monopole (section 5.1). Another configuration is the lift of an OF1 intersecting a D2, that is a wound OM2 intersecting with a transverse M2 with a tension jump of $`+1`$ (sections 4.3, 5.1). An inspection of this example shows that the jump of one unit of F1 charge is actually represented by a physical wound M2 brane which is stretched, with its mirror, like a T shaped brane (this is the M-lift of a double F1 ending on a D2 ).
#### 3.2.2 M-lift of the O2
Once one identifies the M theory objects, the $`OM2^{}`$and the $`OM2^+`$, one can go on and find the M-lift of the O2 planes in Type IIA . Actually, originally it must have been easier to go in the opposite direction and determine the charges of the OM planes from the O2 planes. Like the case of M-lifting the O4, we need to lift the $`𝐙_2`$ action to $`R_{11}`$. This time it must be a reversal - it cannot be the identity because of the susy constraint on $`p`$ (8). So in M theory there are actually two fixed planes located at the two fixed points on the circle.
The different O2 planes and their M theory origin are described by
1. $`O2^{}`$. Charge $`Q=1/8`$. This is a pair of $`OM2^{}`$planes: $`1/8=2\times 1/16`$.
2. $`\stackrel{~}{O2^{}}`$. Charge $`Q=1/8+1/2=+3/8`$. This is a pair of $`OM2^+`$planes: $`3/8=2\times 3/16`$.
3. $`O2^+,\stackrel{~}{O2^+}`$. Both have $`Q=+1/8`$. This is a composite pair of an $`OM2^{}`$with an $`OM2^+`$, $`1/8=1/16+3/16`$. The two possible O2’s correspond to the possible ordering of the OM2’s. This can be seen by intersecting an $`O2^+`$ with a D4 brane. After lifting to M theory and using the intersection rule explained in section (2), one finds that the $`\stackrel{~}{O2^+}`$ has the reversed order of OM2’s.
### 3.3 The OM1 line
Let us consider the M-lift of the orientifold point, the O0, of Type IIA. We take the action on $`R_{11}`$ to be the identity (a reversal is not allowed by equation 8), so we consider M theory on $`R_{11}\times 𝐒^\mathrm{𝟏}\times 𝐑^\mathrm{𝟗}/𝐙_\mathrm{𝟐}`$. A computation of 2d Gravitational anomalies for M theory on $`𝐓^\mathrm{𝟗}/𝐙_\mathrm{𝟐}`$ suggests that there is a chiral fermion on every fixed line as an “M theory twisted sector” . An independent evidence for the existence of a chiral fermion on the fixed line comes from the computation of the Witten index of $`Sp(N),SO(N)`$ matrix quantum mechanics . We may call this line an OM1 orientifold.
It is interesting to get the “twisted sector” which is described above truly from a twisted sector of string theory (see also ). In order to get a 2d model one needs to compactify on an 8 manifold, and Type IIB on $`𝐓^\mathrm{𝟖}/𝐙_\mathrm{𝟐}`$ has the right action on the fields (this is the orbifold which we will call $`OP1_B`$ in section 5.1). Type IIA divided by $`I_8()^{F_L}`$ would also do, but we will stay with the more geometric example. The $`OP1_B`$ has a twisted sector from the 4 form wrapping the 256 resolved $`H^{2,2}`$ cohomologies ( describes the computation of twisted sectors in general). These scalars are chiral as they inherit their self duality from the 4 form. Thus we get one chiral scalar for each of the 256 fixed lines. Comparing with the 512 chiral fermions of M theory on $`𝐓^\mathrm{𝟗}/𝐙_\mathrm{𝟐}`$ we see that they could match by bosonization provided the periodicity of the scalars is at the free fermionic value, as it should. This actually gives a nice realization of 2 dimensional bosonization as implied by a lift to M theory.
As we consider here the OM1 on the $`R_{11}`$ circle there are two possible boundary conditions for the fermion. The Neveu-Schwarz boundary conditions correspond to an O0 with trivial $`RR`$ discrete torsion ($`c_{RR}=0`$), and Ramond corresponds to $`c_{RR}=1`$. This is verified by a computation of the Casimir energy in the two cases, which matches the O0 mass
$$M(O0)=\pm 1/32$$
(15)
where the units are of momentum quanta along the circle.
Let us summarize the different O0 planes and their M theory origin
1. $`O0^{}`$. Charge $`Q=1/32`$. This is an OM1 with NS boundary conditions, and with integral momentum.
2. $`\stackrel{~}{O0^{}}`$. Charge $`Q=1/32+1/2`$. This OM1 has NS boundary conditions, but carries half-integer momentum.
3. $`O0^+,\stackrel{~}{O0^+}`$. Both have $`Q=+1/32`$. These OM1’s are in the Ramond sector, and there are two of them due to the zero mode which generates a degenerate ground state.
## 4 M lifts of Type IIB Orientifolds
In this section we describe (new) M-lifts of various orientifold planes in Type IIB, while paying special attention to the transformation properties under $`SL(2,𝐙)`$.
We start by M-lifting the O3 plane, where we get a nice geometric/ microscopic realization of the $`SL(2,𝐙)`$symmetry. Then we discuss the O1 and O5 planes. Their S duals, which we call OF1 and ON5, are constructed, and will be further discussed in section 5. The ON5 was already discussed in some works (see for a recent review on this plane), and the OF1 was discussed in .
The method is to recall the M theory origin of a Dp brane in terms of M branes and then to find the analogous construction of an Op plane in terms of the OM planes which were reviewed in the previous section. The basic correspondence is between M theory on $`𝐓_𝐌^\mathrm{𝟐}`$ and IIB on a circle of radius $`L_{IIB}`$, so we should always distinguish two cases according to whether the Dp brane wraps $`L_{IIB}`$or not. At weak coupling the $`𝐓_M^2`$has a short side and a long side, such that their ratio is the string coupling (when the RR axion vanishes).
### 4.1 O3
Let us recall the M-lift of the D3 brane. A D3 which does not wrap $`L_{IIB}`$is an M5 wrapping the torus, while a D3 which wraps $`L_{IIB}`$is the M2.
The four kinds of O3 planes have the following charges (in D3 units)
$`Q(O3^{})=1/4,`$
$`Q(\stackrel{~}{O3^{}})=Q(O3^+)=Q(\stackrel{~}{O3^+})=+1/4.`$ (16)
We start by lifting an O3 which wraps $`L_{IIB}`$, for otherwise the circle $`L_{IIB}`$is inverted as well and we get two O3 planes at the two fixed points. Since the M-lift of a D3 brane which wraps $`L_{IIB}`$is an M2 which does not wrap $`𝐓_M^2`$, we should take an OM2 which does not wrap $`𝐓_M^2`$. Because of the compactness of $`𝐓_M^2`$we are considering actually four OM2 planes. The simplest possibility is to take four $`OM2^{}`$planes. One checks that the total charge $`4\times (1/16)=1/4`$ fits the $`O3^{}`$ plane as expected.
To get the other O3 planes <sup>4</sup><sup>4</sup>4A.H. would like to thank Jacques Distler for discussions on related points. we use the brane intersection picture of section 2. For example, to get the M-lift of the $`O3^+`$ we should intersect the $`O3^{}`$ with an NS5 brane. The lift of the NS5 is an M5 wrapping the long side. Using the basic intersection in M theory (section 3.2.1) we see that we end up with two $`OM2^+`$planes along the long side, and two other $`OM2^{}`$planes, as in figure 3. The charges fit since $`2\times (1/16)+2\times (+3/16)=+1/4`$. A similar argument works for the $`\stackrel{~}{O3^{}}`$ and the $`\stackrel{~}{O3^+}`$ planes by replacing the NS5 branes with a D5 brane or a (1,1) brane respectively. Figure (3) summarizes the various configurations.
Recall that the $`SL(2,𝐙)`$properties of O3 planes can be described in a diagram such as figure 3 , where $`SL(2,𝐙)`$acts on the torus in the diagram according the its natural action on $`(𝐙_2)^2`$. This action is clearly visible from our M-lift into $`𝐓_M^2`$.
There is an alternative way of finding the M lift of an O3 which uses T duality (section 2.2). Under T duality an O3 that wraps $`L_{IIB}`$turns into a pair of O2 planes, each one of which can be lifted to a pair of OM2 planes, as in section 3.2.2, giving 4 OM2 planes as above. In this way one can recover the different lifts for the different variants.
One may wonder about other choices for the signs of the four OM2 planes. The ones with an odd number of signs cannot be constructed by intersecting the $`O3^{}`$ with a Type IIB 5-brane as the 5-brane must intersect exactly two OM2 planes (because the M5 is oriented). Nevertheless, such configurations can be constructed making use of the large but compact circle $`L_{IIB}`$, by intersecting the wound O3 (wound on $`L_{IIB}`$) with a KK monopole (on $`L_{IIB}`$) as in (6). One gets a configuration with three $`OM2^{}`$ and one $`OM2^+`$ and total charge zero. Intersecting now with an NS5 would give the other possibility - one $`OM2^{}`$ and three $`OM2^+`$. The case of 4 $`OM2^+`$planes is probably equivalent to four $`OM2^{}`$planes with an additional M2 brane.
Now we turn to an O3 which does not wrap $`L_{IIB}`$. It is actually a pair of O3 planes, and we may T- Dualize them into an O4 plane as in section (2.2). The M-lift of the latter was described already in section 3.1.1 in terms of OM5 planes.
### 4.2 O5 - ON5
Let us discuss O5 planes which wrap $`L_{IIB}`$, both because we are less interested in a pair of O5’s which we would have had if the compact direction were inverted, and since we are interested in configurations which lift to M5 branes rather than KK monopoles. <sup>5</sup><sup>5</sup>5An alternative M lift of the O5, the $`ON5_B`$ and their variants was given in , using OM6 planes and more elaborate quotients rather than OM5 planes. For 5 branes, A $`(p,q)`$5 brane which wraps $`L_{IIB}`$lifts to an M5 wrapping a $`(p,q)`$cycle of $`𝐓_M^2`$. By analogy, we attempt to lift the O5 to an OM5 wrapping the short side of $`𝐓_M^2`$. Such an OM5 plane is actually a pair of OM5’s because of the transverse compact coordinate on $`𝐓_M^2`$. We check that the charges match: $`2\times Q(OM5)=2\times (1/2)=1=Q(O5)`$.
So far we have discussed the $`O5^{}`$. We would like to construct other discrete torsion variants of the O5, ones which are independent of the compactification on $`L_{IIB}`$, namely, those which are not related to forms that were reduced on $`L_{IIB}`$. This can be done by turning on discrete torsions for M theory on the torus.
Performing $`SL(2,𝐙)`$we can get a family of $`(p,q)`$O5 planes, such that the charge of a $`(p,q)`$$`O5^{}`$ plane is $`1`$ in units of a $`(p,q)`$5-brane. In particular we can consider a $`(0,1)`$ O5 plane which we call $`ON5_B`$ because it is charged under the same field that couples to the NS5 brane of Type IIB (the charge is -1 in units of the NS5 charge). The system of $`ON5_B^{}`$ together with an NS5 brane can be identified to be the $`IIB/I_4()^{F_L}`$ orbifold and can be called $`ON5_B^0`$ (section 5.2; see for a more detailed discussion). A set of N NS5 branes in the vicinity of an $`ON5_B`$ results in 6d worldvolume gauge theory, the same as a set of D5 branes near an O5 plane, with the gauge group being one of $`SO(2N),SO(2N+1),Sp(N)`$ according to the type of the $`ON5_B`$.
Since NS5 branes exist both in Type IIA and in Type IIB, one might expect the ON5 to exist in Type IIA as well. Indeed, the $`IIA/I_4()^{F_L}`$ orbifold (section 5.2), which we call an $`ON5_A^0`$, is a system composed of an $`ON5_A`$ with an NS5 brane. The $`ON5_A^{}`$ has the property that when N NS5 branes coincide with it, the worldvolume theory is a (2,0) CFT with a global symmetry group $`SO(2N)`$, and this is the only possible variant.
### 4.3 O1 - OF1
The case of O1 is quite similar to the O5. Consider an O1 which does not wrap $`L_{IIB}`$(so there is actually a pair of O1’s). As a D1 which does not wrap $`L_{IIB}`$M-lifts to a membrane which wrap the long side of the torus, we should try an OM2 plane wrapping the long side (so again there are actually two of them because of the transverse short side). Let us check the charges: $`2\times (1/16)`$ for the pair of O1’s, indeed equals $`2\times (1/16)`$ for the pair of OM2’s.
For O1 planes which wrap $`L_{IIB}`$, we recall that a D1 wrapping $`L_{IIB}`$is described by a unit of momentum along the short side of $`𝐓_M^2`$. So we try to wrap an OM1 along the short side of $`𝐓_M^2`$as its mass scales like units of momentum (actually it is a pair of OM1’s due to the transverse long side). The charges work out for an $`O1^{}`$ being made of a pair of $`OM1^{}`$: $`1/16=2\times (1/32)`$.
One can get other variants of the O1 by lifting brane intersection to M theory. A new configuration happens for a pair of $`O1^{}`$’s which do not wrap $`L_{IIB}`$and are intersected by a D3 which does. After the intersection we get a pair of $`\stackrel{~}{O3^{}}`$, and so the tension jump is $`2\times +1/2=+1`$. By lifting to M theory we learn that the intersection of a wound OM2 with a transverse M2 gives a $`+1`$ tension jump.
Performing $`SL(2,𝐙)`$we can get a family of $`(p,q)`$O1 planes, such that the charge of a $`(p,q)`$$`O1^{}`$ plane is the same as for a $`(p,q)`$string. In particular we can consider a $`(1,0)`$ O1 plane which we call $`OF1_B`$ because it is charged under the same field that couples to the fundamental string. We will see that the $`OF1_B`$ can be identified with the $`IIB/I_8()^{F_L}`$ orbifold (section 5.1). Since fundamental strings exist both in IIA and in IIB, one might expect the OF1 to exist in Type IIA as well. Indeed, this is the $`IIA/()^{F_L}`$ orbifold (section 5.1), and we call it an $`OF1_A`$. The perturbative orbifold variants may be referred to as $`OF1_B^0,OF1_A^0`$.
## 5 Orbifolds, Orientifolds and New Variants
Here we discuss the relations between orientifolds and orbifolds and their variants. We start with lines, then 5 planes and then the O0. Throughout this section, when we identify a perturbative orbifold with some plane which is a dual of an orientifold, it should be borne in mind that the identification holds only for one variant of the plane, possibly with some extra matter, and all other variants are produced by changing non-perturbative discrete torsions.
### 5.1 Orbifold lines
Orbifold lines together with the O1 form a family connected by dualities. Table 1 is our roadmap for these connections. We will first explore this map and then present some results on the tensions of discrete torsion variants and a relation with K theory.
Our starting point is the O1, that is, $`IIB/I_8\mathrm{\Omega }`$. It carries D1 charge and the forms $`B_{NS},C_0,C_4`$ are odd (twisted) under it. It has discrete torsion variants due to $`H_{NS},H_{NS},G_5,G_1`$ or in terms of brane intersections due to the NS5, F1, D3 and a 7 brane. The 7 brane must allow a D1 charged object to end on it, and so it should be a (0,1) 7 brane rather than a (1,0) D7.
We denote the S dual of the O1 by $`OF1_B`$ as S duality replaces a D1 charge with an F1 charge. S duality replaces $`\mathrm{\Omega }`$ with $`()^{F_L}`$ and so this orbifold is $`IIB/I_8()^{F_L}`$<sup>6</sup><sup>6</sup>6To be more precise, one should note that the $`O1^{}`$, which is the O1 plane with no discrete fluxes carries charge $`\frac{1}{16}`$, which must be cancelled by adding non-perturbative discrete torsion and/ or extra matter, since the orbifold Type $`IIB`$ on $`𝐓^\mathrm{𝟖}/𝐙_\mathrm{𝟐}`$ has 0 charge.. The odd (twisted) forms under the projection are $`C_2,C_0,C_4,C_6`$ and we get 16 (!) discrete torsion variants from all four, or in terms of branes from intersections with D1,D3,D5,D7. Since an F string can end on any D brane, it is natural that each intersection with a D brane is allowed and gives a new variant.
Operating on the $`OF1_B`$ with a T duality in a direction transverse to the fixed line gives an orbifold which we denote by $`OF1_A`$, just as this operation acting on the F string of IIB would give the F string of Type IIA. Such a T duality is accompanied by an additional $`()^{F_L}`$, so operating on $`IIB/I_8()^{F_L}OF1_B`$ we get the $`IIA/I_8`$ orbifold. This orbifold has 8 variants due to intersections with D branes $`D2,D4,D6`$ or their respective forms $`G_2,G_4,G_6`$. A D8 intersection is different because the associated $`G_0`$ form has $`H^0(\mathrm{𝐑𝐏}^\mathrm{𝟕})=𝐙`$ cohomology rather than $`𝐙_2`$ and is interpreted as a change in the Type IIA cosmological constant. The D0 is not in the list since it has no cohomology $`H^8(\mathrm{𝐑𝐏}^\mathrm{𝟕})=\mathrm{𝟎}`$ (but it may produce variants nevertheless).
Other orbifolds can be constructed now by compactifying an OF1 line on a circle and performing parallel T duality. This time one does not add an extra $`()^{F_L}`$. We get $`IIB/I_8`$ and $`IIA/I_8()^{F_L}`$. These orbifolds do not have discrete torsion variants (when uncompactified). By T duality they carry a momentum charge, so we denote them by $`OP1_A,OP1_B`$.
The M lift of the O-lines can be found by looking at their charges. After recalling the M lift of the F string we conclude that the OF1 planes must be wrapping modes of the OM2. The OP1 planes, on the other hand, are an unwrapped OM1. Since the OM1 carries a chiral fermion (section 3.3) after being compactified its Casimir energy will give the required momentum charge.
Let us now find the tensions of some of the discrete torsion variants of the OF1 lines (table 2). The tension of a bare OF1 is $`1/16`$ (in F string units) by S duality with the O1 . It is consistent with the M description as an OM2 wrapping the 11’th dimension, since the OM2 has tension $`1/16`$ (in M2 units).
* To compute the tension of an $`OF1_B`$ after intersecting a D5, consider performing S duality to an O1 which upon intersecting an NS5 changes from $`O1^{}`$ of tension $`1/16`$ to an $`O1^+`$ of tension $`+1/16`$ (so the jump is $`+1/8`$).
* The intersection of an $`OF1_A`$ with a D4 can be M lifted to the basic intersection of an OM2 with an M5 (section 3.2.1), and so the tension jump is $`+1/4`$.
* The tension jump of an $`OF1_B`$ intersecting with a D3 is found again by S duality to be $`+1/2`$.
* The tension jump of an $`OF1_A`$ intersecting a D2 is $`+1`$. This is a consequence of the M theory configuration found in section 4.3, where a wound OM2 intersects a transverse M2. Since the tension jump is integral there is nothing to prevent a whole F1 to separate from this variant.
* We seem to get a rule that after intersecting a Dp brane the jump is $`2^{2p}`$, which is consistent with T duality. For the D6 this rule has an independent check. A D6 intersection corresponds after an M lift to a shift in $`R_{11}`$. So this $`OF1_A`$ variant lifts to M theory on the smooth manifold $`(𝐓^\mathrm{𝟖}\times 𝐒^\mathrm{𝟏})/𝐙_\mathrm{𝟐}`$ where the $`𝐙_2`$ acts by inversion on the first factor and by a shift on the second. As a smooth manifold it carries zero tension, in agreement with a $`1/16`$ jump.
* The intersection of an $`OF1_A`$ with a D0 does not have a discrete cohomology as $`H^8(\mathrm{𝐑𝐏}^\mathrm{𝟕})=\mathrm{𝟎}`$ so it is not clear whether it gives a new variant. If there is a new variant corresponding to intersection of the OF1 with D0 the jump in its charge is 4. Conservation of the fundamental string charge then implies that 4 physical fundamental strings must enter the D0 together with the $`OF1^{}`$ plane. This is indeed consistent with the “fork” configuration of which gives some support for its existence.
One can discuss OF1 orbifolds with a discrete torsion from several forms turned on at the same time. It would be interesting to find their tensions.
The above variants have an interesting relation with K theory. It is simpler to consider $`OF1_A=IIA/I_8`$. Discrete torsion variants are classified by the reduced cohomology $`H^{}(\mathrm{𝐑𝐏}^\mathrm{𝟕})=𝐙_\mathrm{𝟐}𝐙_\mathrm{𝟐}𝐙_\mathrm{𝟐}`$. However, it was recently claimed that the correct classifying group is the reduced K group $`K(\mathrm{𝐑𝐏}^\mathrm{𝟕})=𝐙_\mathrm{𝟖}`$ (“reduced” simply means in both cases that we do not write down a trivial $`𝐙`$ factor). The K group is actually a ring which differs from the standard $`𝐙_8`$ and is defined by the following relations on its generator $`x`$
$`8x=0`$
$`x^2=2x`$ (17)
Note that both the cohomology and the K ring have the same order (8) whereas their structures differ. It would be interesting to elucidate the role of the algebraic structure.
### 5.2 Orbifold 5 planes
We can take a similar tour of orbifold 5 planes, with table 3 as our road map.
We start with the O5 orientifold. It has charge -1 in units of D5, and it has a pair of discrete torsions due to the forms $`G_1,H_{NS}`$, or in terms of branes due to intersection with $`D7,NS5`$.
S duality creates the orbifold $`IIB/I_4()^{F_L}`$, which we denote by $`ON5_B`$ since it carries NS5 charge. In order to cancel the charge one needs to add to the $`ON_5`$ plane an NS5 and this is actually the perturbative orbifold. The matter living on the NS5 of Type IIB is, of course, in (1,1) 6d multiplets. It has variants from the RR forms $`G_1,G_3`$, allowing all the $`SO,Sp`$ gauge groups.
Performing T duality parallel to an $`ON5_B`$ gives Type IIA/$`I_4()^{F_L}`$, which we denote by $`ON5_A`$, because it has NS5 charge as well. It carries matter in (2,0) 6d multiplets. If N NS5 branes are added to the orbifold we get a (2,0) theory with $`SO(2N)`$ group.
We can also perform T duality transverse to the orbifold planes. We get $`IIA/I_4`$ and $`IIB/I_4`$. The first orbifold has a variant due to $`G_2`$. These orbifolds are known to carry non perturbative matter \- a (1,1) theory in the first case, and (2,0) in the second.
Let us consider the M lifts of Type IIA orbifold 5-planes. We can imagine two M theory orbifolds that would give us 5 planes in Type IIA, either an OM6 wound on $`R_{11}`$ or an unwound OM5. By comparing the action on the fields we find that the $`ON5_A=IIA/I_4()^{F_L}`$ lifts to the unwound OM5 while $`IIA/I_4`$ lifts to a wound OM6. The case of Type IIB can be discussed as well, but it has more detail because one needs to specify whether the 5 plane wraps $`L_{IIB}`$or not.
### 5.3 New variants of the O0 plane
For the low dimensional orientifolds, the O0 and the O1, a discrete torsion analysis predicts the existence of more than two $`𝐙_2`$ parameters. For the O1 we saw that there are four $`𝐙_2`$ parameters. In addition to the usual $`(b,c)=H_{NS},G_5`$ there are also a pair of $`𝐙_2`$’s from $`G_1`$ and $`H_{NS}`$.
Similarly, for the O0 we can analyze the possible discrete torsions. In addition to the expected pair of $`𝐙_2`$ parameters $`(b,c)=H_{NS},G_6`$, there is a third $`𝐙_2`$ from $`G_2`$. After an M lift, this additional discrete torsion is nothing but a shift in $`R_{11}`$, the possibility which was not considered section 3.3. It would be interesting to find the mass of the 4 variants with $`[G_2]0`$.
## 6 Miscellaneous applications
### 6.1 The spectrum of 4d $`𝒩=4`$ with $`SO,Sp`$ gauge group
We know that every orientifold 3-plane gives rise to a 4 dimensional gauge theory on D3 probes parallel to it. Each O3 plane gives rise to a theory with a different gauge group $`G`$, according to
1. For $`O3^{}G=SO(2N)`$.
2. For $`\stackrel{~}{O3^{}}G=SO(2N+1)`$.
3. For both $`O3^+,\stackrel{~}{O3^+}`$ the gauge group is $`G=Sp(N)`$.
Since both $`O3^+`$and $`\stackrel{~}{O3^+}`$have the same gauge group one may ask how do the two theories differ. It is clear that the theory with $`\stackrel{~}{O3^+}`$is an $`SL(2,𝐙)`$transform of the one with $`O3^+`$by the element
$$T=\left[\begin{array}{cc}1& 1\\ 0& 1\end{array}\right].$$
(18)
We would like to show how this difference manifests itself in one of the basic measurables of the theory, the spectrum of 1/2 BPS states.
By Definition we know that the W bosons lie in the root lattice of the gauge group $`G`$. Knowing the lattice of the W bosons and the $`SL(2,𝐙)`$transformation which acts naturally both on the charges of the states and on the discrete charges of the orientifolds allows us to characterize the 1/2 BPS spectrum as follows
| | $`O3^{}`$ | $`\stackrel{~}{O3^{}}`$ | $`O3^+`$ | $`\stackrel{~}{O3^+}`$ |
| --- | --- | --- | --- | --- |
| $`(b,c)`$= | (0,0) | (0,1) | (1,0) | (1,1) |
| $`(p,q)`$=(1,0) mod 2 | D | B | C | C |
| $`(p,q)`$=(0,1) mod 2 | D | C | B | C |
| $`(p,q)`$=(1,1) mod 2 | D | C | C | B |
(19)
Here $`D`$ denotes that the states lie in a $`D=SO(2N)`$ lattice, $`B`$ is a $`B=SO(2N+1)`$ lattice and $`C`$ is a $`C=Sp(N)`$ lattice. We get that the charge lattice is of type B iff $`(p,q)=(c,b)\mathrm{mod}2`$.
Note that the difference between $`O3^+`$and $`\stackrel{~}{O3^+}`$is manifest in their spectrum of monopoles and dyons - whereas the monopoles of the $`O3^+`$theory lie in the $`B`$ lattice (the dual lattice), the monopoles of the $`\stackrel{~}{O3^+}`$theory lie in a $`C`$ lattice just like the W bosons.
### 6.2 A comment on allowed BPS stated and $`𝐙_2`$ charges
Consider a configuration with an $`Op`$ plane and a physical Dp brane away from it. Our aim is to ask what are the allowed BPS configurations which can stretch in between the Dp brane and its image. For the simplest case, with $`Op^{}`$ the gauge group is $`SO(2)`$. It is known that a fundamental string stretched between the Dp brane and its image does not lead to a BPS state but rather, as Sen shows in detail in , to a non-BPS state. This happens because the BPS ground state is projected out, so the lowest state is the next massive level which is not BPS.
We would like to extend this discussion to monopole solutions and, when possible, to dyonic states. The first example is for $`p=3`$. Using the results of the previous subsection we find that the allowed BPS states are described by the following table
| (b,c) | (0,0) | (1,0) | (0,1) | (1,1) |
| --- | --- | --- | --- | --- |
| Fundamental String | no | yes | no | yes |
| D1 brane | no | no | yes | yes |
(20)
It is easy to see from this table that a brane is allowed as a BPS state whenever the $`𝐙_2`$ of the two form which couples to it electrically is non trivial.
This leads to the following generalization for any $`p`$. Consider a Dp brane and its image under a reflection by an $`Op`$ plane. Then BPS states arise when either a fundamental string or a Dp-2 brane is stretched between the brane and its image, according to the following table:
| (b,c) | (0,0) | (1,0) | (0,1) | (1,1) |
| --- | --- | --- | --- | --- |
| Fundamental String | no | yes | no | yes |
| Dp-2 brane | no | no | yes | yes |
(21)
Unlike the previous case, the simple rule that a brane is allowed to be stretched in between the heavier brane and its image whenever the corresponding $`𝐙_2`$ flux of the form which couples to it electrically is non-trivial does not apply here. It is not clear how the picture generalizes.
Similar statements hold for D3 branes stretching between NS5 branes in the presence of the different types of $`ON5_B`$ planes.
### 6.3 Orientifold webs
Since we identified $`(p,q)`$O1 lines and $`(p,q)`$O5 planes one may wonder whether $`(p,q)`$webs of orientifolds are possible. We would like to show that those may be possible in some special cases, but in general they do not make sense.
Let us consider the basic junction at $`\tau =i`$, where $`\tau `$ is the complex scalar of Type IIB. We have a horizontal $`(0,1)`$ O1 meeting a vertical $`(1,0)`$ OF1 and a $`(1,1)`$ O line which leaves the junction at 45 degrees. Each orientifold plane requires a $`𝐙_2`$ projection. If we ignore the orientation reversal and consider only the spatial inversion we see that these 3 reflections generate a group of order 8, isomorphic to $`D_4`$, the dihedral group of 4 elements. It is not clear how do the different orientation reversals combine.
However, if there are two O lines at an irrational relative angle $`\alpha `$ (measured in radians/$`2\pi `$) then a composition of both reflections gives a rotation by $`2\alpha `$. Since we assumed $`\alpha `$ irrational, then this element generates an infinite group of identifications on the plane, generating a non-discrete image set from a single point. Thus we cannot hope to have an ordinary orbifold/orientifold, and this is the case for generic $`\tau `$.
Acknowledgements
We thank O. Bergman, E. Gimon, S. Gukov, E. Shustin, C. Vafa, E. Witten and A. Zaffaroni for enjoyable discussions.
A.H. would like to thank the Institute for Theoretical Physics at Santa Barbara and Tel-Aviv University for their kind hospitality while various stages of this work were completed. B.K. would like to thank J. Sonnenschein, S. Yankielowicz and the rest of the group at Tel Aviv.
A. H. is partially supported by the National Science Foundation under grant no. PHY94-07194, by the DOE under grant no. DE-FC02-94ER40818, by an A. P. Sloan Foundation Fellowship and by a DOE OJI award. The research of BK was supported in part by the US-Israeli Binational Science Foundation, the German–Israeli Foundation for Scientific Research (GIF), and the Israel Science Foundation.
## Appendix A Appendix - Cohomology of $`\mathrm{𝐑𝐏}^𝐧`$
We distinguish between two kinds of cohomologies. The twisted ones, denoted by $`\stackrel{~}{H}`$ classify “twisted” forms. These are forms which reverse sign under the projection (such forms are not related to a twisted sector). The ordinary “untwisted” cohomologies are denoted just by $`H`$ and are appropriate to classify forms which do not change sign under the projection.
The integral cohomologies of $`\mathrm{𝐑𝐏}^𝐧`$ are
$`H^q=`$ $`\{\begin{array}{cc}𝐙_2\hfill & \text{ }q\text{ even but }q0\hfill \\ 𝐙\hfill & q=0,\text{ and for odd n }q=n\hfill \\ 0\hfill & \text{otherwise}\hfill \end{array}`$ (25)
$`\stackrel{~}{H}^q=`$ $`\{\begin{array}{cc}𝐙_2\hfill & q\text{ odd}\hfill \\ 𝐙\hfill & \text{for even n }q=n\hfill \\ 0\hfill & \text{otherwise}\hfill \end{array}`$ (29)
These results can be easily deduced from the chain complexes
$$\begin{array}{c}C^2=𝐙\stackrel{\times 2}{}C^1=𝐙\stackrel{0}{}C^0=𝐙0\\ \stackrel{\times 2}{}C^{2m+1}=𝐙\stackrel{0}{}C^{2m}=𝐙\stackrel{\times 2}{}\mathrm{}\\ 0C^n=𝐙\mathrm{}\end{array}$$
(30)
and
$$\begin{array}{c}\stackrel{~}{C}^2=𝐙\stackrel{0}{}\stackrel{~}{C}^1=𝐙\stackrel{\times 2}{}\stackrel{~}{C}^0=𝐙0\\ \stackrel{0}{}\stackrel{~}{C}^{2m+1}=𝐙\stackrel{\times 2}{}\stackrel{~}{C}^{2m}=𝐙\stackrel{0}{}\mathrm{}\\ 0\stackrel{~}{C}^n=𝐙\mathrm{}\end{array}$$
(31)
where $`C^q`$ are the q-cochains and $`\stackrel{~}{C}^q`$ are the twisted q-cochains.
For completeness we list the integral homology groups as well
$`H_q=`$ $`\{\begin{array}{cc}𝐙_2\hfill & q\text{ odd but }qn\hfill \\ 𝐙\hfill & q=0,\text{ and for odd n }q=n\hfill \\ 0\hfill & \text{otherwise}\hfill \end{array}`$ (35)
$`\stackrel{~}{H}_q=`$ $`\{\begin{array}{cc}𝐙_2\hfill & q\text{ even }qn\hfill \\ 𝐙\hfill & \text{for even n }q=n\hfill \\ 0\hfill & \text{otherwise}\hfill \end{array}`$ (39)
For odd $`n`$ $`\mathrm{𝐑𝐏}^𝐧`$ is orientable and Poincare duality holds
$$H_i=H^{ni},\stackrel{~}{H}_i=\stackrel{~}{H}^{ni}$$
(40)
For even $`n`$, on the other hand
$$H_i=\stackrel{~}{H}^{ni},\stackrel{~}{H}_i=H^{ni}$$
(41)
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# Output spectrum of a detector measuring quantum oscillations
## I Introduction
The long-standing and still controversial problem of quantum measurements is gradually becoming a practical issue due to the development of experimental techniques capable of measurements at the quantum border (see, e.g. Refs. ). The renewed interest in this subject is caused by the needs of quantum computing, since the measurement of an entangled and possibly evolving quantum state by a realistic detector in a realistic environment presents a nontrivial problem.
Despite the experimental proof of the violation of Bell’s inequality that rejects the idea of hidden variables, the origin of the randomness of a quantum measurement result (the problem known as the “wavefunction collapse”) remains controversial. Important insight into this problem has been provided by the development of the theory of continuous quantum measurement, which generalizes the “orthodox” case of instantaneous measurement. There are two different theoretical approaches. In the first approach based on the theory of interaction with a dissipative environment, the evolution of the density matrix of the measured system is averaged over a complete ensemble of measurements, thus leading to the deterministic equation. This is the best-known approach, at least in the solid-state community, and so can be called “conventional”. The other approach (more developed in quantum optics) studies the random evolution of the quantum state during a single realization of the measurement, so that this evolution is conditioned on (selected by) the particular measurement result. Recently the latter approach has been introduced into the context of solid-state physics using the new derivation based on the simple Bayesian analysis of probabilities.
In the present paper this Bayesian formalism is applied to the calculation of the spectral density of the detector current when a two-level quantum system (qubit) is measured continuously. As a particular example, we consider a double-quantum-dot occupied by one electron, the location of which is measured by a quantum point contact nearby. Another example of the setup is a single-Cooper-pair box, being measured by a single-electron transistor. One more possible example is the continuous measurement of two flux states of a SQUID by another inductively coupled SQUID.
We show that in the weak-coupling case when the quantum (Rabi) oscillations of the qubit state are only slightly perturbed by the detector, the corresponding peak in the spectral density of the detector current can be up to 4 times higher than the noise pedestal (see also Ref. ). As the coupling increases, the Quantum Zeno effect becomes significant leading to the Lorentzian shape of the spectral density centered at zero frequency.
It is important to notice that there should be no difference between the predictions of the conventional approach and the approach of selective evolution unless the measurement result affects the system evolution (for example, via the feedback loop). We prove this equivalence explicitly for the detector spectral density (if there is no feedback). Despite the same final result, the interpretations are different: in the Bayesian formalism a significant contribution to the spectrum comes from the correlation between the detector noise and the system evolution, while this correlation is absent in the conventional approach. In the paper we also discuss the extension of the Bayesian formalism to the case of additional weak interaction of the two-level system with a finite-temperature environment.
## II Bayesian formalism
For a two-level quantum system described by the Hamiltonian
$$_0=\frac{\epsilon }{2}(c_1^{}c_1c_2^{}c_2)+H(c_1^{}c_2+c_2^{}c_1)$$
(1)
(where $`\epsilon `$ is the energy asymmetry and $`H`$ is the tunneling strength) the evolution of its density matrix $`\rho _{ij}`$ due to continuous measurement is given in the conventional approach by the equations
$`\dot{\rho }_{11}=2{\displaystyle \frac{H}{\mathrm{}}}\text{Im}\rho _{12},\rho _{11}+\rho _{22}=1,`$ (2)
$`\dot{\rho }_{12}=ı{\displaystyle \frac{\epsilon }{\mathrm{}}}\rho _{12}+ı{\displaystyle \frac{H}{\mathrm{}}}(\rho _{11}\rho _{22})\mathrm{\Gamma }\rho _{12},`$ (3)
where the dephasing rate $`\mathrm{\Gamma }`$ due to measurement depends on the type of the detector used.
Eqs. (2)–(3) describe the averaged evolution. In contrast, to analyze the single measurement process we need the selective (conditional) evolution of $`\rho _{ij}`$ which in the Bayesian formalism is described by equations
$`\dot{\rho }_{11}=`$ $`2{\displaystyle \frac{H}{\mathrm{}}}\text{Im}\rho _{12}+{\displaystyle \frac{2\mathrm{\Delta }I}{S_0}}\rho _{11}\rho _{22}[I(t)I_0],`$ (4)
$`\dot{\rho }_{12}=`$ $`ı{\displaystyle \frac{\epsilon }{\mathrm{}}}\rho _{12}+ı{\displaystyle \frac{H}{\mathrm{}}}(\rho _{11}\rho _{22})`$ (6)
$`{\displaystyle \frac{\mathrm{\Delta }I}{S_0}}(\rho _{11}\rho _{22})[I(t)I_0]\rho _{12}\gamma \rho _{12},`$
where $`I(t)`$ is the detector current, $`I_0=(I_1+I_2)/2`$, $`I_1`$ and $`I_2`$ are the average currents corresponding to two localized states of the qubit, $`\mathrm{\Delta }I=I_1I_2`$ (notice the different sign in the definition used in Ref. ), and $`S_0`$ is the low-frequency spectral density of the detector shot noise (which is assumed to be constant in the frequency range of interest and to be practically independent of the qubit state). The detector nonideality is described by the dephasing contribution $`\gamma =\mathrm{\Gamma }(\mathrm{\Delta }I)^2/4S_00`$ due to interaction with “pure environment” (which does not affect the detector current). In particular, since $`\mathrm{\Gamma }=(\mathrm{\Delta }I)^2/4S_0`$ for symmetric quantum point contact (see Refs. ), it represents an ideal detector, $`\eta =1`$, where $`\eta =1\gamma /\mathrm{\Gamma }`$ is the ideality factor. In contrast, the single-electron transistor in a typical operation point is a significantly nonideal detector, $`\eta 1`$. (Actually, for the single-electron transistor Eq. (6) can be further improved, however, it does not make a difference for the purposes of the present paper.) The SQUID can be an ideal detector when its sensitivity is quantum-limited. Since the typical output signal from the SQUID is the voltage (not the current), this requires a minor modification of the formalism; so in this paper we do not consider the SQUID case explicitly even though all the results can be easily translated into SQUID language.
Eqs. (4)–(6) allow us to calculate the evolution of $`\rho _{ij}`$ if the detector output $`I(t)`$ is known from the experiment. In order to simulate the measurement we need the replacement
$$I(t)I_0=\mathrm{\Delta }I(\rho _{11}\rho _{22})/2+\xi (t),$$
(7)
where the random process $`\xi (t)`$ has zero average and “white” spectral density $`S_\xi =S_0`$.
Notice that in order to consider the detector as a device with classical output signal, Eqs. (4)–(6) essentially use the Markov approximation and assume that the typical frequency of the internal detector evolution (on the order of $`S_0/e^2I_0/e`$) is much higher than the typical frequency $`\mathrm{max}(\mathrm{\Omega },\mathrm{\Gamma })`$ of $`\rho _{ij}`$ evolution (here $`\mathrm{\Omega }\sqrt{4H^2+\epsilon ^2}/\mathrm{}`$ is the frequency of unperturbed quantum oscillations). In particular, this condition requires the detector to be “weakly responding”, $`|\mathrm{\Delta }I|I_0`$, that allows us to use the linear response approximation.
Averaging of Eqs. (4)–(6) over all possible measurement results (i.e. over random contribution $`\xi (t)`$) reduces them to Eqs. (2)–(3). Notice that the stochastic equations are written in Stratonovich form which preserves the usual calculus rules, while averaging would be easier in Itô form.
## III Weak coupling
Using Eqs. (4)–(7) and the Monte-Carlo method (similar to Ref. ) we can calculate in a straightforward way the spectral density $`S_I(\omega )`$ of the detector current $`I(t)`$. Solid lines in Fig. 2 show the results of such calculations for the ideal detector, $`\eta =1`$, and weak coupling between the qubit and the detector, $`\alpha =0.1`$, where $`\alpha \mathrm{}(\mathrm{\Delta }I)^2/8S_0H`$ ($`\alpha `$ is 8 times less than the parameter $`𝒞`$ introduced in Ref. ). One can see that in the symmetric case, $`\epsilon =0`$, the peak at the frequency of quantum oscillations is 4 times higher than the noise pedestal, $`S_I(\mathrm{\Omega })=5S_0`$ while the peak width is determined by the coupling strength $`\alpha `$ (see Fig. 5 below). In the asymmetric case, $`\epsilon 0`$, the peak height decreases (Fig. 2), while the additional Lorentzian-shape increase of $`S_I(\omega )`$ appears at low frequencies. The origin of this low-frequency feature is the slow fluctuation of the asymmetry of $`\rho _{11}`$ oscillations (Fig. 3). In case $`\epsilon =0`$ the amplitude of $`\rho _{11}`$ oscillations is maximal (see thick line in Fig. 4a), hence there is no such asymmetry and the low-frequency feature is absent, while the spectral peak at the frequency of quantum oscillation is maximally high.
In order to understand the factor 4 for the maximum peak height, let us consider the case $`\alpha 1`$, $`\epsilon =0`$, and $`\eta =1`$. Then the selective evolution can be written as the quantum oscillations with slowly fluctuating phase $`\phi (t)`$:
$`z(t)`$ $`\rho _{11}(t)\rho _{22}(t)=\mathrm{cos}\varphi (t),\varphi =\mathrm{\Omega }t+\phi (t),`$ (9)
$`\rho _{12}=(ı/2)\mathrm{sin}\varphi (t)`$
(the state becomes pure after an initial transient period since $`\eta =1`$, while the real part of $`\rho _{12}`$ decays because of $`\epsilon =0`$). From Eqs. (4) and (7) we obtain the random phase dynamics:
$$\dot{\phi }=\mathrm{sin}\varphi \frac{\mathrm{\Delta }I}{S_0}[\mathrm{cos}\varphi \frac{\mathrm{\Delta }I}{2}+\xi (t)].$$
(10)
Since $`(\mathrm{\Delta }I)^2/2S_0\mathrm{\Omega }`$, we can neglect the first term in the square brackets and average the second contribution over $`\varphi `$ that leads to the white spectrum of $`\dot{\phi }`$: $`S_{\dot{\phi }}(\omega )=(\mathrm{\Delta }I)^2/2S_0`$. Then the correlation function $`K_z(\tau )z(t)z(t+\tau )`$ can be calculated as $`K_z(\tau )=\mathrm{cos}(\mathrm{\Omega }\tau )\mathrm{cos}\mathrm{\Delta }\phi (\tau )/2=\mathrm{cos}(\mathrm{\Omega }\tau )\mathrm{exp}[(\mathrm{\Delta }I)^2\tau /8S_0]/2`$ and the spectral density $`S_z(\omega )2_{\mathrm{}}^{\mathrm{}}K_z(\tau )\mathrm{exp}(ı\omega \tau )𝑑\tau `$ has a peak in the vicinity of the oscillation frequency, $`|\omega \mathrm{\Omega }|\mathrm{\Omega }`$:
$$S_z(\omega )=\frac{8S_0}{(\mathrm{\Delta }I)^2}\frac{1}{1+[8S_0(\omega \mathrm{\Omega })/(\mathrm{\Delta }I)^2]^2}.$$
(11)
The detector current is given by Eq. (7), so its spectral density consists of three contributions:
$$S_I(\omega )=S_0+\frac{\mathrm{\Delta }I^2}{4}S_z(\omega )+\frac{\mathrm{\Delta }I}{2}[S_{\xi z}(\omega )+S_{z\xi }(\omega )],$$
(12)
where the last contribution originates from the correlation between $`\rho _{ij}`$ evolution and the detector noise $`\xi (t)`$. To calculate the correlation function $`K_{\xi z}(\tau )\xi (t)z(t+\tau )`$ for $`\tau >0`$, we need to take into account the phase shift $`d\phi =\mathrm{sin}\varphi \mathrm{\Delta }IS_0^1\xi (t)dt`$ during even infinitesimally small time duration $`dt`$, since the amplitude of the stochastic function $`\xi (t)`$ grows with the timescale decrease, $`\xi (t)^2dt=S_0/2`$. Using trigonometric formulas and linear expansion in $`d\phi `$ we obtain $`\xi (t)\mathrm{cos}[\varphi (t)+d\varphi +\mathrm{\Omega }\tau +\mathrm{\Delta }\phi (\tau )]=\mathrm{\Delta }IS_0^1\xi (t)^2dt\mathrm{sin}\varphi (t)\mathrm{sin}[\varphi (t)+\mathrm{\Omega }\tau ]\mathrm{cos}\mathrm{\Delta }\phi (\tau )`$ and finally $`K_{\xi z}(\tau )=\mathrm{\Delta }I\mathrm{cos}(\mathrm{\Omega }\tau )\mathrm{exp}[(\mathrm{\Delta }I)^2\tau /8S_0]/4`$. After the Fourier transformation one finds that the correlation between $`\xi (t)`$ and $`z(t)`$ brings exactly the same contribution to the detector spectral density (see Eq. (12)) as the term due to $`z(t)`$ evolution, so that
$$S_I(\omega )=S_0+\frac{4S_0}{1+[8S_0(\omega \mathrm{\Omega })/(\mathrm{\Delta }I)^2]^2}.$$
(13)
Thus, the peak corresponding to quantum oscillations is 4 times higher than the noise background, while its full width at half height is equal to $`(\mathrm{\Delta }I)^2/4S_0=\alpha \mathrm{\Omega }`$ (the same peak width has been calculated in Ref. ). The integral under the peak,
$$_0^{\mathrm{}}[S(\omega )S_0]\frac{d\omega }{2\pi }=\frac{(\mathrm{\Delta }I)^2}{4},$$
(14)
has an obvious relation to the average square of the detector current variation due to oscillations in the measured system. Notice, however, that this integral is twice as large as one would expect from the classical harmonic signal, since one half of the peak height comes from nonclassical correlation between the qubit evolution and the detector noise. Classically, Eq. (14) would be easily understood if the signal was not harmonic but rectangular-like, which is obviously not the case. Actually, the detector current shows neither clear harmonic nor rectangular signal distinguishable from the intrinsic noise contribution. Figure 4a shows the simulation of $`\rho _{11}`$ evolution (thick line) together with the detector current $`I(t)`$. Since $`I(t)`$ contains white noise, it necessarily requires some averaging. Thin solid, dotted, and dashed lines show the detector current averaged with different time constants $`\tau _a`$: $`\tau _a\mathrm{\Omega }/2\pi =0.3`$, 1, and 3, respectively. For weak averaging the signal is too noisy, while for strong averaging individual oscillation periods cannot be resolved either, so quantum oscillations can never be observed directly by a continuous measurement (although they can be calculated using Eqs. (4)–(6)). This unobservability is revealed in the relatively low peak height of the spectral density of the detector current.
## IV Arbitrary coupling
The situation changes as the coupling between the detector and qubit increases, $`\alpha 1`$. The strong influence of measurement destroys quantum oscillations, and the Quantum Zeno effect develops, so that for $`\alpha 1`$ the qubit performs random jumps between two localized states (see Fig. 4b). In this case the properly averaged detector current follows pretty well the evolution of the qubit (however, the unsuccessful tunneling “attempts” still cannot be directly resolved), and the spectral density of $`I(t)`$ can be calculated using the classical theory of telegraph noise leading to the Lorentzian shape of $`S_I(\omega )`$. Figure 5a shows the gradual transformation of the spectral density with the increase of the coupling $`\alpha `$ for a symmetric qubit, $`\epsilon =0`$, and an ideal detector, $`\eta =1`$. The results for an asymmetric qubit, $`\epsilon /H=1`$, are shown in Fig. 5b.
The curves in Fig. 5 as well as the dashed curves in Fig. 2 are calculated using the conventional master equation approach which gives the same results for the detector spectral density as the Bayesian formalism (we will prove this later). In the conventional approach we should assume no correlation between the detector noise and the qubit evolution (the last term in Eq. (12) is absent) while the correlation function $`K_{\widehat{z}}(\tau )`$ should be calculated considering $`z(t)`$ not as an ordinary function but as an operator. Then the calculation of $`\widehat{z}(t)\widehat{z}(t+\tau )`$ can be essentially interpreted as follows. The first (in time) operator $`\widehat{z}(t)`$ collapses the qubit into one of two eigenstates which correspond to localized states, then during time $`\tau `$ the qubit performs the evolution described by conventional Eqs. (2)–(3), and finally the second operator $`\widehat{z}(t+\tau )`$ gives the probability for the qubit to be measured in one of two localized states. (Of course, this procedure can be done purely formally, without any interpretation.) Notice that there is complete symmetry between states “1” and “2” even for $`\epsilon 0`$ (in particular, in the stationary state $`\rho _{11}=\rho _{22}=1/2`$), so the evolution after the first collapse can be started from any localized state leading to the same contribution to the correlation function. In this way we obviously get $`K_{\widehat{z}}(\tau )=\rho _{11}(\tau )\rho _{22}(\tau )`$ where $`\rho _{ii}`$ is the solution of Eqs. (2)–(3) with the initial conditions $`\rho _{11}(0)=1`$ and $`\rho _{12}(0)=0`$.
For the symmetric qubit, $`\epsilon =0`$, these equations can be easily solved analytically, and finally we obtain:
$$S_I(\omega )=S_0+\frac{\mathrm{\Omega }^2(\mathrm{\Delta }I)^2\mathrm{\Gamma }}{(\omega ^2\mathrm{\Omega }^2)^2+\mathrm{\Gamma }^2\omega ^2},$$
(15)
where $`\mathrm{\Gamma }=\eta ^1(\mathrm{\Delta }I)^2/4S_0=\alpha \eta ^1\mathrm{\Omega }`$. This equation obviously transforms into Eq. (13) for $`\eta =1`$ and $`\alpha 1`$. Notice that for weak coupling with a nonideal detector, $`\eta <1`$ and $`\alpha \eta ^11`$, the peak height of $`S_I(\omega )`$ is equal to $`4\eta S_0`$, while the linewidth $`\alpha \eta ^1\mathrm{\Omega }`$ of the peak is $`\eta ^1`$ times wider than for the ideal detector. As the coupling increases, the linewidth grows and the oscillation frequency decreases: $`\omega _{osc}=\mathrm{\Omega }[1(\alpha /2\eta )^2)^{1/2}]`$. The transition into the overdamped regime occurs at $`\alpha \eta ^1>2`$ while the peak-like feature disappears at $`\alpha \eta ^1>\sqrt{2}`$. For $`\alpha \eta ^1>2`$ the spectral density consists of two Lorentzians \[$`\omega _{1,2}=\mathrm{\Gamma }/2(\mathrm{\Gamma }^2/4\mathrm{\Omega }^2)^{1/2}`$\] centered at zero frequency, with the negative sign and the smaller amplitude $`A_2`$ of the second Lorentzian which has higher cutoff frequency, $`A_2/A_1=\omega _1/\omega _2`$. In the case $`\alpha \eta ^11`$, which corresponds to the well developed Quantum Zeno effect, $`S_I(\omega )S_0`$ has purely Lorentzian shape $`(\mathrm{\Delta }I)^2\omega _1/(\omega ^2+\omega _1^2)`$ with $`\omega _1=\mathrm{\Omega }^2/\mathrm{\Gamma }=\mathrm{\Omega }\eta /\alpha `$.
For the asymmetric qubit, $`\epsilon 0`$, the spectral density can in principle also be calculated analytically but the expressions would be too lengthy and it is simpler to use numerical solution of Eqs. (2)–(3). The analytical formula for the weak coupling limit is
$`S_I(\omega )=`$ $`S_0+{\displaystyle \frac{\eta S_0\epsilon ^2/H^2}{1+(\omega \mathrm{}^2\mathrm{\Omega }^2/4H^2\mathrm{\Gamma })^2}}`$ (17)
$`+{\displaystyle \frac{4\eta S_0(1+\epsilon ^2/2H^2)^1}{1+[(\omega \mathrm{\Omega })/\mathrm{\Gamma }(12H^2/\mathrm{}^2\mathrm{\Omega }^2)]^2}}.`$
The spectral peak and the low-frequency Lorentzian become wider with the coupling increase since $`\mathrm{\Gamma }=\alpha \eta ^1\mathrm{\Omega }`$, and for $`|\epsilon /H|<1/\sqrt{2}`$ the overdamped regime starts from $`\mathrm{\Gamma }=\mathrm{\Gamma }_1`$ where $`\mathrm{\Gamma }_{1,2}^2=(\mathrm{\Omega }^2/2a)[b(b^24a)^{1/2}]`$, $`b1/427a^2/4+9a/2`$, $`a\epsilon ^2/(4H^2+\epsilon ^2)`$. At $`\mathrm{\Gamma }>\mathrm{\Gamma }_2`$ the dynamics formally returns to the underdamped regime, however, the peak linewidth is much larger than the frequency and so $`S_I(\omega )`$ is monotonous. For $`|\epsilon /H|>1/\sqrt{2}`$ the overdamped regime does not occur. In both cases in the limit of large $`\mathrm{\Gamma }`$ the spectral density has almost Lorentzian shape with the cutoff frequency $`\omega _1=4H^2/\mathrm{}^2\mathrm{\Gamma }`$.
One can check that the spectral densities given by Eqs. (15) and (17) satisfy the integral condition (14), which remains valid for arbitrary parameters $`\alpha `$, $`\epsilon /H`$, and $`\eta `$, because of the equation $`K_I(+0)=(\mathrm{\Delta }I/2)^2`$.
## V Equivalence of two approaches
Comparing two derivations of $`S_I(\omega )`$ in the case $`\alpha 1`$, $`\eta =1`$, and $`\epsilon =0`$, we see that $`K_{\widehat{z}}(\tau )`$ is twice larger than $`K_z(\tau )`$ because in the conventional approach the evolution always starts from the localized state while in the Bayesian approach it starts from arbitrary phase of quantum oscillations. This difference exactly compensates for the absence of the last term in Eq. (12) in the conventional approach.
Let us prove explicitly that the two approaches give the same result for $`S_I(\omega )`$ in a general case. In order to calculate $`K_{\xi z}(\tau )`$ for $`\tau >0`$ using the Bayesian formalism, let us first average the product $`\xi (t_0)z(t_0+\tau )`$ over random $`\xi (t)`$ during time period $`t_0<t<t_0+\tau `$, fixing the same conditions at $`t=t_0`$. Then we can use conventional Eqs. (2)–(3) \[regarded as Eqs. (4)–(7) averaged over random $`\xi (t)`$\] with the initial condition $`\rho _{ij}(t_0+0)=\rho _{ij}(t_0)+d\rho _{ij}`$ where
$`dz=\mathrm{\Delta }IS_0^1[1z^2(t_0)]\xi (t_0)dt,`$ (18)
$`d\rho _{12}=\mathrm{\Delta }IS_0^1z(t_0)\rho _{12}(t_0)\xi (t_0)dt`$ (19)
(for simplicity we will refer to $`z\rho _{11}\rho _{22}`$ as a component of $`\rho _{ij}`$). Since the sign of $`\xi (t_0)`$ is arbitrary and averaged evolution equations are linear, we need only the fluctuating contribution to $`\rho _{ij}(t_0+0)`$ and, hence, can formally assume that the evolution starts from $`\rho _{ij}(t_0+0)=d\rho _{ij}`$ (notice that we can forget the condition $`\rho _{11}+\rho _{22}=1`$ and use only $`z`$ and $`\rho _{12}`$). Using the relation $`\xi (t_0)^2dt=S_0/2`$ and the evolution linearity, we can formally write
$$K_{\xi z}(\tau )=(\mathrm{\Delta }I/2)\stackrel{~}{z}(t_0+\tau ),$$
(20)
where $`\stackrel{~}{\rho }_{ij}`$ satisfies Eqs. (2)–(3) with $`\stackrel{~}{z}(t_0)=1z(t_0)^2`$ and $`\stackrel{~}{\rho }_{12}(t_0)=z(t_0)\rho _{12}(t_0)`$, and the averaging over the initial conditions at $`t=t_0`$ should still be done later. Before that let us do a similar formal trick for $`K_z(\tau )`$ representing it as $`\stackrel{~}{z}(t_0+\tau )`$ where the evolution starts from $`\stackrel{~}{z}(t_0)=z(t_0)^2`$ and $`\stackrel{~}{\rho }_{12}(t_0)=z(t_0)\rho _{12}(t_0)`$. Now combining two terms in the detector current correlation function $`K_I(\tau )I(t_0)I(t_0+\tau )I^2=(\mathrm{\Delta }I/2)^2K_z(\tau )+(\mathrm{\Delta }I/2)K_{\xi z}(\tau )`$ (here $`\tau >0`$), we see that it can be written as $`(\mathrm{\Delta }I/2)^2\stackrel{~}{z}(t_0+\tau )`$ where $`\stackrel{~}{z}(t_0)=1`$ and $`\stackrel{~}{\rho }_{12}(t_0)=0`$. Thus we have exactly arrived at the expression of the conventional approach, in which the evolution always starts from the localized state, regardless of the actual quantum state at $`t=t_0`$. This proof is obviously valid for arbitrary $`\alpha `$, $`\eta `$, and $`\epsilon /H`$. Despite the same result in two approaches, the interpretations are quite different since the Bayesian approach allows us to follow the qubit evolution during the measurement process, while the conventional approach gives only the average characteristics.
## VI Finite-temperature environment
In this section we will discuss how to introduce the finite-temperature environment into Eqs. (4)–(7) of the Bayesian formalism. Notice that so far there has been complete symmetry between the states “1” and “2” even for finite energy difference $`\epsilon `$, while the finite temperature effects would be expected to lead to different average populations of these states. Such symmetry requires an implicit assumption that the typical energy in the detector (voltage or temperature) is much larger than the energies involved in the qubit dynamics. So, the absence of temperature in the formalism does not mean that it is zero or very large, it just means that the temperature effects are not important. Now let us assume that besides the detector, the qubit is coupled to an additional finite-temperature environment, which creates an asymmetry between states “1” and “2” when $`\epsilon 0`$.
While the case of finite coupling of a two-level system with an environment presents a really difficult problem, the weak coupling limit can be treated in a simple way. In the standard method it is described by the equations
$`\dot{\rho }_{++}=\gamma _1(\rho _{++}p_{st}),\rho _{++}+\rho _{}=1,`$ (21)
$`\dot{\rho }_+=ı\mathrm{\Omega }\rho _+\gamma _2\rho _+,`$ (22)
which are written in the diagonal basis (“+” corresponds to the ground state). The temperature $`T`$ determines the stationary occupation $`p_{st}=[1+\mathrm{exp}(\mathrm{}\mathrm{\Omega }/T)]^1`$ of the ground state, and the relaxation rates obey inequality $`\gamma _1/2\gamma _2\mathrm{\Omega }`$.
If the coupling of the qubit with the detector is also weak, $`\alpha \eta ^11`$, the evolution due to extra finite-temperature environment can be simply added to the evolution due to measurement. For this purpose Eqs. (21)–(22) should be translated into the basis of localized states, so the terms
$`(A^2\gamma _1+B^2\gamma _2)(\rho _{11}1/2)\gamma _1A(1/2p_{st})`$ (23)
$`+AB(\gamma _1\gamma _2)\text{Re}\rho _{12},`$ (24)
$`\text{where}A\epsilon /\mathrm{}\mathrm{\Omega },B2H/\mathrm{}\mathrm{\Omega },`$ (25)
should be added into Eq. (4) for $`\dot{\rho }_{11}`$ and the terms
$`(A^2\gamma _2+B^2\gamma _1)\text{Re}\rho _{12}+AB(\rho _{11}\rho _{22})(\gamma _1\gamma _2)/2`$ (26)
$`+\gamma _1B(1/2p_{st})ı\gamma _2\text{Im}\rho _{12}`$ (27)
should be added into Eq. (6) for $`\dot{\rho }_{12}`$. The same terms should obviously be added into Eqs. (2)–(3) for the conventional approach. (Of course, this generalization is purely phenomenological and is limited to the weak coupling regime, so the effect of Eqs. (24)–(27) can be considered only at the timescale longer than oscillation period.)
In the generalized case it is still possible to prove that the results of the Bayesian formalism for the detector current spectral density $`S_I(\omega )`$ exactly coincide with the results of the conventional approach. The essential difference from the proof above is nonzero stationary solution ($`z_{st},\rho _{12,st}`$) of modified Eqs. (2)–(3) when $`p_{st}1/2`$. It is convenient to consider homogeneous evolution equations (with $`p_{st}=1/2`$) simply shifting $`z(t)`$ and $`\rho _{12}(t)`$ by the stationary values. Using the same idea as in the proof above we can show that in the Bayesian approach $`K_I(\tau )`$ for $`\tau >0`$ can be written as $`(\mathrm{\Delta }I/2)^2\stackrel{~}{z}(t_0+\tau )`$ where $`\stackrel{~}{\rho }_{ij}`$ satisfies homogeneous modified Eqs. (2)–(3) with $`\stackrel{~}{z}(t_0)=12z_{st}z(t_0)+z_{st}^2`$ and $`\stackrel{~}{\rho }_{12}(t_0)=z(t_0)\rho _{12,st}z_{st}\rho _{12}(t_0)+z_{st}\rho _{12,st}`$. After the averaging over initial states these initial conditions can be replaced with $`\stackrel{~}{z}(t_0)=1z_{st}^2`$ and $`\stackrel{~}{\rho }_{12}(t_0)=z_{st}\rho _{12,st}`$.
Now let us show that we get the same $`K_I(\tau )`$ in the conventional approach. With the probability $`(1+z_{st})/2`$ the first operator $`\widehat{z}(t_0)`$ localizes the qubit into the state “1”. Then the initial state for the homogeneous equations is $`\stackrel{~}{z}(t_0)=1z_{st}`$, $`\stackrel{~}{\rho }_{12}(t_0)=\rho _{12,st}`$. With the probability $`(1z_{st})/2`$ the evolution starts from the state “2”, i.e. $`\stackrel{~}{z}(t_0)=1z_{st}`$, $`\stackrel{~}{\rho }_{12}(t_0)=\rho _{12,st}`$. Adding two contributions with opposite signs we see again that for $`\tau >0`$, $`K_I(\tau )=(\mathrm{\Delta }I/2)^2\stackrel{~}{z}(t_0+\tau )`$ where $`\stackrel{~}{z}`$ can be found as a solution of Eqs. (2)–(3) modified by Eqs. (24)–(27) without inhomogeneous terms, with the initial condition $`\stackrel{~}{z}(t_0)=1z_{st}^2`$ and $`\stackrel{~}{\rho }_{12}(t_0)=z_{st}\rho _{12,st}`$. Thus, the correlation function $`K_I(\tau )`$ and, hence, the spectral density $`S_I(\omega )`$ coincide in the two approaches.
It is technically simpler to consider the averaged evolution in the diagonal basis rather than in the basis of localized states \[for this purpose we need to translate the term $`\mathrm{\Gamma }\rho _{12}`$ from Eq. (3) into the diagonal basis and add it into Eqs. (21)–(22)\]. So, to calculate the correlation function $`K_{\widehat{z}}(\tau )`$ analytically, we start the evolution from one of the localized states, then consider the averaged evolution in the diagonal basis (neglecting the rapidly oscillating terms due to measurement), and make the second projection onto localized states at $`t=\tau `$. Finally we obtain the result
$`S_I(\omega )=`$ $`S_0+{\displaystyle \frac{(\mathrm{\Delta }I)^2}{W_t}}\left[{\displaystyle \frac{\epsilon ^2}{\mathrm{}^2\mathrm{\Omega }^2}}z_{st}^2\right]{\displaystyle \frac{1}{1+(w/W_t)^2}}`$ (29)
$`+{\displaystyle \frac{2(\mathrm{\Delta }I)^2H^2}{W_0\mathrm{}^2\mathrm{\Omega }^2}}{\displaystyle \frac{1}{1+[(\omega \mathrm{\Omega })/W_0]^2}}`$
where
$`z_{st}={\displaystyle \frac{\epsilon }{\mathrm{}\mathrm{\Omega }}}{\displaystyle \frac{1}{1+4H^2\mathrm{\Gamma }/\gamma _1\mathrm{}^2\mathrm{\Omega }^2}}\mathrm{tanh}({\displaystyle \frac{\mathrm{}\mathrm{\Omega }}{2T}}),`$ (30)
$`W_t=\gamma _1+{\displaystyle \frac{4\mathrm{\Gamma }H^2}{\mathrm{}^2\mathrm{\Omega }^2}},`$ (31)
$`W_0=\gamma _2+{\displaystyle \frac{\mathrm{\Gamma }}{2}}(1+{\displaystyle \frac{\epsilon ^2}{\mathrm{}^2\mathrm{\Omega }^2}}).`$ (32)
Let us emphasize that the effect of finite-temperature environment is not generally equivalent to the nonideality of the detector described by finite $`\gamma `$ in Eq. (6). As an example, in the case of extra environment the right hand part of Eq. (14) should be multiplied by the factor $`1z_{st}^2`$ which disappears ($`z_{st}=0`$) only if $`T=\mathrm{}`$ or $`\epsilon =0`$.
Comparing Eqs. (29) and (17) one can see that within the accuracy of weak coupling approximation the change of $`S_I(\omega )`$ due to extra environment can be reduced to the detector nonideality, $`\eta <1`$, in two cases. If $`|\epsilon /H|1`$, then Eqs. (29) and (17) coincide at arbitrary temperature $`T`$ for $`\eta =(1+2\gamma _2/\alpha \mathrm{\Omega })^1`$. For asymmetric qubit, $`|\epsilon /H|1`$, the equivalence is possible only at high temperatures, $`T\mathrm{}\mathrm{\Omega }`$, and requires conditions
$`\gamma _2=\gamma _1(1+\epsilon ^2/2H^2)/2,`$ (33)
$`\eta ^1=1+(1+\epsilon ^2/4H^2)^{3/2}\gamma _1/\alpha \mathrm{\Omega }.`$ (34)
Figure 6 shows the numerically calculated spectral density $`S_I(\omega )`$ of the detector current for a nonideal detector, $`\eta =0.5`$ (dashed lines) and for an ideal detector but extra coupling of the qubit to the environment at temperature $`T=H`$ (solid lines). The rates $`\gamma _1`$ and $`\gamma _2`$ are chosen according to Eqs. (33) and (34). For the symmetric qubit, $`\epsilon =0`$, the results of two models practically coincide. In contrast, the solid and dashed lines for $`\epsilon =2H`$ significantly differ from each other at low frequencies, while the spectral peak at $`\omega \mathrm{\Omega }`$ is fitted quite well.
## VII Conclusion
Using both Bayesian and conventional approaches, we have calculated the spectral density $`S_I(\omega )`$ of the detector current when a two-level quantum system (qubit) is measured continuously. Depending on the coupling strength, there is a gradual transition from quantum oscillations to quantum jumps. This results in a transition from the peak-like spectral density to the Lorentzian shape of $`S_I(\omega )`$. The maximal height of the peak at the frequency of quantum oscillation is shown to be 4 times the shot noise pedestal (see also Ref. ).
In the simple case of weak coupling between a symmetric qubit and an ideal detector, the height of the spectral peak is twice as high as the classical result for a harmonic signal. In the Bayesian approach this is explained by the significant correlation between the detector noise and the further evolution of the measured system due to quantum back-action. In contrast, in the conventional approach this fact is the consequence of discrete eigenvalues of $`\widehat{z}`$ operator, which corresponds to the magnitude measured by the detector. (In other words, this operator “collapses” the system into one of two eigenstates, and that is why the averaged product of two operators is twice as large as that for a classical harmonic signal.) So, even though the results for $`S_I(\omega )`$ coincide in two approaches, the interpretations are significantly different since the “abrupt” collapse is replaced in the Bayesian approach by the “continuous” collapse related to the noisy detector output.
It is important to notice that the Bayesian formalism allows us to monitor continuously the phase of quantum oscillations. This makes it possible to tune the phase using the feedback control of the qubit parameters. If the real-time calculations using Eqs. (4)–(6) and fast feedback loop were available in an experiment (the typical bandwidth should be larger than $`\mathrm{\Gamma }`$) then the random diffusion of the oscillation phase could be eliminated and the qubit could “stay fresh” for a very long time. The suppression of qubit dephasing using the feedback control of the tunneling strength $`H`$ has been confirmed by the Monte-Carlo simulations. The elimination of the phase diffusion gives rise to the $`\delta `$-function peak in the detector spectral density $`S_I(\omega )`$ at the frequency $`\mathrm{\Omega }`$. A detailed analysis of this situation is beyond the scope of the present paper.
The author thanks D. V. Averin, J. E. Lukens, and K. K. Likharev for fruitful discussions. The work was partly supported by AFOSR.
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# THE b-u SKEWED PARTON DISTRIBUTIONS
## 1 INTRODUCTION
A good theoretical understanding of the $`B\pi `$ transition form factors are of utmost interest. Accurate predictions of these form factors would permit a determination of the less well-known Cabbibo-Kobayashi-Maskawa matrix element $`V_{\mathrm{ub}}`$ from experimental rates of semi-leptonic $`B\pi `$ transitions. The $`B\pi `$ form factors also form an important ingredient of the calculation of $`B\pi \pi `$ decay rates; the factorization hypothesis relates the form factors at maximum recoil to the decay rates. Thus, not surprisingly, the $`B\pi `$ and other heavy-to-light form factors, attracted the attention of many theoreticians. In several of these approaches two distinct dynamical mechanism are considered which build up the $`B\pi `$ form factors: The $`B\pi `$ resonances which control the form factors at small recoil and the overlap of the meson wave functions that dominates at large recoil. Other mechanisms, like the perturbative one, provide only small and often negligible corrections. The crucial questions are how to match these two prominent contributions at intermediate recoil and how to avoid double counting. In a recent article , on which I am going to report here, we proposed to start from skewed parton distributions . The SPDs allow an unambigous superposition of $`B\pi `$ resonances and the overlap contribution.
## 2 $`\mathrm{b}`$-$`\mathrm{u}`$ SKEWED PARTON DISTRIBUTIONS
To be specific let us consider the semi-leptonic decay $`\overline{B}^0\pi ^+\mathrm{}^{}\overline{\nu }_l`$. Instead of the usual form factors, $`F_+`$ and $`F_0`$, (see e.g. ) it is more appropriate here to use the definition
$$\pi ^+;p^{}|\overline{\mathrm{u}}(0)\gamma _\mu \mathrm{b}(0)|\overline{B}^0;p=F^{(1)}(q^2)p_\mu ^{}+F^{(2)}(q^2)\left(q_\mu \frac{q^2}{M_B^2}p_\mu \right),$$
(1)
where $`q=pp^{}`$ and $`M_B`$ ($`M_\pi `$) is the $`B`$ ($`\pi `$) mass. A convenient frame of reference is a generalized Breit frame in which the mesons move collinearly. In this frame the momentum transfer is given by
$$q^2=\zeta M_B^2\left(1\frac{M_\pi ^2}{M_B^2(1\zeta )}\right),$$
(2)
where the so-called skewedness parameter, $`\zeta `$, is defined by the ratio of light-cone plus components
$$\zeta =\frac{q^+}{p^+}=1\frac{p^+}{p^+}.$$
(3)
The skewedness parameter $`\zeta `$ covers the interval $`[0,1M_\pi /M_B]`$ in parallel with the variation of the momentum transfer from zero (the lepton mass is neglected) to $`q_{\mathrm{max}}^2=(M_BM_\pi )^2`$ in the physical region of the $`B\pi `$ transitions. The pion mass is neglected in the calculation of SPDs and form factors. The advantage of the generalized Breit frame is that the $`B\pi `$ matrix element of the current’s plus component is related to the form factor $`F^{(1)}`$ solely while that of the minus component is related to $`F^{(2)}`$. The matrix elements of the transverse currents are zero.
The flavour non-diagonal $`\mathrm{b}`$-$`\mathrm{u}`$ SPDs $`\stackrel{~}{}_\zeta ^{(i)}`$, $`i=1,2`$ are defined by the $`B\pi `$ matrix elements of bilocal products of quark field operators, e.g.
$$\frac{\mathrm{d}z^{}}{2\pi }e^{ixp^+z^{}}\pi ^+;p^{}|\overline{\mathrm{u}}(0)\gamma ^+\mathrm{b}(z^{})|\overline{B}^0;p=(1\zeta )\stackrel{~}{}_\zeta ^{(1)}(x,q^2).$$
(4)
$`x=k^+/p^+`$ is the usual fraction of plus-components of the $`\mathrm{b}`$-quark and $`B`$-meson momenta. The second SPD, $`\stackrel{~}{}_\zeta ^{(2)}`$, is analogously defined with $`\gamma ^+`$ being replaced by $`\gamma ^{}`$. The variable $`q^2`$ is redundant in the generalized Breit frame, see Eq. (2). Integration of (4) over $`x`$ reduces the bilocal operator product to the local one that defines the form factors, see (1). Hence, one has the reduction formula
$$F^{(i)}(q^2)=_{1+\zeta }^1dx\stackrel{~}{}_\zeta ^{(i)}(x)$$
(5)
for $`i=1,2`$.
Depending on the value of $`x`$, the SPDs describe different physical situations:
i) For $`1x\zeta `$ a $`\mathrm{b}`$-quark with momentum fraction $`x`$ is emitted from the $`B`$-meson and a $`\mathrm{u}`$-quark carrying a momentum fraction $`x^{}=k^+{}_{}{}^{}/p^+^{}`$ is absorbed, turning the $`B`$-meson into a pion. This part of the SPDs can be modelled as overlaps of $`B`$ and $`\pi `$ light-cone wave functions. For the valence Fock states, for instance, the overlap reads
$$\stackrel{~}{}_{\zeta \mathrm{ove}}^{(1)}(x)=\frac{2}{1\zeta }\frac{\mathrm{d}^2𝐤_{}}{16\pi ^3}\mathrm{\Psi }_\pi ^{}(x^{}=\frac{x\zeta }{1\zeta },𝐤_{})\mathrm{\Psi }_B(x,𝐤_{}),$$
(6)
where $`𝐤_{}`$ is the intrinsic transverse momentum of the $`\mathrm{b}`$ ($`\mathrm{u}`$) quark with respect to the $`B`$ ($`\pi `$)-meson momentum.
ii) For $`0x<\zeta `$ the fraction $`x^{}`$ is negative. Interpreting a parton with a negative momentum fraction as an antiparton with a positive fraction, one sees that the physical situation is now the emission of a $`\mathrm{b}\overline{\mathrm{u}}`$ pair from $`B`$-meson and the formation of the pion from the remaining partons. This contribution can be desribed by overlaps for $`N+2N`$ parton processes; it is found to be very small numerically. $`B\pi `$ resonances contribute to the SPDs in that region as well . The contribution of the most important resonance, the $`B^{}`$ vector meson, reads
$`\stackrel{~}{}_{\zeta \mathrm{res}}^{(1)}(x)`$ $`=`$ $`{\displaystyle \frac{f_B^{}g_{BB^{}\pi }}{M_B^{}}}\left(M_B^{}^2{\displaystyle \frac{1}{2}}\zeta M_B^2\right){\displaystyle \frac{\varphi _B^{}(x/\zeta )}{M_B^{}^2\zeta M_B^2}},`$
$`\stackrel{~}{}_{\zeta \mathrm{res}}^{(2)}(x)`$ $`=`$ $`{\displaystyle \frac{1}{2}}M_B^2{\displaystyle \frac{f_B^{}g_{BB^{}\pi }}{M_B^{}}}{\displaystyle \frac{\varphi _B^{}(x/\zeta )}{M_B^{}^2\zeta M_B^2}},`$ (7)
where $`\varphi _B^{}(y)`$ is the valence Fock state distribution amplitude of the $`B^{}`$-meson. The argument of the $`B^{}`$ distribution amplitude, $`x/\zeta `$, equals the momentum fraction, $`k_+/q_+`$, which the $`\mathrm{b}`$-quarks carries w.r.t. the $`B^{}`$-meson.
iii) For $`1+\zeta x<0`$ where $`x^{}`$ is negative too, the SPDs describe the emission of a $`\overline{\mathrm{b}}`$-quark and the absorption of a $`\overline{\mathrm{u}}`$ one. Since the probability of finding a $`\mathrm{b}\overline{\mathrm{b}}`$ sea-quark pair in the $`B`$-meson is practically zero, $`\stackrel{~}{}_\zeta ^{(i)}(x)0`$ in this region.
Combining all these contributions, one finds for the $`\mathrm{b}\mathrm{u}`$ SPDs the superposition
$$\stackrel{~}{}_\zeta ^{(i)}(x)=\theta (x\zeta )\stackrel{~}{}_{\zeta \mathrm{ove}}^{(i)}(x)+\theta (\zeta x)\theta (x)\left[\stackrel{~}{}_{\zeta \mathrm{ann}}^{(i)}(x)+\stackrel{~}{}_{\zeta \mathrm{res}}^{(i)}(x)\right].$$
(8)
The relative importance of the overlap contribution to the SPDs on the one side and the sum of annihilation and resonance one on the other side, change with the momentum transfer as a consequence of the relation (2). At large recoil, $`q^20`$, the annihilation and resonance parts do not contribute while they dominate at small recoil, $`q^2q_{\mathrm{max}}^2`$. The superposition (8) is controlled by the skewedness parameter $`\zeta `$ in an unambiguous way, i.e. there is no danger of double counting.
For the numerical estimate of the overlap contribution a simple Gaussian wave function is used to describe the pion’s valence Fock state
$$\mathrm{\Psi }_\pi (x,𝐤_{})=\frac{\sqrt{6}}{f_\pi }\mathrm{exp}\left[\frac{1}{8\pi ^2f_\pi ^2}\frac{𝐤_{}^2}{x(1x)}\right].$$
(9)
$`f_\pi `$ (=132 MeV) is the usual pion decay constant. The wave function (9) has been tested against experiment and found to work satisfactorily in many exclusive reactions involving pions (e.g. ). It is also supported by theoretical studies, e.g. .
For the $`\mathrm{b}\overline{\mathrm{d}}`$ wave function of the $`B`$ meson a slightly modified version of the Bauer-Stech-Wirbel (BSW) function is used
$$\mathrm{\Psi }_B(x,𝐤_{})f_Bx(1x)\mathrm{exp}\left[a_B^2[M_B^2\left(xx_0\right)^2+𝐤_{}^2]\right].$$
(10)
$`m_\mathrm{b}`$ is taken to be 4.8 GeV, $`a_B=1.51\mathrm{GeV}^1`$, $`f_B=180\mathrm{MeV}`$ and the wave function is normalized to unity. The distribution amplitude exhibits a pronounced peak, its position is approximately at $`xx_0=m_\mathrm{b}/M_B`$.
In principle, the overlap parts of the SPDs receive contributions from all Fock states. The generalization of the overlap representation (6) to higher Fock states is a straightforward application of the methods outlined in . Using suitably generalized $`N`$-particle wave functions, one can show that the higher Fock state contributions to the SPD $`\stackrel{~}{}_{\zeta \mathrm{ove}}^{(1)}`$ are very small and can be neglected; they represent power corrections $`(\overline{\mathrm{\Lambda }}/M_B)^{n(N)}`$.
In order to estimate the resonance contribution the same ansatz as for the $`B`$-meson is employed for the $`B^{}`$-meson distribution amplitude. Its explicit form is however irrelevant for the transition form factors. The product of the $`B^{}`$ decay constant, $`f_B^{}`$ and the $`BB^{}\pi `$ coupling constant is taken to be $`20f_B`$ .
The numerical results for the $`\mathrm{b}`$-$`\mathrm{u}`$ SPD $`\stackrel{~}{}_\zeta ^{(1)}`$ are shown in Fig. 1. Both the contributions exhibit charcteristic bumps which are generated by the pronounced peaks in the $`B`$ and $`B^{}`$ distribution amplitudes.
## 3 $`B`$-$`\pi `$ FORM FACTORS
The $`B\pi `$ form factors $`F^{(i)}`$ can be evaluated form the SPDs through the reduction formula (5). Since the form factors $`F_+`$ and $`F_0`$, being linearly related to the $`F^{(i)}`$, are more suitable in applications to decay processes, I only present results for them. In addition to the overlap and resonance contributions the form factors also receive contributions from perturbative QCD where a hard gluon, with a virtuality of the order of $`M_B^2`$, is exchanged between the struck and the spectator quark. In Ref. the perturbative contributions have been evaluated at large recoil within the modified perturbative approach and one can make use of these results. At small recoil the perturbative contributions cease to be reliable because of the small virtualities some of the internal off-shell quarks and gluons acquire in this region.
Numerical results for the form factors are displayed in Fig. 2. In the case of the form factor $`F_+`$ one sees the dominance of the overlap contribution at large recoil while the resonance contribution takes the lead at small recoil (cf. Eq. (8)). The perturbative contribution provides a correction to $`F_+`$ of the order of $`10\%`$ at large recoil and can be neglected at small recoil. The sum of the three contributions to $`F_+`$ is in fair agreement with lattice QCD results . Due to the absence of the $`B^{}`$ pole the form factor $`F_0`$ behaves differently; it is rather flat over the entire range of momentum transfer. The perturbative contribution makes up a substantial fraction of the total result for $`F_0`$ at intermediate momentum transfer. Since, as is mentioned above, it becomes unreliable for $`q^2\stackrel{>}{}18\mathrm{GeV}^2`$ $`F_0`$ cannot reliably be predicted at large $`q^2`$. A calculation of $`F_0`$ in that region would also require a detailed investigation of the scalar $`B\pi `$ resonances of which not much is known at present. Despite of this drawback the results for this form factor are also in fair agreement with the lattice QCD results and, in tendency, seem to extrapolate to the $`B`$-sector analogue of the Callan-Treiman value.
An assessment of the theoretical uncertainties of the predictions for the $`B\pi `$ transition form factors leads to an uncertainty of about $`2025\%`$. It includes estimates of: Sudakov suppressions in the end-point region ($`x1`$), contributions from two-particle twist-three wave functionss, deviations from the asymptotic form of the pion distribution amplitude, uncertainties of the input parameters ($`g_{BB^{}\pi }`$,$`f_B`$, $`f_B^{}`$, $`m_\mathrm{b}`$) and from the order $`\overline{\mathrm{\Lambda }}/M_B`$ corrections, In particular for the form factors at maximum recoil, $`F_+(0)=F_0(0)`$ a value of $`0.22\pm 0.05`$ is obtained.
With the form factors at hand one can evaluate the semi-leptonic decay rates $`\overline{B}^0\pi ^+\mathrm{}^{}\overline{\nu }_{\mathrm{}}`$. For the branching ratio of the light-lepton modes one finds
$$[\overline{B}^0\pi ^+e\overline{\nu }_e][\overline{B}^0\pi ^+\mu \overline{\nu }_\mu ]=1.910^4\left(\frac{|V_{\mathrm{ub}}|}{0.0035}\right)^2.$$
(11)
The theoretical uncertainty of this prediction, dominated by that of the overlap contribution, amounts to about $`30\%`$. This result is to be compared with the CLEO measurement : $`(1.8\pm 0.4\pm 0.3\pm 0.2)10^4`$. For the $`\tau `$ channel one obtains a value of $`1.510^4`$ for the branching ratio ($`|V_{\mathrm{ub}}|=0.0035`$).
The exclusive $`B`$-decays into pairs of pions are usually calculated on the basis of a factorization hypothesis according to which the decay amplitudes can be written as a product of two weak current matrix elements
$$=\frac{G_F}{\sqrt{2}}V_{\mathrm{ud}}^{}V_{\mathrm{ub}}\pi ^{};q|J_W^\mu |0\pi ^+;p^{}|J_\mu ^W|\overline{B}^0;p.$$
(12)
The first matrix element defines the usual pion decay constant ($`f_\pi q^\mu `$) while the second one defines the $`B\pi `$ transition form factors (1). The factorizing contribution alone leads to the following branching ratio
$$(\overline{B}^0\pi ^+\pi ^{})=10.510^6\left(\frac{|V_{\mathrm{ub}}|}{0.0035}\right)^2\left|\frac{F_+(0)}{0.33}\right|^2.$$
(13)
Ignoring the short-distance corrections which seem to amount to about $`1020\%`$ , and choosing $`|V_{\mathrm{ub}}|=0.0035`$, one finds agreement between the prediction for $`F_+(0)`$ from the SPD approach and the recent CLEO measurement for the $`\overline{B}^0\pi ^+\pi ^{}`$ branching ratio of $`(4.3{}_{1.4}{}^{+1.6}\pm 0.5)10^6`$. The experimental value is much smaller than expected (based on $`F_+(0)0.30.33`$) and a revision of the theoretical analysis of exclusive $`B`$-decays seems to be required.
## 4 CONCLUSIONS
The $`\mathrm{b}`$-$`\mathrm{u}`$ SPDs are calculated from light-cone wave function overlaps and a contribution from the $`B^{}`$ resonance. The chief advantage of the SPD approach is that the skewedness parameter clearly separates the overlap from the resonance contribution and both the contributions can be added unambigously. The $`B\pi `$ transition form factors are calculated from the $`\mathrm{b}\mathrm{u}`$ SPDs by means of reduction formulas. $`F_+`$ is obtained in the entire range of momentum transfer and $`F_0`$ up to about 18 GeV<sup>2</sup>. In particular, a value of $`0.22\pm 0.05`$ is found for the form factors at maximum recoil. This value appears to be in agreement with the recent CLEO measurement of $`B\pi \pi `$ decays (if the latter process is analysed on the basis of the factorization hypothesis). The prediction for the total decay for the process $`\overline{B}^0\pi ^+e\overline{\nu }_e`$ is also in agreement with a CLEO measurement . In both the cases, the $`\pi \pi `$ and the semi-leptonic decay, a value of 0.0035 is used for the CKM matrix element $`V_{\mathrm{ub}}`$.
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# BAND DESCRIPTION OF KNOTS AND VASSILIEV INVARIANTS
> Abstract
>
> In 1993 K. Habiro defined $`C_k`$-move of oriented links and around 1994 he proved that two oriented knots are transformed into each other by $`C_k`$-moves if and only if they have the same Vassiliev invariants of order $`k1`$. In this paper we define Vassiliev invariant of type $`(k_1,\mathrm{},k_l)`$, and show that, for $`k=k_1+\mathrm{}+k_l`$, two oriented knots are transformed into each other by $`C_k`$-moves if and only if they have the same Vassiliev invariants of type $`(k_1,\mathrm{},k_l)`$. We introduce a concept ‘band description of knots’ and give a diagram-oriented proof of this theorem. When $`k_1=\mathrm{}=k_l=1`$, the Vassiliev invariant of type $`(k_1,\mathrm{},k_l)`$ coincides with the Vassiliev invariant of order $`l1`$ in the usual sense. As a special case, we have Habiro’s theorem stated above.
2000 Mathematics Subject Classification: 57M25Keywords and Phrases: knot, $`C_n`$-move, Vassiliev invariant, finite type invariant, band description
Introduction
In 1993 K. Habiro defined $`C_k`$-move of oriented links for each natural number $`k`$ . A $`C_k`$-move is a kind of local move of oriented links. Around 1994 he proved that two oriented knots have the same Vassiliev invariants of order $`k1`$ if and only if they are transformed into each other by $`C_k`$-moves. Thus he has succeeded in deducing a geometric conclusion from an algebraic condition. However this theorem appears only in his recent paper . In he develops his original clasper theory and obtains the theorem as a consequence of clasper theory. We note that the ‘if’ part of the theorem is also shown in , , and , and in T. Stanford gives another characterization of knots with the same Vassiliev invariants of order $`k1`$.
In this paper we define Vassiliev invariant of type $`(k_1,\mathrm{},k_l)`$, and show that, for $`k=k_1+\mathrm{}+k_l`$, two oriented knots are transformed into each other by $`C_k`$-moves if and only if they have the same Vassiliev invariants of type $`(k_1,\mathrm{},k_l)`$. When $`k_1=\mathrm{}=k_l=1`$, the Vassiliev invariant of type $`(k_1,\mathrm{},k_l)`$ coincides with the Vassiliev invariant of order $`l1`$ in the usual sense. As a special case, we have Habiro’s theorem. We use a concept ‘band description of knots’ and give a diagram-oriented proof of the theorem. The proof is elementary and completely self-contained. Note that the prototypes of band description appear in , and . In particular in the second author showed that any knot can be expressed as a band sum of a trivial knot and some Borromean rings. The concept of band description is a development of this fact. More generally the authors defined ‘band description of spatial graphs’ in . The concept of Vassiliev invariant of type $`(k_1,\mathrm{},k_l)`$ is defined in . A related result to the case $`k_1=\mathrm{}=k_l=2`$ is shown in .
1. Definitions and Main Result
Throughout this paper we work in the piecewise linear category.
A tangle $`T`$ is a disjoint union of properly embedded arcs in the unit $`3`$-ball $`B^3`$. A tangle $`T`$ is trivial if there exists a properly embedded disk in $`B^3`$ that contains $`T`$. A local move is a pair of trivial tangles $`(T_1,T_2)`$ with $`T_1=T_2`$ such that for each component $`t`$ of $`T_1`$ there exists a component $`u`$ of $`T_2`$ with $`t=u`$. Such a pair of components is called a corresponding pair. Two local moves $`(T_1,T_2)`$ and $`(U_1,U_2)`$ are equivalent, denoted by $`(T_1,T_2)(U_1,U_2)`$, if there is an orientation preserving self-homeomorphism $`\psi :B^3B^3`$ such that $`\psi (T_i)`$ and $`U_i`$ are ambient isotopic in $`B^3`$ relative to $`B^3`$ for $`i=1,2`$. Here $`\psi (T_i)`$ and $`U_i`$ are ambient isotopic in $`B^3`$ relative to $`B^3`$ if $`\psi (T_i)`$ is deformed to $`U_i`$ by an ambient isotopy of $`B^3`$ that is pointwisely fixed on $`B^3`$.
Let $`(T_1,T_2)`$ be a local move, $`t_1`$ a component of $`T_1`$ and $`t_2`$ a component of $`T_2`$ with $`t_1=t_2`$. Let $`N_1`$ and $`N_2`$ be regular neighbourhoods of $`t_1`$ and $`t_2`$ in $`(B^3T_1)t_1`$ and $`(B^3T_2)t_2`$ respectively such that $`N_1B^3=N_2B^3`$. Let $`\alpha `$ be a disjoint union of properly embedded arcs in $`B^2\times [0,1]`$ as illustrated in Fig. 1.1. Let $`\psi _i:B^2\times [0,1]N_i`$ be a homeomorphism with $`\psi _i(B^2\times \{0,1\})=N_iB^3`$ for $`i=1,2`$. Suppose that $`\psi _1(\alpha )=\psi _2(\alpha )`$ and $`\psi _1(\alpha )`$ and $`\psi _2(\alpha )`$ are ambient isotopic in $`B^3`$ relative to $`B^3`$. Then we say that a local move $`((T_1t_1)\psi _1(\alpha ),(T_2t_2)\psi _2(\alpha ))`$ is a double of $`(T_1,T_2)`$ with respect to the components $`t_1`$ and $`t_2`$. Note that a double of $`(T_1,T_2)`$ with respect to $`t_1`$ and $`t_2`$ is well-defined up to equivalence.
Fig. 1.1
A $`C_1`$-move is a local move $`(T_1,T_2)`$ as illustrated in Fig. 1.2. A double of a $`C_k`$-move is called a $`C_{k+1}`$-move. Note that for each natural number $`k`$ there are only finitely many $`C_k`$-moves up to equivalence. We note that by the definition a $`C_k`$-move is Brunnian. That is, if $`(T_1,T_2)`$ is a $`C_k`$-move and $`t_1`$, $`t_2`$ components of $`T_1`$ and $`T_2`$ respectively with $`t_1=t_2`$, then the tangles $`T_1t_1`$ and $`T_2t_2`$ are ambient isotopic in $`B^3`$ relative to $`B^3`$.
Fig. 1.2
Let $`(T_1,T_2)`$ be a local move. Then $`(T_2,T_1)`$ is also a local move. We call $`(T_2,T_1)`$ the inverse of $`(T_1,T_2)`$. It is easy to see that the inverse of a $`C_1`$-move is equivalent to itself. Then it follows inductively that the inverse of a $`C_n`$-move is equivalent to a $`C_n`$-move (but possibly not equivalent to itself).
Let $`K_1`$ and $`K_2`$ be oriented knots in the oriented three-sphere $`S^3`$. We say that $`K_1`$ and $`K_2`$ are related by a local move $`(T_1,T_2)`$ if there is an orientation preserving embedding $`h:B^3S^3`$ such that $`K_ih(B^3)=h(T_i)`$ for $`i=1,2`$ and $`K_1h(B^3)=K_2h(B^3)`$ together with orientations. Then we also say that $`K_2`$ is obtained from $`K_1`$ by an application of $`(T_1,T_2)`$. Two oriented knots $`K_1`$ and $`K_2`$ are $`C_k`$-equivalent if $`K_1`$ and $`K_2`$ are related by a finite sequence of $`C_k`$-moves and ambient isotopies. This relation is an equivalence relation on knots. It is known that $`C_k`$-equivalence implies $`C_{k1}`$-equivalence , see Remark 2.4.
Let $`l`$ be a positive integer and $`k_1,\mathrm{},k_l`$ positive integers. Suppose that for each $`P\{1,\mathrm{},l\}`$ an oriented knot $`K_P`$ in $`S^3`$ is assigned. Suppose that there are orientation preserving embeddings $`h_i:B^3S^3`$ $`(i=1,\mathrm{},l)`$ such that
(1) $`h_i(B^3)h_j(B^3)=\mathrm{}`$ if $`ij`$,
(2) $`K_P_{i=1}^lh_i(B^3)=K_P^{}_{i=1}^lh_i(B^3)`$ together with orientation for any subsets $`P,P^{}\{1,\mathrm{},l\}`$,
(3) $`(h_i^1(K_{\mathrm{}}),h_i^1(K_{\{1,\mathrm{},l\}}))`$ is a $`C_{k_i}`$-move $`(i=1,\mathrm{},l)`$, and
(4) $`K_Ph_i(B^3)=\{\begin{array}{cc}K_{\{1,\mathrm{},l\}}h_i(B^3)\hfill & \text{if }iP,\hfill \\ K_{\mathrm{}}h_i(B^3)\hfill & \text{otherwise}.\hfill \end{array}`$
Then we call the set of knots $`\{K_P|P\{1,\mathrm{},l\}\}`$ a singular knot of type $`(k_1,\mathrm{},k_l)`$. Let $`𝒦`$ be the set of all oriented knot types in $`S^3`$ and $`𝒦`$ the free abelian group generated by $`𝒦`$. We sometimes identify a knot and its knot type without explicite mention. For a singular knot $`K=\{K_P|P\{1,\mathrm{},l\}\}`$ of type $`(k_1,\mathrm{},k_l)`$, we define an element $`\kappa (K)`$ of $`𝒦`$ by
$$\kappa (K)=\underset{P\{1,\mathrm{},l\}}{}(1)^{|P|}K_P.$$
Let $`𝒱(k_1,\mathrm{},k_l)`$ be the subgroup of $`𝒦`$ generated by all $`\kappa (K)`$ where $`K`$ varies over all singular knots of type $`(k_1,\mathrm{},k_l)`$.
Let $`K_1\mathrm{\#}K_2`$ be the composite knot of two knots $`K_1`$ and $`K_2`$. Then $`K_1\mathrm{\#}K_2K_1K_2𝒦`$ is called a composite relator following Stanford . Let $`_\mathrm{\#}`$ be the subgroup of $`𝒦`$ generated by all composite relators.
Let $`\iota :𝒦𝒦`$ be the natural inclusion map. Let $`\pi _{(k_1,\mathrm{},k_l)}:𝒦𝒦/𝒱(k_1,\mathrm{},k_l)`$ and $`\lambda _{(k_1,\mathrm{},k_l)}:𝒦/𝒱(k_1,\mathrm{},k_l)𝒦/(𝒱(k_1,\mathrm{},k_l)+_\mathrm{\#})`$ be the quotient homomorphisms. Then the composite maps $`\pi _{(k_1,\mathrm{},k_l)}\iota :𝒦𝒦/𝒱(k_1,\mathrm{},k_l)`$ and $`\lambda _{(k_1,\mathrm{},k_l)}\pi _{(k_1,\mathrm{},k_l)}\iota :𝒦𝒦/(𝒱(k_1,\mathrm{},k_l)+_\mathrm{\#})`$ are called the universal Vassiliev invariant of type $`(k_1,\mathrm{},k_l)`$ and the universal additive Vassiliev invariant of type $`(k_1,\mathrm{},k_l)`$ respectively. We denote them by $`v_{(k_1,\mathrm{},k_l)}`$ and $`w_{(k_1,\mathrm{},k_l)}`$ respectively.
Since a $`C_1`$-move is a crossing change we have that a singlar knot of type $`(\underset{l}{\underset{}{1,\mathrm{},1}})`$ is essentially the same as a singlar knot with $`l`$ crossing vertices in the usual sense. Therefore we have that $`v_{(1,\mathrm{},1)}`$ is the universal Vassiliev invariant of order $`l1`$ and $`w_{(1,\mathrm{},1)}`$ is the universal additive Vassiliev invariant of order $`l1`$. Note that $`v_{(1,\mathrm{},1)}(K_1)=v_{(1,\mathrm{},1)}(K_2)`$ if and only if $`v(K_1)=v(K_2)`$ for any Vassiliev invariant $`v`$ of order $`l1`$. Similarly $`w_{(1,\mathrm{},1)}(K_1)=w_{(1,\mathrm{},1)}(K_2)`$ if and only if $`w(K_1)=w(K_2)`$ for any additive Vassiliev invariant $`w`$ of order $`l1`$. We also note that $`v_{(2,\mathrm{},2)}`$ is essentially same as that defined in , . In the authors defined a finite type invariant of order $`(k;n)`$, that is essentially same as $`v_{(\underset{k+1}{\underset{}{n1,\mathrm{},n1}})}`$.
Now we state our main results.
Theorem 1.1. Let $`k_1,\mathrm{},k_l`$ be positive integers and $`k=k_1+\mathrm{}+k_l`$. Then $`𝒱(k)𝒱(k_1,\mathrm{},k_l)`$.
Theorem 1.2. Let $`k_1,\mathrm{},k_l`$ be positive integers and $`k=k_1+\mathrm{}+k_l`$. Then $`𝒱(k_1,\mathrm{},k_l)𝒱(k)+_\mathrm{\#}`$.
By Theorems 1.1 and 1.2, we have the following corollary.
Corollary 1.3. Let $`k_1,\mathrm{},k_l`$ be positive integers and $`k=k_1+\mathrm{}+k_l`$. Then $`𝒱(k_1,\mathrm{},k_l)+_\mathrm{\#}=𝒱(k)+_\mathrm{\#}`$. $`\mathrm{}`$
The following theorem was proved by Habiro , and the authors gave a proof by using band description . In section 3, we give the same proof as in for the convenience of the reader.
Theorem 1.4. (Habiro ) The $`C_k`$-equivalence classes of oriented knots in $`S^3`$ forms an abelian group under connected sum of oriented knots.
We denote this group by $`𝒦/C_k`$.
Theorem 1.5. Let $`\eta _k:𝒦/C_k𝒦/(𝒱(k)+_\mathrm{\#})`$ be a map induced by the inclusion $`\iota `$. Then $`\eta _k`$ is an isomorphism.
Let $`K_1`$ and $`K_2`$ be oriented knots and $`k=k_1+\mathrm{}+k_l`$. If $`K_1`$ and $`K_2`$ are $`C_k`$-equivalent, then by Theorem 1.1, $`K_1K_2𝒱(k)𝒱(k_1,\mathrm{},k_l)`$. Therefore we have $`v_{(k_1,\mathrm{},k_l)}(K_1)=v_{(k_1,\mathrm{},k_l)}(K_2)`$. On the other hand, if $`v_{(k_1,\mathrm{},k_l)}(K_1)=v_{(k_1,\mathrm{},k_l)}(K_2)`$, then $`w_{(k_1,\mathrm{},k_l)}(K_1)=w_{(k_1,\mathrm{},k_l)}(K_2)`$. Then by Corollary 1.3, $`K_1K_2𝒱(k)+_\mathrm{\#}`$. Then by Theorem 1.5 we have that $`K_1`$ and $`K_2`$ are $`C_k`$-equivalent. Hence we have the following theorem.
Theorem 1.6. Let $`k_1,\mathrm{},k_l`$ be positive integers and $`k=k_1+\mathrm{}+k_l`$. Let $`K_1`$ and $`K_2`$ be oriented knots in $`S^3`$. Then the following conditions are mutually equivalent.
(1) $`K_1`$ and $`K_2`$ are $`C_k`$-equivalent,
(2) $`v_{(k_1,\mathrm{},k_l)}(K_1)=v_{(k_1,\mathrm{},k_l)}(K_2)`$,
(3) $`w_{(k_1,\mathrm{},k_l)}(K_1)=w_{(k_1,\mathrm{},k_l)}(K_2)`$. $`\mathrm{}`$
As a special case of Theorem 1.6 we have the following theorem.
Theorem 1.7. (Habiro ) Two oriented knots $`K_1`$ and $`K_2`$ are $`C_k`$-equivalent if and only if their values of the universal (additive) Vassiliev invariant of order $`k1`$ are equal. $`\mathrm{}`$
The remainder of this paper, we prove Theorems 1.1, 1.2, 1.4 and 1.5. We give a proof of Theorem 1.1 in section 2 and proofs of Theorems 1.2, 1.4 and 1.5 in section 3. The reader who wishes may change the order of sections 2 and 3 since they are independent except for Remark 2.4.
2. $`C_k`$-moves
A graph is called a tree if it is connected and simply connected as a topological space. A graph is uni-trivalent if each vertex has degree one or three. Let $`G`$ be a uni-trivalent tree embedded on the unit disk $`D^2`$ such that $`GD^2`$ is exactly the set of degree-one vertices $`\{v_1,\mathrm{},v_{k+1}\}`$ of $`G`$. Suppose that an edge $`e`$ of $`G`$ is specified. We will assign a $`C_k`$-move to $`G`$ with respest to $`e`$ and a pair of corresponding components of it to each $`v_i`$ as follows. If $`G`$ is a tree with just two vertices $`v_1,v_2`$ and one edge $`e`$ joining them then we assign a diagram of a $`C_1`$-move as illustrated in Fig. 2.1. Suppose that for each uni-trivalent tree on $`k`$ vertices with a specified edge, a diagram of a $`C_{k1}`$-move is assigned. Suppose that $`v_k`$ and $`v_{k+1}`$ are degree-one vertices of $`G`$ such that there is a degree-three vertex $`u`$ of $`G`$ that is adjacent to both of $`v_k`$ and $`v_{k+1}`$. Suppose that neither $`uv_k`$ nor $`uv_{k+1}`$ is a specified edge $`e`$. Let $`G^{}`$ be a uni-trivalent tree obtained from $`G`$ by deleting $`v_{k+1}`$ and the edge $`uv_{k+1}`$ and forgetting $`u`$. Let $`𝒟^{}`$ be a diagram of a $`C_{k1}`$-move assigned to $`G^{}`$ with respect to $`e`$. Let $`t_1`$ and $`t_2`$ be a pair of corresponding components assigned to $`v_k`$. Then we replace their parts in $`D^{}`$ as the diagram illustrated in Fig. 2.2. We assign the new pairs of corresponding components to $`v_k`$ and $`v_{k+1}`$ respecting the cyclic order of them on $`D^2`$. Thus we have assigned a $`C_k`$-move to $`G`$ with respect to the specified edge. See for example Fig. 2.3. Note that any $`C_k`$-move is assigned to a uni-trivalent tree with respect to a specified edge up to equivalence of local moves.
Fig. 2.1
Fig. 2.2
Fig. 2.3
Lemma 2.1. Let $`(T_1,T_2)`$ be a $`C_k`$-move assigned to a uni-trivalent tree $`G`$ with respect to a specified edge $`e`$. Suppose that an edge $`e^{}`$ is incident to $`e`$. Then there is a re-embedding $`f:GD^2`$ such that the $`C_k`$-move assigned to $`f(G)`$ with respect to $`f(e^{})`$ is equivalent to $`(T_1,T_2)`$.
Proof. Let $`v`$ be the common vertex of $`e`$ and $`e^{}`$. First suppose that only $`v`$ is the degree-three vertex of $`G`$. Then we have the result by the deformation illustrated in Fig. 2.4. By taking appropriate doubles we have the general case. $`\mathrm{}`$
Fig. 2.4
A path of $`G`$ is a subgraph of $`G`$ homeomorphic to a closed interval. Let $`\sigma (G)`$ be the maximal of the number of the edges of a path of $`G`$ and call it the diameter of $`G`$. We say that a $`C_k`$-move $`(T_1,T_2)`$ is one-branched if it is assigned to a uni-trivalent tree of diameter $`k`$. Note that if $`G`$ is a uni-trivalent tree with $`k+1`$ degree-one vertices and $`\sigma (G)=k`$, then all of the degree-three vertices lie on the path with $`k`$ edges.
Let $`S_1`$ and $`S_2`$ be tangles. We say that $`S_1`$ and $`S_2`$ are related by a local move $`(T_1,T_2)`$ if there is an orientation preserving embedding $`h:B^3\mathrm{int}B^3`$ such that $`S_ih(B^3)=h(T_i)`$ for $`i=1,2`$ and $`S_1h(B^3)=S_2h(B^3)`$. Then we say that $`S_2`$ is obtained from $`S_1`$ by an application of $`(T_1,T_2)`$.
The following result was shown by Habiro . Since the article is written in Japanese, we give a proof by using our terms.
Lemma 2.2. (Habiro ) Let $`(T_1,T_2)`$ be a $`C_k`$-move. Then $`T_1`$ and $`T_2`$ are related by a finite sequence of one-branched $`C_k`$-moves and ambient isotopies relative to $`B^3`$.
Proof. Suppose that $`(T_1,T_2)`$ is assigned to a uni-trivalent tree $`G`$. If $`\sigma (G)=k`$ then $`(T_1,T_2)`$ itself is a one-branched $`C_k`$-move. Therefore we may suppose that $`\sigma (G)<k`$. It is sufficient to show that $`T_1`$ and $`T_2`$ are related by $`C_k`$-moves each of which is assigned to a uni-trivalent tree of diameter $`\sigma (G)+1`$. Let $`P`$ be a path of $`G`$ containing $`\sigma (G)`$ edges. Since $`\sigma (G)<k`$, there are degree-three vertices of $`G`$ that are not on $`P`$. Let $`v`$ be one of them such that $`v`$ is adjacent to a vertex $`w`$ on $`P`$. Let $`w_1`$ and $`w_2`$ be the vertices on $`P`$ adjacent to $`w`$. Let $`v_1`$ and $`v_2`$ be the other vertices adjacent to $`v`$. By Lemma 2.1 we may assume that the specified edge is $`ww_1`$. We temporarily forget the embedding of $`G`$ into $`D^2`$ and define two abstract graphs $`G_1`$ and $`G_2`$ from $`G`$ as follows. Let $`G_i`$ be a uni-trivalent tree obtained from $`G`$ by deleting the edge $`vv_i`$, forgetting $`v`$, adding a vertex $`u`$ on $`ww_2`$ and adding an edge $`uv_i`$ $`(i=1,2)`$. Then we have $`\sigma (G_1)=\sigma (G_2)=\sigma (G)+1`$. We will show that $`T_1`$ and $`T_2`$ are related by a $`C_k`$-move assigned to some embedding of $`G_1`$ and a $`C_k`$-move assigned to some embedding of $`G_2`$. In the following we consider the simplest case illustrated in Fig. 2.5. General cases follow immediately by taking appropriate doubles. The $`C_k`$-move of Fig. 2.5 is assigned to the graph of Fig. 2.5 with respect to $`ww_1`$. Note that the $`C_k`$-move of Fig. 2.5 is equivalent to a local move of Fig. 2.6. Then it is realized by two local moves Fig. 2.7 (a) and (b). Then each of them is equivalent to a desired $`C_k`$-move. Fig. 2.8 indicates it for the case (a). The case (b) is similar and we omit it. $`\mathrm{}`$
Fig. 2.5
Fig. 2.6
Fig. 2.7
Fig. 2.8
Corollary 2.3. (Habiro ) Two oriented knots $`K_1`$ and $`K_2`$ are $`C_k`$-equivalent if and only if they are related by a finite sequence of one-branched $`C_k`$-moves and ambient isotopies. $`\mathrm{}`$
Let $`T_0`$ be a $`j`$-component tangle. Let $`𝒯(T_0)`$ be the set of the tangles each element of which is homotopic to $`T_0`$ relative to $`B^3`$. Let $`𝒯(T_0)`$ be the free abelian group generated by $`𝒯(T_0)`$. Let $`k_1,\mathrm{},k_l`$ be natural numbers. We define a singular tangle of type $`(k_1,\mathrm{},k_l)`$ and a subgroup $`𝒱(k_1,\mathrm{},k_l)(T_0)`$ of $`𝒯(T_0)`$ as follows. Suppose that for each $`P\{1,\mathrm{},l\}`$ a tangle $`T_P𝒯(T_0)`$ is assigned. Suppose that there are orientation preserving embeddings $`h_i:B^3\mathrm{int}B^3`$ $`(i=1,\mathrm{},l)`$ such that
(1) $`h_i(B^3)h_j(B^3)\mathrm{}`$ if $`ij`$,
(2) $`T_P_{i=1}^lh_i(B^3)=T_P^{}_{i=1}^lh_i(B^3)`$ for any subsets $`P,P^{}\{1,\mathrm{},l\}`$,
(3) $`(h_i^1(T_{\mathrm{}}),h_i^1(T_{\{1,\mathrm{},l\}}))`$ is a $`C_{k_i}`$-move $`(i=1,\mathrm{},l)`$, and
(4) $`T_Ph_i(B^3)=\{\begin{array}{cc}T_{\{1,\mathrm{},l\}}h_i(B^3)\hfill & \text{if }iP,\hfill \\ T_{\mathrm{}}h_i(B^3)\hfill & \text{otherwise}.\hfill \end{array}`$
Then we call the set of tangles $`\{T_P|P\{1,\mathrm{},l\}\}`$ a singular tangle of type $`(k_1,\mathrm{},k_l)`$. For a singular tangle $`T=\{T_P|P\{1,\mathrm{},l\}\}`$ of type $`(k_1,\mathrm{},k_l)`$, we define an element $`\kappa (T)`$ of $`𝒯(T_0)`$ by
$$\kappa (T)=\underset{P\{1,\mathrm{},l\}}{}(1)^{|P|}T_P.$$
Let $`𝒱(k_1,\mathrm{},k_l)(T_0)`$ be the subgroup of $`𝒯(T_0)`$ generated by all $`\kappa (T)`$ where $`T`$ varies over all singular tangles of type $`(k_1,\mathrm{},k_l)`$.
Proof of Theorem 1.1. By Corollary 2.3 it is sufficient to show that if $`(T_1,T_2)`$ is a one-branched $`C_k`$-move, then $`T_1T_2𝒱(k_1,\mathrm{},k_l)(T_1)`$. We first show that $`T_1T_2𝒱(kk_l,k_l)(T_1)`$. By Lemma 2.1 we have that any one-branched $`C_k`$-move is equivalent to a move illustrated in Fig. 2.9 (a) or (b). Therefore we may assume that $`T_1`$ and $`T_2`$ are as illustrated in Fig. 2.9 (a) or (b). Suppose that $`T_1`$ and $`T_2`$ are as in Fig. 2.9 (a) (resp. (b)). Let $`T_3,T_4,T_5`$ and $`T_6`$ be tangles as illustrated in Fig. 2.10 (a) (resp. (b)). Note that $`T_3`$ and $`T_4`$ are ambient isotopic relative to $`B^3`$. Then we have that $`T_1T_2=(T_1T_5T_3+T_6)+(T_5T_2T_6+T_4)`$. Note that $`T_1`$ and $`T_3`$, $`T_2`$ and $`T_4`$, and $`T_5`$ and $`T_6`$ are related by a (one-branched) $`C_{k_l}`$-move. Similarly $`T_1`$ and $`T_5`$, $`T_3`$ and $`T_6`$, $`T_5`$ and $`T_2`$, and $`T_6`$ and $`T_4`$ are related by a one-branched $`C_{kk_l}`$-move. It is not hard to see that both $`T=\{T_1,T_5,T_3,T_6\}`$ and $`T^{}=\{T_5,T_2,T_6,T_4\}`$ are singular tangles of type $`(kk_l,k_l)`$ and that $`\epsilon \kappa (T)=T_1T_5T_3+T_6`$ and $`\epsilon ^{}\kappa (T^{})=T_5T_2T_6+T_4`$, where $`\epsilon =\pm 1`$, and $`\epsilon ^{}=\pm 1`$. Thus $`T_1T_2𝒱(kk_l,k_l)(T_1)`$.
Similarly we have that if $`(T_1^{},T_2^{})`$ is a one-branched $`C_{kk_l}`$-move, then $`T_1^{}T_2^{}𝒱(kk_lk_{l1},k_{l1})(T_1)`$. By substituting this to one-branched $`C_{kk_l}`$-moves in $`B^3`$ with respect to $`T_1`$ and $`T_5`$, $`T_3`$ and $`T_6`$, $`T_5`$ and $`T_2`$, and $`T_6`$ and $`T_4`$, we have that $`T_1T_2𝒱(kk_lk_{l1},k_{l1},k_l)(T_1)`$. Repeating this argument we finally have the desired conclusion. $`\mathrm{}`$
Fig. 2.9
Fig. 2.10
Remark 2.4. In Proof of Theorem 1.1, since $`T_1`$ and $`T_5`$, and $`T_5`$ and $`T_2`$ are related by a one-branched $`C_{kk_l}`$-move, a one-branced $`C_k`$-move is realized by twice applications of one-branched $`C_{kk_l}`$-moves. By Lemma 2.2, for any positive integer $`k,k^{}`$ $`(k^{}<k)`$, a $`C_k`$-move is realized by finitely many $`C_k^{}`$-moves. Hence $`C_k`$-equivalence implies $`C_k^{}`$-equivalence.
3. Band description
A $`C_1`$-link model is a pair $`(\alpha ,\beta )`$ where $`\alpha `$ is a disjoint union of properly embedded arcs in $`B^3`$ and $`\beta `$ is a disjoint union of arcs on $`B^3`$ with $`\alpha =\beta `$ as illustrated in Fig. 3.1. Suppose that a $`C_k`$-link model $`(\alpha ,\beta )`$ is defined where $`\alpha `$ is a disjoint union of $`k+1`$ properly embedded arcs in $`B^3`$ and $`\beta `$ is a disjoint union of $`k+1`$ arcs on $`B^3`$ with $`\alpha =\beta `$ such that $`\alpha \beta `$ is a disjoint union of $`k+1`$ circles. Let $`\gamma `$ be a component of $`\alpha \beta `$ and $`W`$ a regular neighbourhood of $`\gamma `$ in $`(B^3(\alpha \beta ))\gamma `$. Let $`V`$ be an oriented solid torus, $`D`$ a disk in $`V`$, $`\alpha _0`$ properly embedded arcs in $`V`$ and $`\beta _0`$ arcs on $`D`$ as illustrated in Fig. 3.2. Let $`\psi :VW`$ be an orientation preserving homeomorphism such that $`\psi (D)=WB^3`$ and $`\psi (\alpha _0\beta _0)`$ bounds disjoint disks in $`B^3`$. Then we call the pair $`((\alpha \gamma )\psi (\alpha _0),(\beta \gamma )\psi (\beta _0))`$ a $`C_{k+1}`$-link model. We also say that the pair $`((\alpha \gamma )\psi (\alpha _0),(\beta \gamma )\psi (\beta _0))`$ is a double of $`(\alpha ,\beta )`$ with respect to the component $`\gamma `$. A link model is a $`C_k`$-link model for some $`k`$.
Fig. 3.1
Fig. 3.2
Let $`(\alpha _1,\beta _1),\mathrm{},(\alpha _l,\beta _l)`$ be link models. Let $`K`$ be an oriented knot (resp. a tangle). Let $`\psi _i:B^3S^3`$ (resp. $`\psi _i:B^3\mathrm{int}B^3`$) be an orientation preserving embedding for $`i=1,\mathrm{},l`$ and $`b_{1,1},b_{1,2},\mathrm{},b_{1,\rho (1)},b_{2,1},b_{2,2},\mathrm{},b_{2,\rho (2)},\mathrm{},b_{l,1},b_{l,2},\mathrm{},b_{l,\rho (l)}`$ mutually disjoint disks embedded in $`S^3`$ (resp. $`B^3`$). Suppose that they satisfy the following conditions;
(1) $`\psi _i(B^3)\psi _j(B^3)=\mathrm{}`$ if $`ij`$,
(2) $`\psi _i(B^3)K=\mathrm{}`$ for each $`i`$,
(3) $`b_{i,k}K=b_{i,k}K`$ is an arc for each $`i,k`$,
(4) $`b_{i,k}(_{j=1}^l\psi _j(B^3))=b_{i,k}\psi _i(B^3)`$ is a component of $`\psi _i(\beta _i)`$ for each $`i,k`$,
(5) ($`_{k=1}^{\rho (i)}b_{i,k})\psi _i(B^3)=\psi _i(\beta _i)`$ for each $`i`$.
Let $`J`$ be an oriented knot (resp. a tangle) defined by
$$J=K(\underset{i,k}{}b_{i,k})(\underset{i=1}{\overset{l}{}}\psi _i(\alpha _i))\underset{i,k}{}\mathrm{int}(b_{i,k}K)\underset{i=1}{\overset{l}{}}\psi _i(\mathrm{int}\beta _i),$$
where the orientation of $`J`$ coincides that of $`K`$ on $`K_{i,k}b_{i,k}`$ if $`K`$ is oriented. Then we say that $`J`$ is a band sum of $`K`$ and link models $`(\alpha _1,\beta _1),\mathrm{},(\alpha _l,\beta _l)`$. We call each $`b_{i,k}`$ a band. Each image $`\psi _i(B^3)`$ is called a link ball. In particular if $`(\alpha _i,\beta _i)`$ is a $`C_k`$-link model then $`b_{i,k}`$ is called a $`C_k`$-band and $`\psi _i(B^3)`$ is called a $`C_k`$-link ball. We set $`_i=((\alpha _i,\beta _i),\psi _i,\{b_{i,1},\mathrm{},b_{i,\rho (i)}\})`$ and call $`_i`$ a chord. In particular $`_i`$ is called a $`C_k`$-chord when $`(\alpha _i,\beta _i)`$ is a $`C_k`$-link model. We denote $`J`$ by $`J=\mathrm{\Omega }(K;\{_1,\mathrm{},_l\})`$ and call it a band description of $`J`$. We also say that $`J`$ is a band sum of $`K`$ and chords $`_1,\mathrm{},_l`$.
Sublemma 3.1. Let $`(T_1,T_2)`$ be a $`C_k`$-move. Then there is a $`C_k`$-link model $`(\alpha ,\beta )`$ such that $`(T_1,T_2)(\alpha ,\widehat{\beta })`$ where $`\widehat{\beta }`$ is a slight push in of $`\beta `$.
Proof. It is clearly true for $`k=1`$. Suppose that $`(T_1,T_2)`$ is a double of a $`C_{k1}`$-move $`(U_1,U_2)`$ with respect to the components $`u_1`$ and $`u_2`$ and $`(\alpha ^{},\beta ^{})`$ is a $`C_{k1}`$-link model such that $`(U_1,U_2)(\alpha ^{},\widehat{\beta ^{}})`$. Then by the deformation illustrated in Fig. 3.3, we have that there is a double $`(\alpha ,\beta )`$ of $`(\alpha ^{},\beta ^{})`$ with respect to the component which corresponds to $`u_1`$ and $`u_2`$ such that $`(\alpha ,\widehat{\beta })`$ is equivalent to a double of $`(\alpha ^{},\widehat{\beta ^{}})`$ and therefore equivalent to $`(T_1,T_2)`$. $`\mathrm{}`$
Fig. 3.3
Fig. 3.3 indicates that for a double $`(\alpha ,\beta )`$ of $`(\alpha ^{},\beta ^{})`$, $`(\alpha ,\widehat{\beta })`$ is equivalent to a double of $`(\alpha ^{},\widehat{\beta ^{}})`$. Thus we have
Sublemma 3.2. If $`(\alpha ,\beta )`$ is a $`C_k`$-link model, then $`(\alpha ,\widehat{\beta })`$ is equivalent to a $`C_k`$-move. $`\mathrm{}`$
Sublemma 3.3. Let $`(T_1,T_2)`$ be a $`C_k`$-move. Then there is a $`C_k`$-link model $`(\alpha ,\beta )`$ such that a band sum of $`T_1`$ and $`(\alpha ,\beta )`$ is ambient isotopic to $`T_2`$ relative to $`B^3`$.
Proof. Since the inverse $`(T_2,T_1)`$ is also a $`C_k`$-move we have by Sublemma 3.1 that there is a $`C_k`$-link model $`(\alpha ,\beta )`$ such that $`(T_2,T_1)(\alpha ,\widehat{\beta })`$. It is easy to see that $`\alpha `$ is a band sum of $`\widehat{\beta }`$ and $`(\alpha ,\beta )`$ up to ambient isotopy relative to $`B^3`$. Therefore we have the result. $`\mathrm{}`$
From now on we consider knots up to ambient isotopy of $`S^3`$ and tangles up to ambient isotopy of $`B^3`$ relative to $`B^3`$ without explicit mention. As an immeditate consequence of Sublemmas 3.2 and 3.3 we have
Sublemma 3.4. Let $`K`$ and $`J`$ be oriented knots. Then $`K`$ and $`J`$ are related by a $`C_k`$-move if and only if $`J`$ is a band sum of $`K`$ and a $`C_k`$-link model. $`\mathrm{}`$
Let $`K`$ be a knot and $`J=\mathrm{\Omega }(K;\{_1,\mathrm{},_l\})`$ a band sum of $`K`$ and $`C_{k_i}`$-chords $`_i`$ $`(i=1,\mathrm{},l)`$. We define an element $`\kappa (J)`$ of $`𝒦`$ by
$$\kappa (J)=\underset{P\{1,\mathrm{},l\}}{}(1)^{|P|}\mathrm{\Omega }(K;\underset{iP}{}\{_i\}).$$
By Subemmas 3.2 and 3.3 we have that the subgroup $`𝒱(k_1,\mathrm{},k_l)`$ of $`𝒦`$ is generated by all $`\kappa (J)`$ where $`J`$ varies over all band sums of knots and $`C_{k_i}`$-chords $`(i=1,\mathrm{},l)`$.
Sublemma 3.5. Let $`K`$, $`J`$ and $`I`$ be oriented knots. Suppose that $`J=\mathrm{\Omega }(K;\{_1,\mathrm{},_l\})`$ for some chords $`_1,\mathrm{},_l`$ and $`I=\mathrm{\Omega }(J;\{\})`$ for some $`C_k`$-chord $``$. Then there is a $`C_k`$-chord $`^{}`$ such that $`I=\mathrm{\Omega }(K;\{_1,\mathrm{},_l,^{}\})`$. Moreover, if there is a subset $`P`$ of $`\{1,\mathrm{},l\}`$ such that the link ball and the bands of $``$ intersect neither the link ball nor the bands of $`_j`$ for any $`j\{1,\mathrm{},l\}P`$, then $`\mathrm{\Omega }(\mathrm{\Omega }(K;_{iP}\{_i\});\{\})=\mathrm{\Omega }(K;(_{iP}\{_i\})\{^{}\})`$.
Proof. If the bands and the link ball of $``$ are disjoint from those of $`_1,\mathrm{},_l`$ then we have that $`I=\mathrm{\Omega }(K;\{_1,\mathrm{},_l,\})`$. If not then we deform $`I`$ up to ambient isotopy as follows. First we thin and shrink the bands and the link ball of $``$ so that they are thin enough and small enough respectively. If the link ball of $``$ intersects the bands and the link balls of $`_1,\mathrm{},_l`$ then we move the link ball of $``$ so that they does not intersect. Then we slide the bands of $``$ along $`J`$ so that the intersection of the bands with $`J`$ is disjoint from the bands and the link balls of $`_1,\mathrm{},_l`$. Then we sweep the bands of $``$ out of the link balls of $`_1,\mathrm{},_l`$. Note that this is always possible since the tangles are trivial. Finally we sweep the intersection of the bands of $``$ and the bands of $`_1,\mathrm{},_l`$ out of the intersection of the bands of $`_1,\mathrm{},_l`$ and $`K`$. Let $`^{}`$ be the result of the deformation of $``$ described above. Then it is not hard to see that $`^{}`$ is a desired chord. $`\mathrm{}`$
By repeated applications of Sublemmas 3.4 and 3.5 we immediately have the following lemma.
Lemma 3.6. Let $`k`$ be a positive integer and let $`K`$ and $`J`$ be oriented knots. Then $`K`$ and $`J`$ are $`C_k`$-equivalent if and only if $`J`$ is a band sum of $`K`$ and some $`C_k`$-link models. $`\mathrm{}`$
As in the definition of $`C_k`$-move we define iteratedly doubled strings as follows. A $`0`$-double pattern is $`\{𝐨\}\times [0,1]`$ in $`B^2\times [0,1]`$ where $`𝐨`$ is the center of $`B^2`$. Suppose that a $`k`$-double pattern $`A`$ in $`B^2\times [0,1]`$ is defined. Let $`N`$ be a regular neighbourhood of a component $`\gamma `$ of $`A`$ in $`(B^2\times [0,1]A)\gamma `$ such that $`N(B^2\times [0,1])=\mathrm{}`$. Let $`\psi :B^2\times [0,1]N`$ be a homeomorphism with $`\psi (B^2\times \{0,1\})=N(B^2\times \{0,1\})`$. Let $`\alpha `$ be a disjoint union of properly embedded arcs in $`B^2\times [0,1]`$ as illustrated in Fig. 1.1. Then $`(A\gamma )\psi (\alpha )`$ is called a $`(k+1)`$-double pattern. Let $`N`$ be a regular neighbourhood of a properly embedded arc in $`B^3`$. Let $`\psi :B^2\times [0,1]N`$ be a homeomorphism. Then the image $`\psi (A)`$ of a $`k`$-double pattern $`A`$ is called a $`k`$-double string. Note that a ‘crossing change’ between $`k`$-double strng and $`j`$-double string is equivalent to a $`C_{k+j+1}`$-move.
Sublemma 3.7. Let $`(\alpha ,\beta )`$ be a $`C_k`$-link model and $`\gamma `$ a component of $`\alpha \beta `$. Let $`D`$ be a disk in $`B^3`$ such that $`D=\gamma `$ and $`\mathrm{int}DB^3=\mathrm{}`$. Let $`\delta `$ be a properly embedded unknotted arc in $`B^3`$ that intersects $`D`$ transversally at one point in $`\mathrm{int}D`$. Let $`N`$ be a regular neighbourhood of $`\delta `$ in $`B^3\gamma `$. Then there exist a $`(k1)`$-double string $`𝒟`$ in $`N`$ and an orientation preserving homeomorphism $`\phi :B^3B^3`$ such that $`\phi |_{\gamma \beta }=\mathrm{id}|_{\gamma \beta }`$, and $`\phi (\alpha \gamma )=𝒟\gamma `$.
In the lemma above, we note that $`𝒟(\alpha \gamma )`$ is a $`k`$-double string.
Proof. The case $`k=1`$ is clear. Suppose that it is shown for $`k1`$. Let $`(\alpha ,\beta )`$ be a double of a $`C_{k1}`$-link model $`(\alpha ^{},\beta ^{})`$ with respect to the component $`\gamma ^{}`$ of $`\alpha ^{}\beta ^{}`$.
If $`\gamma `$ is already a component of $`\alpha ^{}\beta ^{}`$, then by the deformation of $`\alpha `$ in a regular neighbourhood of $`\gamma ^{}`$ as illustrated in Fig. 3.4 and by the assumption, we have the result.
Next suppose that $`\gamma `$ is not a component of $`\alpha ^{}\beta ^{}`$. In other words $`\gamma `$ is contained in a regular neighbourhood $`W`$ of $`\gamma ^{}`$. Let $`D^{}`$ be a disk in $`B^3`$ such that $`D^{}=\gamma ^{}`$ and $`\mathrm{int}D^{}B^3=\mathrm{}`$. Let $`\delta ^{}`$ be a properly embedded unknotted arc in $`B^3`$ that intersects $`D^{}`$ transversally at one point in $`\mathrm{int}D^{}`$. Let $`N^{}`$ be a regular neighbourhood of $`\delta ^{}`$ in $`B^3\gamma ^{}`$. Then by the assuption there is a $`(k2)`$-double string $`𝒟^{}`$ in $`N^{}`$ and an orientation preserving homeomorphism $`\phi ^{}:B^3B^3`$ such that $`\phi ^{}|_{\gamma ^{}\beta ^{}}=\mathrm{id}|_{\gamma ^{}\beta ^{}}`$ and $`\phi ^{}(\alpha ^{}\gamma ^{})=𝒟^{}\gamma ^{}`$. By modifying $`\phi ^{}`$ if necessary we may suppose that $`\phi ^{}|_{WB^3}=\mathrm{id}|_{WB^3}`$. Then we have that $`\phi ^{}(\alpha \gamma )`$ is as illustrated in Fig. 3.5 (a). Then by the deformation illustrated in Fig. 3.5 we have the result. $`\mathrm{}`$
Fig. 3.4
Fig. 3.5
Lemma 3.8. Let $`K`$, $`J=\mathrm{\Omega }(K;\{_1,\mathrm{},_l,_0\})`$ and $`I=\mathrm{\Omega }(K;\{_1,\mathrm{},_l,_0^{}\})`$ be oriented knots, where $`_1,\mathrm{},_l`$ are chords and $`_0,_0^{}`$ are $`C_k`$-chords. Suppose that $`J`$ and $`I`$ differ locally as illustrated in Fig $`3.6`$ (a), (b), i.e., $`I`$ is obtained from $`J`$ by a crossing change between $`K`$ and a band of $`_0`$. Then $`J`$ and $`I`$ are related by a $`C_{k+1}`$-move. Moreover, there is a $`C_{k+1}`$-chord $``$ such that $`\mathrm{\Omega }(K;(_{iP}\{_i\})\{_0\})=\mathrm{\Omega }(K;(_{iP}\{_i\})\{_0^{},\})`$ for any subset $`P`$ of $`\{1,\mathrm{},l\}`$.
Proof. By shrinking the band and pulling the link ball as illustrated in Fig. 3.6 it is sufficient to show the case that $`K`$ is near the link ball. Then by Sublemma 3.7 we deform the strings in the link ball without disturbing a neighbourhood of the band. Then the crossing change is realized by a $`C_{k+1}`$-move. See Fig. 3.7. In Fig. 3.7, there is a $`3`$-ball $`B`$ in $`S^3`$ and a homeomorphism $`h:BB^3`$ such that $`(h(J),h(I))`$ is a $`C_{k+1}`$-move and $`B`$ is disjoint from the link ball and the bands of any chord $`_i`$ $`(i=1,\mathrm{},l)`$. Thus by Sublemmas 3.3 and 3.5, we have the latter assertion. $`\mathrm{}`$
Fig. 3.6
Fig. 3.7
Lemma 3.9. Let $`K`$, $`J=\mathrm{\Omega }(K;\{_1,\mathrm{},_l,_{0j},_{0k}\})`$ and $`I=\mathrm{\Omega }(K;\{_1,\mathrm{},_l,_{0j}^{},_{0k}^{}\})`$ be oriented knots, where $`_1,\mathrm{},_l`$ are chords and $`_{0j},_{0j}^{}`$ $`(`$resp. $`_{0k},_{0k}^{})`$ are $`C_j`$-chords $`(`$resp. $`C_k`$-chords$`)`$. Suppose that $`J`$ and $`I`$ differ locally as illustrated in Fig. 3.8. Then $`J`$ and $`I`$ are related by a $`C_{j+k}`$-move. Moreover, there is a $`C_{j+k}`$-chord $``$ such that $`\mathrm{\Omega }(K;(_{iP}\{_i\})\{_{0j},_{0k}\})=\mathrm{\Omega }(K;(_{iP}\{_i\})\{_{0j}^{},_{0k}^{},\})`$ for any subset $`P`$ of $`\{1,\mathrm{},l\}`$.
We call the change from $`J`$ to $`I`$ in Lemma 3.9 a band exchange.
Fig. 3.8
Proof. First we deform the strings in the link balls as stated in Sublemma 3.7 then we slide one of the two bands along the other band and then perform a $`C_{j+k}`$-move. See Fig. 3.9. In Fig. 3.9, there is a $`3`$-ball $`B`$ in $`S^3`$ and a homeomorphism $`h:BB^3`$ such that $`(h(J),h(I))`$ is a $`C_{j+k}`$-move and $`B`$ is disjoint from the link ball and the bands of any chord $`_i`$ $`(i=1,\mathrm{},l)`$. Thus by Sublemmas 3.3 and 3.5, we have the latter assertion. $`\mathrm{}`$
Fig. 3.9
Proof of Theorem 1.2. Let $`k_1,\mathrm{},k_l`$ $`(l2)`$ be positive integer and $`k=k_1+\mathrm{}+k_l`$. Let $`K_0`$ be a knot and $`K_1`$ a band sum of $`K_0`$ and $`C_{k_j}`$-chords $`_{k_j,j}`$ $`(j=1,\mathrm{},l)`$. It is sufficient to show that
$$\mu _k\left(\underset{P\{1,\mathrm{},l\}}{}(1)^{|P|}\mathrm{\Omega }(K_0;\underset{jP}{}\{_{k_j,j}\})\right)=0𝒦/(𝒱(k)+_\mathrm{\#}),$$
where $`\mu _k:𝒦𝒦/(𝒱(k)+_\mathrm{\#})`$ is the quotient homomorphism.
Set
$$K_P=\mathrm{\Omega }(K_0;\underset{jP}{}\{_{k_j,j}\}).$$
In the following we consider knots up to $`C_k`$-equivalence and by a symbol $`K_1`$ $`(`$resp. $`K_P)`$ we express a knot that is $`C_k`$-equivalent to $`K_1`$ $`(`$resp. $`K_P)`$. We will deform the form of band description of $`K_1`$ step by step. At each step $`K_1`$ is expressed as a band sum of $`K_0`$ and some chords such that each $`K_P`$ is a band sum of $`K_0`$ and some subset of the chords of $`K_1`$. To be more precise
$$K_1=\mathrm{\Omega }(K_0;\underset{i,j}{}\{_{i,j}\})$$
at each step where $`_{i,j}`$ is a $`C_i`$-chord for some $`i`$ with $`1i<k`$ and it has an associated subset $`\omega (_{i,j})\{1,\mathrm{},l\}`$ with $`_{t\omega (_{i,j})}k_ti`$ such that for each $`P\{1,\mathrm{},l\}`$
$$()K_P=\mathrm{\Omega }(K_0;\underset{\omega (_{i,j})P}{}\{_{i,j}\}).$$
A chord $`_{i,j}`$ is called a local chord if there is a 3-ball $`B`$ such that $`B`$ contains all of the bands and the link ball of $`_{i,j}`$, $`B`$ does not intersect any other bands and link balls, and that $`(B,BK_0)`$ is a trivial ball-arc pair. Such a local chord $`_{ij}`$ represents a knot $`K_{ij}`$ connected summed to $`K_0`$. The final goal of the step by step deformation is a band sum of $`K_0`$ and some local chords $`_{ij}`$’s so that $`K_1`$ is a connected sum of $`K_0`$ and $`K_{ij}`$’s. Since $`\mu _k(K\mathrm{\#}K^{})=\mu _k(K+K^{})𝒦/(𝒱(k)+_\mathrm{\#})`$, we have
$$\begin{array}{cc}& \mu _k\left(\underset{P\{1,\mathrm{},l\}}{}(1)^{|P|}\mathrm{\Omega }(K_0;\underset{\omega (_{i,j})P}{}\{_{i,j}\})\right)\hfill \\ \hfill =& \mu _k\left(\underset{P\{1,\mathrm{},l\}}{}(1)^{|P|}\left(K_0+\underset{\omega (_{i,j})P}{}K_{i,j}\right)\right)\hfill \\ \hfill =& \mu _k\left(\underset{P\{1,\mathrm{},l\}}{}(1)^{|P|}K_0+\underset{P\{1,\mathrm{},l\}}{}(1)^{|P|}\left(\underset{\omega (_{i,j})P}{}K_{i,j}\right)\right)\hfill \\ \hfill =& \mu _k\left(0+\underset{i,j}{}\left(\underset{P\{1,\mathrm{},l\},\omega (_{i,j})P}{}(1)^{|P|}K_{i,j}\right)\right).\hfill \end{array}$$
We consider the coefficient of $`K_{i,j}`$. Since $`_{t\omega (_{i,j})}k_t<k`$, $`\omega (_{i,j})`$ is a proper subset of $`\{1,\mathrm{},l\}`$. We may assume that $`\omega (_{i,j})`$ does not contain $`a\{1,\mathrm{},l\}`$. Then we have that
$$\begin{array}{cc}\hfill \underset{P\{1,\mathrm{},l\},\omega (_{i,j})P}{}(1)^{|P|}=& \underset{P\{1,\mathrm{},l\}\{a\},\omega (_{i,j})P}{}(1)^{|P|}\hfill \\ & +\underset{P\{1,\mathrm{},l\}\{a\},\omega (_{i,j})P}{}(1)^{|P\{a\}|}=0.\hfill \end{array}$$
Thus, we have the conclusion if we can get a desired band description.
Now we will deform the band description of $`K_1`$ into a desired form. We first set $`\omega (_{k_j,j})=\{j\}`$ for $`j=1,\mathrm{},l`$. Then we have $`_{t\omega (_{k_j,j})}k_t=k_j`$ and
$$K_P=\mathrm{\Omega }(K_0;\underset{\omega (_{k_j,j})P}{}\{_{k_j,j}\}).$$
Note that a crossing change between bands or a self-crossing change of a band can be realized by crossing changes between $`K_0`$ and a band as illustrated in Fig. 3.10. Therefore we can deform each chord into a local chord by band exchanges and crossing changes between $`K_0`$ and bands.
When we perform a crossing change between $`K_0`$ and a $`C_p`$-band of a $`C_p`$-chord $`_{p,q}`$ with $`pk2`$ we introduce a new $`C_{p+1}`$-chord $`_{p+1,r}`$ and we set $`\omega (_{p+1,r})=\omega (_{p,q})`$ so that the condition $`()`$ still holds for all subset $`P`$ of $`\{1,\mathrm{},l\}`$. To do this we use Lemma 3.8. When we perform a band exchange between a $`C_p`$-chord $`_{p,q}`$ and a $`C_r`$-chord $`_{r,s}`$ with $`p+rk1`$ we introduce a new $`C_{p+r}`$-chord $`_{p+r,n}`$ and set $`\omega (_{p+r,n})=\omega (_{p,q})\omega (_{r,s})`$ so that the condition $`()`$ still holds for all subset $`P`$ of $`\{1,\mathrm{},l\}`$. To do this we use Lemma 3.9. Note that the condition $`_{t\omega (_{i,j})}k_ti`$ still holds for all chords.
By Lemma 3.8, a crossing change between $`K_0`$ and a $`C_{k1}`$-band is realized by a $`C_k`$-move and therefore does not change the $`C_k`$-equivalence class. Similarly, by Lemma 3.9, a band exchange between a $`C_p`$-chord $`_{p,q}`$ and a $`C_r`$-chord $`_{r,s}`$ with $`p+rk`$ is realized by a $`C_{p+r}`$-move. Therefore, by Remark 2.4, it also does not change the $`C_k`$-equivalence class.
We note that the process definitely ends at last because a deformation of a $`C_p`$-chords does not produces $`C_q`$-chords for $`qp`$. $`\mathrm{}`$
Fig. 3.10
Proof of Theorem 1.4. It is clear that if $`K_1`$ and $`K_2`$ are $`C_k`$-equivalent and if $`J_1`$ and $`J_2`$ are $`C_k`$-equivalent then the connected sums $`K_1\mathrm{\#}J_1`$ and $`K_2\mathrm{\#}J_2`$ are $`C_k`$-equivalent. Thus the binary operation is well-defined. We note that the ambient isotopy classes of oriented knots forms a commutative monoid under connected sum of knots with unit element a trivial knot. Therefore it is sufficient to show the existence of the inverse element. Let $`K`$ be an oriented knot. First we note that $`K`$ itself is $`C_1`$-equivalent to a trivial knot. Suppose that there is an oriented knot $`J`$ such that $`K\mathrm{\#}J`$ is $`C_{k1}`$-equivalent to a trivial knot $`O`$. Then by Lemma 3.6 we have that $`O`$ is a band sum of $`K\mathrm{\#}J`$ and some $`C_{k1}`$-link models. We choose a 3-ball $`B`$ in $`S^3`$ such that $`BK\mathrm{\#}J`$ is an unknotted arc in $`B`$. By an ambient isotopy we deform the link balls into $`B`$ and slide the ends of bands into $`B`$. Then using Lemma 3.8 we deform $`O`$ up to $`C_k`$-equivalence so that the whole of the bands are contained in $`B`$. Then we have that the result is a connected sum of $`K\mathrm{\#}J`$ and some knot $`L`$. Namely $`K\mathrm{\#}J\mathrm{\#}L`$ is $`C_k`$-equivalent to $`O`$. Thus $`J\mathrm{\#}L`$ is the desired knot. $`\mathrm{}`$
Proof of Theorem 1.5. Let $`\xi _k:𝒦/(𝒱(k)+_\mathrm{\#})𝒦/C_k`$ be a homomorphism defined by $`\xi _k(K)=[K]_{C_k}`$ for $`K𝒦`$ where $`[K]_{C_k}`$ denote the $`C_k`$-equivalence class of $`K`$. It follows from Theorem 1.4 that both $`\eta _k`$ and $`\xi _k`$ are well-defined homomorphisms. Then it is clear that both $`\xi _k\eta _k`$ and $`\eta _k\xi _k`$ are identities. Therefore both $`\eta _k`$ and $`\xi _k`$ are isomorphisms. $`\mathrm{}`$
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# Soliton excitations as emitted clusters on nuclear surfaces
## I Introduction
There are different theoretical models which describe cluster formation and emission from nuclei, and most of them use nonlinear partial differential equations. A fundamental understanding of nonlinear terms in a nuclear model reveals new phenomena and shapes more complicated than linear theory suggests. In this paper we introduce a model for the cluster emission, based on soliton solutions of nonlinear equation.
Soliton structures have been analyzed within the frameworks of hydrodynamics, nonlinear optics, solid state and plasma physics. Experimental and theoretical results suggest that solitons are non-dispersive, localized waves executing uniform motion that can be described by three interrelated parameters: amplitude $`A`$, half-width $`L`$, and velocity $`V`$. Furthermore, these structures arise as analytical solutions of non-linear dynamical systems, like the Korteweg de Vries (KdV) or Nonlinear Schrödinger (NLS) equations. Soliton theory has been applied many times in nuclear physics, so far. For exemple, it provided very good descriptions of some localized stable surface excitations , of spectroscopic factors in cluster and $`\alpha `$ emission , and of quasimolecular spectra for $`\alpha `$ plus heavy nuclei collisions, . Also, the cubic and quintic NLS equations were used in three-dimensional models for cluster emission , providing results in very good agreement with TDHF simulations.
In order to look for the possibility of describing nuclear phenomena such as cluster emission by soliton formation on the nuclear surface, it is necessary to assign a microscopic structure to the parent heavy nucleus and the emitted soliton cluster. The microscopic substructure further allows one to add shell corrections to the usual macroscopic liquid drop energy and thus give a complete descripton of the system, from the initial nucleus with no soliton substructure to one with a soliton-like structure on its surface and on out to possible cluster emission.
A straightforward way to accomplish this is to calculate shell effects obtained from the single-particle levels of an asymmetric two-center shell model. One of the centers is placed in the middle of a small emitted sphere and the other is the center of the heavy fragment. This approach allows a microscopic description of the nuclear evolution from one to two independent quantum systems.
The procedure involves calculating the total potential energy as the sum of the macroscopic energy and shell corrections which is then minimized, which yields a barrier that increases as a function of the amplitude of the soliton. Calculations have been performed for two possible reactions:
<sup>248</sup>No $``$ <sup>208</sup>Pb +<sup>40</sup>Ca
<sup>224</sup>Th $``$ <sup>208</sup>Pb +<sup>16</sup>O
We select spherical daughters and emitted clusters, in order to fit the asymmetric two-center shell model we have constructed, which goes from one sphere (parent) to two necked spheres and then to separation.
In the present model, we describe cluster emission processes by using such soliton-like shapes on the nuclear surface of the heavy fragment. For a given cluster geometry, we calculate the corresponding soliton parameters ($`A`$, $`L`$, $`V`$) as functions of the separation parameter, that is along the static path of the cluster emission process.
Solitary waves have been shown recently to exist on liquid drops, bubbles, and shells . The non-linear hydrodynamic equations are related with the KdV and mKdV equation generating localized patterns ranging from small oscillations to nonlinear ones, including solitons. This model has a Hamiltonian structure, and such soliton-like excitations were observed experimentally when the shape oscillations of a droplet became nonlinear. It is therefore natural to extend that this model to other drop-like systems, from neutron stars to hyperdeformed nuclei and fission.
In the first section we define the deformation space we that we work with and the shapes that can be obtained therein. In the second section we give a short description of the macroscopic-microscopic model we used in the calculations. Emphazis has been given on the asymmetric two-center shell model we constructed in order to approach solitonic shapes. Results for the two reactions given above are presented in the third section.
## II Space of deformation
The problem of describing cluster emission (or nuclear fission) by a convenient parametrization of the shape is not new, and it is a decisive factor determining the amount of calculations. In the present model we use a new type of parametrization, described by soliton shapes. As well as the other parametrizations are in relation with the nuclear system and its quantization, the soliton parametrization is related to the KdV and mKdV dynamics of the soliton. Moreover, the soliton model takes profite from the high stability in time of such shapes. While in general one associates solitons with macroscopic pictures, it is nevertheless true that the soliton dynamics can be quantized and the KdV (or mKdV) equations can be related to quantum systems like those described by Schrödinger equations, , as we noted above.
Solitons on the surface are described by two asymmetrical spheres smoothly joined one to another through a neck region (Fig. 1). There are three independent geometrical parameters which form the space of deformation: the distance $`R`$ between the centers of the two fragments, the emitted small sphere radius $`R_2`$ and the neck sphere radius $`R_3`$. The neck region is obtained by rolling a sphere of radius $`R_3`$ around the symmetry axis. Such shapes are generated by the following equation written in cylindrical coordinates:
$$\rho (z)=\{\begin{array}{cc}\sqrt{R_1^2(zz_1)^2},& z_1R_1zz_{c1}\\ \rho _3\sqrt{R_3^2(zz_3)^2},& z_{c1}zz_{c2}\\ \sqrt{R_2^2(zz_2)^2},& z_{c2}zz_2+R_2\end{array}$$
(1)
where the quantities not shown in Fig. 1 are: $`z_3`$–the position of the center of the neck sphere on the symmetry axis, $`z_{c1}`$ and $`z_{c2}`$–the positions of the intersection planes of the two fragements spheres with the neck sphere, and $`z_1`$ and $`z_2`$–the positions of the two spherical fragment centers. The heavy fragment radius $`R_1`$ is obtained from total nuclear volume conservation. The soliton solution along the $`\theta `$ direction
$$r_{surface}(\varphi ,\theta ,t)=A(\varphi ,t)\left[sech\frac{\theta Vt}{L}\right]^2$$
(2)
is characterized by the amplitude $`A`$, or the relative amplitude $`a=A/R_1`$, the half-width $`L`$, and the angular velocity $`V`$. The soliton solutions have a special shape–kinematic dependence, $`VA`$ and $`L1/\sqrt{A}`$, that is, a higher soliton is narrower and travels faster . This relation can be used to experimentally distinguish solitons from other normal modes of excitations (for example by calculating the reciprocal moment of inertia) . The amplitude of the soliton is related to the two-center shell model by the relation with $`A=RR_1+R_2`$. The halfwidth of the soliton is approximated with $`2\rho (z_{c2})`$, or with the diameter of the circular surface within the separation plane between the emitted sphere and the rest of the shape
$$L=\frac{2R_2A}{(1+a)(R_1R_2)}.$$
(3)
We mention that the cubic and quartic NLS equations are related to the mKdV equation by a very simple exponential transform, and actually there is no essential difference between the NLS and the mKdV solitons . As an example, we notice the connection between the KdV equation and the nuclear potential in the Schrödinger equation , or the relation between NLS solitons and coherent states . Some applications of the KdV or mKdV-solitons are macroscopic, but the soliton solutions can be quantified by standard procedures .
## III Macroscopic energy
The deformation energy $`E_{def}`$ is calculated in a macroscopic-microscopic approach:
$$E_{def}=E_{macro}+\delta E_{shell}+\delta P.$$
(4)
The macroscopic part $`E_{macro}`$ includes the shape-dependent components of the charged liquid drop:
$$E_{macro}=E_C+E_{Y+E},$$
(5)
where $`E_C`$ is the Coulomb energy and $`E_{Y+E}`$ is the surface or nuclear energy calculated within the Yukawa-plus-exponential model .
The Coulomb energy general form is the double-volume integral:
$$E_C=\frac{1}{2}_V_V\frac{\rho _e(r_1)\rho _e(r_2)d^3r_1d^3r_2}{r_{12}},$$
(6)
which is split into four parts, two of them being equal to one another :
$$E_C=\frac{\rho _{1e}^2}{2}_{V_1}d^3r_1_{V_1}\frac{d^3r_2}{r_{12}}+\rho _{1e}\rho _{2e}_{V_1}d^3r_1_{V_2}\frac{d^3r_2}{r_{12}}+\frac{\rho _{2e}^2}{2}_{V_2}d^3r_1_{V_2}\frac{d^3r_2}{r_{12}}$$
(7)
where $`r_{12}=|r_1r_2|`$. The first and the last terms represent the self-energies of the two fragments, and the middle term is the Coulomb interaction between the two fragments.
In cylindrical coordinates the three terms are given by:
$`E_c`$ $`=`$ $`{\displaystyle \frac{\rho _e^2}{10}}{\displaystyle _z^{}^{z^{\prime \prime }}}dz{\displaystyle _z^{}^{z^{\prime \prime }}}dz_1{\displaystyle _0^{2\pi }}d\phi {\displaystyle _0^{2\pi }}d\phi _1(\rho ^2{\displaystyle \frac{z}{2}}{\displaystyle \frac{\rho ^2}{z}})[\rho _1^2`$ (10)
$`\rho \rho _1\mathrm{cos}(\phi \phi _1)+\rho {\displaystyle \frac{\rho _1}{\phi _1}}\mathrm{sin}(\phi \phi _1)+{\displaystyle \frac{(zz_1)}{2}}{\displaystyle \frac{\rho _1^2}{z_1}}][\rho ^2+`$
$`\rho _1^22\rho \rho _1\mathrm{cos}(\phi \phi _1)+(zz_1)^2]^{1/2},`$
where $`\rho =\rho (z,\phi )`$ is the nuclear surface equation, and $`z^{}`$ and $`z^{\prime \prime }`$ are the intersections of the surface with $`Oz`$ axis.
The general form of the Yukawa-plus-exponential energy is:
$$E_{Y+E}=\frac{a_2}{8\pi ^2r_0^2a^4}_V_V\left(\frac{r_{12}}{a}2\right)\frac{\mathrm{exp}(r_{12}/a)}{r_{12}/a}d^3r_1d^3r_2,$$
(11)
where $`r_{12}=|𝐫_\mathrm{𝟏}𝐫_\mathrm{𝟐}|`$, $`a`$=0.68 fm accounts for the finite range of nuclear forces, and $`a_2=a_s(1\kappa I^2)`$. $`\kappa `$ is the asymmetry energy constant, and the surface energy constant is $`a_s`$=21.13 MeV. In a similar way to the Coulomb part, one obtains three terms:
$`E_{Y+E}`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{2}{}}}{\displaystyle \frac{a_{2i}}{8\pi ^2r_0^2a^4}}{\displaystyle _{V_1}}d^3r_1{\displaystyle _{V_1}}\left({\displaystyle \frac{r_{12}}{a}}2\right){\displaystyle \frac{\mathrm{exp}(r_{12}/a)}{r_{12}/a}}d^3r_2`$ (13)
$`{\displaystyle \frac{2\sqrt{a_{21}a22}}{8\pi r_0^2a^4}}{\displaystyle _{V_1}}d^3r_1{\displaystyle _{V_2}}\left({\displaystyle \frac{r_{12}}{a}}2\right){\displaystyle \frac{\mathrm{exp}r_{12}/a}{r_{12}/a}}d^3r_2`$
where $`a_{2i}=a_s(1KI_i^2)`$, $`I_i=(N_iZ_i)/A_i`$. For shapes with axial symmetry, each of these terms involving a double-folded integration over the nuclear volume can be reduced to a three-dimensional integral.
## IV The asymmetric two-center shell model
The two-center shell model was developed by the Frankfurt school for symmetric splitting and for low asymmetry . Here we present the main steps of the two-center shell model we developed for large asymmetry starting from another symmetrical two-center model .
The general Hamiltonian describing the evolution of the level scheme of the two-center shell model is based on two oscillators which split from an initial common oscillator. The usual spin-orbit interaction and the $`𝐥^\mathrm{𝟐}`$ term are constructed as depending on the mass asymmetry. Thus, the total Hamiltonian reads:
$$H=H_{osc}+V(𝐥𝐬)+V(𝐥^\mathrm{𝟐})$$
(14)
where $`H_{osc}`$ is the two-oscillator Hamiltonian, and $`V(𝐥𝐬)`$ and $`V(𝐥^\mathrm{𝟐})`$ are the spin-orbit and the $`𝐥^\mathrm{𝟐}`$ potentials.
### A The diagonalization basis
The oscillator part of the Hamiltonian, $`H_{osc}`$, is given in cylindrical coordinates by:
$$H_{osc}=\frac{\mathrm{}^2}{2m_0}\left[\frac{^2}{\rho ^2}+\frac{1}{\rho }\frac{}{\rho }+\frac{1}{\rho ^2}\frac{^2}{\varphi ^2}+\frac{^2}{z^2}\right]+V(\rho ,z),$$
(15)
where the asymmetric two-center oscillator potential has the form:
$$V(\rho ,z)=\frac{1}{2}m_0\{\begin{array}{ccc}\omega _{\rho _1}^2\rho ^2+\omega _{z_1}^2(z+z_1)^2& ,\hfill & z<z_0\\ \omega _{\rho _2}^2\rho ^2+\omega _{z_2}^2(zz_2)^2& ,\hfill & zz_0.\end{array}$$
(16)
Here $`z_0`$ is the separation plane coordinate between the two asymmetric systems on the symmetry axis Oz. Since we consider the case of two asymmetric spheres, $`\omega _{\rho _1}=\omega _{z_1}=\omega _1`$ and $`\omega _{\rho _2}=\omega _{z_2}=\omega _2`$.
To obtain an appropiate basis for this system, we consider the intermediate case when $`\omega _{\rho _1}=\omega _{\rho _2}=\omega _1`$. Then the potential is:
$$V(\rho ,z)=\frac{1}{2}m_0\{\begin{array}{ccc}\omega _1^2\rho ^2+\omega _1^2(z+z_1)^2& ,\hfill & z<0\\ \omega _1^2\rho ^2+\omega _2^2(zz_2)^2& ,\hfill & z0.\end{array}$$
(17)
At this point the intermediate Hamiltonian is separable, and one gets the eigenfunctions :
$$\mathrm{\Phi }_m(\varphi )=\frac{1}{\sqrt{2\pi }}\mathrm{exp}(im\varphi )$$
(18)
for the axial degree of freedom and
$$R_{n_\rho |m|}(\rho )=\sqrt{\frac{2\mathrm{\Gamma }(n_\rho +1)\alpha _1^2}{\mathrm{\Gamma }(n_\rho +|m|+1)}}\mathrm{exp}\left(\frac{\alpha _1^2\rho ^2}{2}\right)(\alpha _1^2\rho ^2)^{\frac{|m|}{2}}L_{n_\rho }^{|m|}(\alpha _1^2\rho ^2)$$
(19)
for describing radial motion where $`\alpha _i=(m\omega _i/\mathrm{})`$, $`\mathrm{\Gamma }(x)`$ is the gamma function, and $`L_{n_\rho }^{|m|}(x)`$ is the Laguerre polynomial. As one knows, the oscillator energy for oscillations in the plane perpendicular on the symmetry axis is:
$$E_{\rho ,\varphi }=\mathrm{}\omega _\rho (2n_\rho +|m|+1)$$
(20)
where $`\omega _\rho =\omega _1`$.
Solving the third equation which accounts for oscillations along the symmetry axis, we have different solutions for the two regions of the nuclear shape. According to the $`z`$-dependent potential,
$$V(z)=\frac{1}{2}m_0\{\begin{array}{ccc}\omega _1^2(z+z_1)^2& ,\hfill & z<0\\ \omega _2^2(zz_2)^2& ,\hfill & z0\end{array}$$
(21)
where $`z_1`$ and $`z_2`$ are the centers of the heavy and light spherical fragments, respectively, one obtains the differential equations
$`\left[{\displaystyle \frac{d^2}{dz^2}}+{\displaystyle \frac{2m_0E_z}{\mathrm{}^2}}{\displaystyle \frac{m_0^2\omega _1^2(z+z_1)^2}{\mathrm{}^2}}\right]Z_{\nu _1}(z)=0`$ $`,z<0`$ (22)
$`\left[{\displaystyle \frac{d^2}{dz^2}}+{\displaystyle \frac{2m_0E_z}{\mathrm{}^2}}{\displaystyle \frac{m_0^2\omega _2^2(zz_2)^2}{\mathrm{}^2}}\right]Z_{\nu _2}(z)=0`$ $`\text{, z}\text{ 0 }.`$ (23)
At this point it is important to mention that the $`z=0`$ plane is the intersection plane between the two systems with $`\omega _{\rho _1}=\omega _{\rho _2}`$, whereas the “real” intersection between the asymmetric spherical systems is at $`z=z_0`$. The solution for the $`z`$-dependent Hamilton equation will be:
$$Z_\nu (z)=\{\begin{array}{ccc}C_{\nu _1}\mathrm{exp}\left[\frac{\alpha _1^2(z+z_1)^2}{2}\right]H_{\nu _1}[\alpha _1(z+z_1)]& ,\hfill & z<0\\ C_{\nu _2}\mathrm{exp}\left[\frac{\alpha _2^2(zz_2)^2}{2}\right]H_{\nu _2}[\alpha _2(zz_2)]& ,\hfill & z0\end{array}$$
(24)
where $`C_{\nu _1}`$ and $`C_{\nu _2}`$ are normalization constants, and $`H_\nu (z)`$ are the Hermite functions.
As can be seen from these results, four quantities need to be determined: the two quantum numbers $`\nu _1`$ and $`\nu _2`$ and the two normalization constants $`C_{\nu _1}`$ and $`C_{\nu _2}`$. These quantities can be calculated from a system of four equations. From the normalization condition
$$_{\mathrm{}}^{\mathrm{}}|Z_\nu (z)|^2𝑑z=1,$$
(25)
from the continuity of the $`z`$-wave function and its derivative at $`z`$=0
$$Z_{\nu _1}(z=0)=Z_{\nu _2}(z=0),$$
(26)
$$Z_{\nu _1}^{}(z=0)=Z_{\nu _2}^{}(z=0),$$
(27)
and from the energy matching condidtion along the O<sub>z</sub> axis
$$\mathrm{}\omega _1(\nu _1+0.5)=\mathrm{}\omega _2(\nu _2+0.5).$$
(28)
From these, a basis for diagonalization of the potential differences to obtain the real energy values can be calculated.
### B The asymmetric oscillator system
Once we have total wave functions, we have to determine differences between the diagonal Hamiltonian and the real one. First, the oscillator Hamiltonian has to provide the initial oscillator potential when there is only one heavy sphere (starting point). For this initial configuration the difference that needs to be diagonalized is
$$\mathrm{\Delta }V^{sphere}(z)=\{\begin{array}{ccc}\frac{1}{2}m_0[\omega _1^2(z+z_1)^2\omega _2^2(zz_2)^2]& ,\hfill & z0\\ 0& ,\hfill & z<0.\end{array}$$
(29)
For the next stages of deformation, the difference between the $`z`$-dependent oscillator potentials that needs to be diagonalized is:
$$\mathrm{\Delta }V(z)=\{\begin{array}{ccc}0& ,\hfill & z<0\\ \frac{1}{2}m_0[\omega _1^2(z+z_1)^2\omega _2^2(zz_2)^2& ,\hfill & 0zz_0\\ 0& ,\hfill & z>z_0.\end{array}$$
(30)
As for the difference in the $`\rho `$-dependent oscillator potential, this only exists for intersecting spheres and is given by:
$$\mathrm{\Delta }V(\rho )=\{\begin{array}{ccc}0& ,\hfill & zz_0\\ \frac{1}{2}m_0(\omega _1^2\omega _2^2)\rho ^2& ,\hfill & z>z_0\end{array}$$
(31)
or, if written as an operator, the quantity to be diagonalized is given by:
$$\mathrm{\Delta }V(\rho )=\frac{1}{2}m_0(\omega _1^2\omega _2^2)\rho ^2\mathrm{\Theta }(zz_0)$$
(32)
where $`\mathrm{\Theta }(z)`$ is the Heaviside function. The difference $`\mathrm{\Delta }V(\rho )`$ is zero for the initial spherical configuration.
Once $`\mathrm{\Delta }V(z)`$ and $`\mathrm{\Delta }V(\rho )`$ are diagonalized and added to the oscillator energy of the sphere + ellipsoid system, which is
$$E=\mathrm{}\omega _1[2n_\rho +|m|+\nu _1+1.5],$$
(33)
the level schemes of the two intersecting asymmetric oscillators with frequencies $`\omega _1`$ and $`\omega _2`$ are obtained.
### C Spin-orbit and orbit-orbit interactions
The spin-orbit ($`𝐥𝐬`$) and orbit-orbit ($`𝐥^\mathrm{𝟐}`$) interaction terms generate the necessary single-particle level splitting to obtain the correct schemes of the individual fragments after separation.
The use of deformation dependent form of these operators has been introduced in for the Nilsson model and in for the two center shell model; instead of the $`𝐥`$ operator one introduces:
$$𝐥=\frac{V\times 𝐩}{m_0\omega ^2}$$
(34)
where $`V`$ is the asymmetric two-center oscillator potential. The usual expression for the two operators are:
$`V(𝐥𝐬)`$ $`=`$ $`2\kappa \mathrm{}\omega 𝐥𝐬`$ (35)
$`V(𝐥^\mathrm{𝟐})`$ $`=`$ $`\kappa \mu \mathrm{}\omega 𝐥^\mathrm{𝟐}.`$ (36)
Since one obtains the level schemes of two nuclei which lie in different mass regions, the strength parameters of the interactions $`\kappa `$ and $`\mu `$ will be different. The values we use for these parameters are:
$$\begin{array}{ccc}\kappa _n=0.0588& & \kappa _p=0.0592\\ \mu _n=0.328& & \mu _p=0.335\end{array}$$
(37)
for actinide region, and for light nuclei region:
$$\begin{array}{ccc}\kappa _n& =& \kappa _p=0.0601\\ \mu _n& =& \mu _p=0.448.\end{array}$$
(38)
Since the strength parameters are different for the asymmetric regions of the nuclear shape, they become $`z`$-dependent operators as follows:
$$\begin{array}{ccc}\kappa \mathrm{}\omega (z)=\kappa _1\mathrm{}\omega _1+(\kappa _2\mathrm{}\omega _2\kappa _1\mathrm{}\omega _1)\mathrm{\Theta }(zz_0)& & \\ \kappa \mu \mathrm{}\omega (z)=\kappa _1\mu _1\mathrm{}\omega _1+(\kappa _2\mu _2\mathrm{}\omega _2\kappa _1\mu _1\mathrm{}\omega _1)\mathrm{\Theta }(zz_0).& & \end{array}$$
(39)
To obtain a Hermitian operator for $`V(𝐥𝐬)`$ and $`V(𝐥^\mathrm{𝟐})`$ one has to use the anticommutator :
$$\begin{array}{ccc}V(𝐥𝐬)=[\kappa \mathrm{}\omega (z),\frac{V\times 𝐩}{m_0\omega ^2}𝐬],& & \\ V(𝐥^\mathrm{𝟐})=\frac{1}{2}[\kappa \mu \mathrm{}\omega (z),\left(\frac{V\times 𝐩}{m_0\omega ^2}\right)^2].& & \end{array}$$
(40)
For the dependence of $`\kappa `$ and $`\mu `$ with respect to the elongation, we choose a linear dependence for the oscillator frequency along the $`z`$-axis:
$$\kappa _i=\kappa _0+\frac{\omega _i\omega _0}{\omega _{if}\omega _0}(\kappa _{if}\kappa _0)$$
(41)
and the same law of variation for $`\mu `$. Here $`i`$=1,2, $`\kappa _0`$ is the value for the initial nucleus and $`\kappa _{if}`$ for the final one.
Finally one has to diagonalize the potential:
$$\mathrm{\Delta }V(\rho ,z)=\mathrm{\Delta }V(z)+\mathrm{\Delta }V(\rho )+V(𝐥𝐬)+V(𝐥^\mathrm{𝟐})$$
(42)
together with the diagonal term of the two-center oscillator potential. The model provides the evolution from an initial level scheme toward two asymptotically independent single-particle schemes. With the introduction of large asymmetry between fragments, the shapes can simulate the existence of a soliton on the nuclear surface and assign it a microscopic structure.
## V Shell corrections
The level scheme of a soliton shape is used to obtain the shell corrections of the system. As the soliton is assimilated with an emerging fragment, it will provide the shell correction value of the independent nucleus of similar shape. Shell corrections are obtained by means of the Strutinsky procedure . One defines the shell correction energy as the difference between the total sum of the energy levels and a smoothed part of the spectrum:
$$\delta E=\underset{\nu }{}2E_\nu \stackrel{~}{U}.$$
(43)
One calculates the smoothed part $`\stackrel{~}{U}`$ with the help of a smoothed level density function $`\stackrel{~}{g}(ϵ)`$, which is obtained by averaging the real distribution $`g(ϵ)`$ over the whole energy spectrum:
$`\stackrel{~}{g}(ϵ)`$ $`=`$ $`{\displaystyle \frac{1}{\gamma }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}\zeta \left({\displaystyle \frac{ϵϵ^{}}{\gamma }}\right)g(ϵ^{})𝑑ϵ^{}`$ (44)
$`=`$ $`{\displaystyle \frac{1}{\gamma }}{\displaystyle \underset{i=1}{\overset{\mathrm{}}{}}}\zeta \left({\displaystyle \frac{ϵϵ^{}}{\gamma }}\right),`$ (45)
where $`\gamma =\mathrm{\Gamma }/\mathrm{}\omega `$ and $`\zeta (x)`$ is the smoothing function. A common smoothing function is provided by
$$\zeta (x)=\frac{1}{\sqrt{\pi }}e^{x^2}\underset{k=0}{\overset{m}{}}a_{2k}H_{2k}(x),$$
(46)
where $`H_n(x)`$ are the Hermite polynomials. The coefficients $`a_{2k}`$ are
$$a_{2k}=\frac{H_{2k}(0)}{2^{2k}(2k)!}.$$
(47)
The Fermi energy $`\stackrel{~}{\lambda }`$ of the smoothed level distribution is calculated as a solution of the particle number conservation:
$$N_p=2_{\mathrm{}}^{\stackrel{~}{\lambda }}\stackrel{~}{g}(ϵ)𝑑ϵ.$$
(48)
Then, the total energy of the uniform level distribution $`\stackrel{~}{U}`$, reproducing the microscopic part which is not subjected to local fluctuations of the spectrum, is obtained as:
$$\stackrel{~}{U}=2\mathrm{}\omega _{\mathrm{}}^{\stackrel{~}{\lambda }}\stackrel{~}{g}(ϵ)ϵ𝑑ϵ.$$
(49)
After performing the calculations, one obtains the following formula, which can be used directly
$`\delta U`$ $`=`$ $`{\displaystyle \underset{\nu }{}}\{ϵ_\nu [1erf(x_{F_\nu })]`$ (51)
$`+{\displaystyle \frac{e^{x_{F_\nu }^2}}{\sqrt{\pi }}}[2ϵ_\nu {\displaystyle \underset{k=1}{\overset{m}{}}}a_{2k}H_{2k1}+\stackrel{~}{\gamma }a_{2m}H_{2m}]\},`$
where $`erf(x)`$ is the error function. Usually one chooses the upper order of the Hermite polynomials to be $`m`$=3. The variable $`x_{F_\nu }`$ is given by
$$x_{F_\nu }=\frac{ϵ_\nu ϵ_F}{\gamma },$$
(52)
Shell corrections are calculated separately for protons and neutrons, and the results are added.
## VI Results
A first look at the energetic behavior of an emerging soliton is given in Fig. 2 in terms of the macroscopic energy surfaces. The LHS macroscopic energy surface corresponds to the formation of <sup>40</sup>Ca on the surface of <sup>248</sup>No, whereas the RHS represents the <sup>16</sup>O-like soliton on the surface of <sup>224</sup>Th. Variation along the elongation $`R`$ corresponds to the increment in the soliton amplitude along the symmetry axis. A larger neck radius $`R_3`$ corresponds to a larger half-width $`L`$. The rear plane at $`R`$=0 is the spherical state of the system. Then the energy increases monotonously with a higher slope for small values of the neck radius. As $`R_3`$ increases, the energy increase is smoothed by the necking. With the enhancement of the kinetic energy of the soliton the half-width becomes larger, except in the first stages of the process where the neck radius is very small. The ridge in energy has a maximum at the near touching spheres configuration for both reactions. The slope continues to increase for large $`R_3`$, beyond the touching point value of the elongation $`R`$.
The addition of shell corrections yields the total deformation energy shown in Fig. 3. As a first observation note the pronounced deformed ground state of <sup>248</sup>No as the first minimum in energy moves to $`R>`$0, and a much less but still deformed ground state for <sup>224</sup>Th. For both emerging solitons it is obvious that the energy path corresponds to large half-width values up to the top of the energy ridge; then they abruptly turn towards rupture point shapes ($`R_3`$=0). Hence, these potential energy surfaces suggest a three-dimensional curve as the path of minimum energy in the cluster emission. The potential barrier formed along the path of minimum energy values is obtained by minimization of the total energy in the multi-dimensional deformation space.
The static paths, which a soliton with the internal structure of an emitted cluster has to follow, have been plotted on the contour maps of the energy surface in ($`R`$,$`R_3`$) coordinate space in Fig. 4. Again the LHS plot is the <sup>40</sup>Ca emission, and the RHS one corresponds to <sup>16</sup>O. Apart from the first $`R`$ values, where $`R_3`$ is small, the solitons bypass the first energy peak by taking large neck radius values, i.e. large half-widths. As the energy increases on the large $`R_3`$ side, the static path for both cases changes direction reaching the scission point where the clusters are emitted.
The two barriers are plotted in Fig. 5 with a full line, together with the macroscopic energy (dotted line) and the shell corrections (dashed line). One can see how the deformed ground state of <sup>248</sup>No is formed (LHS plot): due to shell corrections, the first minimum is at about 6.8 MeV of the total energy. This point becomes the ground state and the whole barrier in front of the emerging soliton is shifted with respect to this value. A two-humped barrier no higher than about 1.2 MeV blocks the <sup>40</sup>Ca emission.
The situation is different for the emission of <sup>16</sup>O from <sup>224</sup>Th. The ground state is only slightly deformed. A rather high one-hump barrier of about 11 MeV extends along the whole range of elongation $`R`$ up to the scission point. Shell corrections decrease slightly the macroscopic energy values. The decrease is mainly due to the double-magic character of <sup>208</sup>Pb which forms as the cluster emerges.
One can state that <sup>40</sup>Ca-like solitons are energetically favored to form on the nuclear surface of a very heavy nucleus as <sup>248</sup>No. The formation of <sup>16</sup>O-like soliton on <sup>224</sup>Th is not energetically favored due to the high and large potential barrier it has to penetrate.
The relative velocity distribution $`V`$ of the two presumed solitons along the minimum energy path, together with the scaled values of the halfwidth $`L`$ and the relative amplitude $`a=A/R_1`$, are plotted in Figs. 6. We investigated the evolution of these soliton parameters (as defined in section I), which are a function of the static energy evolution, parametrized by the distance between centers $`R`$.
In the first stages, the tendency is that the amplitude and half-width increase with the elongation parameter, when the emitted cluster is emerging out from the parent nucleus (since their non-overlaping sector is increasing). During the formation of the cluster the half-width remains practically constant, since the surface energy controls this stage. When the two nuclei are well separated, the soliton envelope hardly fits the two spheres, and in this limit, the half-width approaches zero value. This gives the limiting configuration for this soliton model. These values of the half-width $`L`$ (solid line) are compared with those obtained analyticaly in directly from the soliton amplitude, within the frame of the nonlinear liquid drop model ($`L`$-dashed line). We notice a good agreement for the half-widths within the range $`R4.514`$ Fm. The hydrodynamic soliton model is not valid anymore for separation parameter $`R`$ smaller than 4-5 Fm, because of the dominating shell effects in this range. This can be noticed in a comparison between Fig. 2 and Fig. 3 for $`R45`$ Fm. For the first $`R`$ values, the static paths follow the first energy peak, and jump from small toward large values for $`R_3`$, Fig. 4, providing small half-widths (Figs. 6, $`L`$-solid line), while a pure hydrodynamic soliton would have larger half-widths for this range ($`L`$-dashed line). In the above range of validity of the soliton model, we calculate the relative velocity of the soliton ($`V`$-dashed line), . The velocity is increasing with the amplitude of the soliton, hence with the elongation of the cluster-like emission shape. Fig. 6a displays the <sup>40</sup>Ca emission, and Fig. 6b represents the <sup>16</sup>O emission. Soliton shapes at the begining and the end of the process are also shown. For lighter nuclei (like <sup>16</sup>O) the evolution of the parameters is smooth and monotonic. In the case of heavier nuclei (<sup>40</sup>Ca) we obtained some oscillations in width and velocity, during the first half of the emission process, which can be related with the oscillations produced by the shell effects in the $`R_3`$ parameter.
This work was supported by the U.S. National Science Foundation through a regular grant, No. 9970769, and a Cooperative Agreement, No. EPS-9720652, that includes matching from the Louisiana Board of Regents Support Fund. RAG is grateful to Prof. J. P. Draayer for a postdoctoral appointment in the Department of Physics and Astronomy, Louisiana State University, Baton Rouge.
Figure captions
Fig.1 Deformation space for two-center shell-model calculations. The neck radius $`R_3`$ can vary from zero (intersecting spheres) to infinity (compact shapes). $`R`$ (the distance between the two centers) and $`R_2`$ (radius of the emitted fragment) are also independent coordinates.
Fig.2 Macroscopic potential energy surfaces for <sup>40</sup>Ca emission from <sup>248</sup>No (left-hand-side plot) and for <sup>16</sup>O emission from <sup>224</sup>Th (right-hand-side plot), as function of elongation $`R`$ and neck radius $`R_3`$. First maximum appears close to the touching point configuration in both cases ($`R_3`$=0).
Fig.3 Total potential energy surfaces (macroscopic plus shell corrections) for <sup>40</sup>Ca emission from <sup>248</sup>No (left-hand-side plot) and <sup>16</sup>O from <sup>224</sup>Th (right-hand-side plot). The deformed ground state of <sup>248</sup>No is revealed as a minimum along $`R_3`$ axis at its origin. For both surfaces the closest energy maximum occurs at the tangent sphere configuration. The maximum energy value for larger $`R_3`$ is not reached in the figure.
Fig. 4 Contour plot of Fig. 3 with static path (dashed line) for <sup>40</sup>Ca emission (left-hand-side plot) and <sup>16</sup>O emission (right-hand-side plot). As the elongation increases the amplitude of the soliton increases together with the half-width which is proportional to the neck radius $`R_3`$. Once the touching point maximum is bypassed, both shapes decrease rather abruptly through necking towards scission.
Fig. 5 The barriers along the static path for <sup>40</sup>Ca emission (left-hand-side plot) and <sup>16</sup> (right-hand-side plot), together with macroscopic energy (dotted lines) and shell corrections (dashed lines). Shell corrections increase the total energy of the first energy minimum for <sup>248</sup>No, thus the two-humped barrier is not higher than about 1.2 MeV. <sup>16</sup>O emission from <sup>224</sup>Th has a barrier of about 11 MeV.
Figs. 6 The evolution of $`a=A/R_1`$, $`L`$ (with shell corrections solid line, and without shell corrections dashed line) and $`V`$ parameters in relative units, versus the elongation $`R`$ in fm. The corresponding nuclear configurations are plotted for two situations: for the initial stage when the emitted cluster is only slightly displaced off the common center, and the final stage when the two nuclei are almost separated. Fig. 6a displays the <sup>40</sup>Ca emission and Fig. 6b the <sup>16</sup>O emission. In the later case, oscillations can be seen in the soliton parameters related to the shell corrections.
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# Kondo effect in multielectron quantum dots at high magnetic fields.
## I Introduction
Kondo effect has been recently observed in quantum dots (QD) coupled to leads by analyzing the temperature dependence of the conductance . Kondo physics has been mainly detected for an odd number of electrons $`N`$ in the QD while for even $`N`$, the temperature dependence of the conductance is usually that corresponding to a normal system. This is interpreted in terms of the formation of singlets by pairs of electrons so that only when $`N`$ is odd a single electron remains unpaired. This last electron, in a state doubly degenerate by spin, is responsible for the Kondo behavior proposed long time ago . However, some recent experiments suggests a more complicated situation because Kondo effect is also observed for even $`N`$ in some situations. A particularly interesting experimental feature is the appearance, for a given $`N`$, of alternating high and low conductance valleys as a function of an external magnetic field $`B`$. When $`N`$ is varied in $`\pm 1`$, the high and low conductance valleys are interchanged. Therefore, the representation of the conductance (in a grey scale) as a function of both $`B`$ and a gate potential which allows to vary $`N`$, takes the aspect of a chess board.
Three models have been already proposed for understanding how Kondo effect might be possible for states of two electrons. In these approaches, the authors consider the situation of double degeneracy that a magnetic field can create between the singlet ($`S_z=0`$) and the triplet ($`S_z=1`$) states of the electronic pair. However, as shown below, the actual situation is, in general, rather more complicated due to correlation and tunneling amplitude effects.
In this paper we present a general description of the Kondo effect in a QD in the presence of a high magnetic field such that the filling factor is $`1<\nu <2`$. This is a regime in which experiments clearly show Kondo effect. The description is valid for any number of electrons (even or odd without restriction to 1 or 2) and any value of the QD spin $`S_z`$ (not restricted to 0, 1 for even $`N`$ or 1/2 for odd $`N`$). Instead of describing spin-flip scattering in terms of spin-ladder operators $`S^{(\pm )}`$, we find a set of spin-flip Hubbard operators describing a collective spin effect of all the $`N`$ electrons contained in the QD. Despite both $`N`$ and $`S_z`$ can be very large, the scattering of carriers only produces transitions between two many-body states of the QD with $`S_z`$ differing in $`\pm 1`$. These many-body states are well characterized theoretically in the regime $`1<\nu <2`$, and we find that the spin-flip process can be described by a Kondo Hamiltonian with antiferromagnetic couplings which depend on both tunneling amplitudes and correlation effects. As a consequence, the system presents a Kondo behavior with a Kondo temperature $`T_K`$ which we analyze in various limiting cases.
In section II we discuss the ground state (GS) of an isolated QD in the regime $`1<\nu <2`$. in Section III an effective Kondo Hamiltonian is obtained, by means of a scattering description of tunneling through the QD. Section IV contains a discussion on the exchange couplings and the Kondo temperature in the regime $`1<\nu <2`$. Section V is devoted to the explanation of the chess board aspect of the experimental conductance. A summary is given in section VI.
## II QD spectrum
We consider $`N`$ electrons in the presence of a magnetic field and confined in a QD coupled to leads. The Hamiltonian is
$`H=H_{QD}+H_L+H_{TUN}.`$ (1)
Within a lowest Landau level approach, an isolated parabolic QD with magnetic field and interaction between the electrons is described by (hereafter we take $`\mathrm{}=1`$)
$`H_{QD}`$ $`={\displaystyle \frac{N\mathrm{\Omega }}{2}}+{\displaystyle \frac{(\mathrm{\Omega }\omega _c)M}{2}}+g\mu _BBS_z+`$ (3)
$`{\displaystyle \frac{1}{2}}{\displaystyle \underset{m_i,\sigma _i}{}}V_{m_1m_2m_3m_4}d_{m_1,\sigma _1}^{}d_{m_2,\sigma _2}^{}d_{m_3,\sigma _3}d_{m_4,\sigma _4}.`$
The first two terms describe a single particle contribution depending on both the QD confinement $`\omega _0`$ and cyclotron $`\omega _c`$ frequencies through $`\mathrm{\Omega }=\sqrt{\omega _c^2+4\omega _0^2}`$. $`M`$ is the third component of the total angular momentum. The third term is the Zeeman energy, depending on the Landé $`g`$-factor. In the last term, describing electron-electron repulsion, $`d_{m,\sigma }^{}`$ creates an electron with angular momentum $`m`$ and spin $`\sigma `$ in the QD. The Coulomb interaction matrix elements $`V_{m_1m_2m_3m_4}`$ have a typical energy scale $`e^2/\epsilon l_B`$, where $`\epsilon `$ is the dielectric constant and $`l_B=1/\sqrt{m\mathrm{\Omega }}`$ the magnetic length. These interactions decrease with the increasing width $`w`$ of the quantum well in which the QD has been built up.
$`H_L=_{k,\sigma }\epsilon _{k,\sigma }c_{k,\sigma }^{}c_{k,\sigma }`$ describes, in a single-particle approach, the leads having electrons with quantum numbers $`k,\sigma `$ occupying states up to the Fermi energy $`\epsilon _F`$.
A crucial role is going to be played by the tunneling part of the Hamiltonian
$`H_{TUN}={\displaystyle \underset{k,m,\sigma }{}}V_m\left(d_{m,\sigma }^{}c_{k,\sigma }+c_{k,\sigma }^{}d_{m,\sigma }\right)`$ (4)
in which we neglect any dependence on $`k`$ of the tunneling amplitudes $`V_m`$ (taken as real positive) but we retain the dependence on $`m`$ because is going to produce physical consequences.
$`H_{QD}`$ can be numerically diagonalized for a significantly broad range of $`N`$, as discussed in many theoretical papers. An example (similar results are obtained for an odd number of electrons) is shown in Fig. 1 and 2, for $`N=8`$, which is the regime in which some experiments have been performed, although some others involve a larger number of electrons. In order to have access to all the possibilities for the GS, it is convenient to discuss a case with strong interaction effects, so that we present results for a QD with $`w=0`$. The case of an arbitrary finite width $`w`$ can be qualitatively discussed from these results: correlated GS tend to disappear when $`w`$ increases. Fig. 1 shows the phase diagram of the possible GS between $`\nu =1`$ and $`\nu =2`$. The wider regions of the upper part of the phase diagram correspond to compact states which are favored by large Zeeman coupling
$`|C_{NK}^K={\displaystyle \underset{m=0}{\overset{K1}{}}}d_{m,}^{}{\displaystyle \underset{m=0}{\overset{NK1}{}}}d_{m,}^{}|0`$ (5)
where $`|0`$ is the vacuum state. In going from left to right, one finds successively the states $`|C_4^4`$ ($`\nu =2`$), $`|C_5^3`$, $`|C_6^2`$, $`|C_7^1`$ and $`|C_8^0`$ ($`\nu =1`$). Strong interaction effects are manifested in the appearance, lower in the phase diagram, of the narrower regions corresponding to different skyrmion-like states
$`|SK_{N,K,\pm }^P=𝒩_{P,\pm }\left(\mathrm{\Lambda }_{\pm 1}^{}\right)^P|C_{NK}^K`$ (6)
of topological charge 1. $`𝒩_{P,\pm }`$ is a normalization constant and
$`\mathrm{\Lambda }_{\pm 1}^{}={\displaystyle \underset{m}{}}(m+1)^{1/2}d_{m+1,}^{}d_{m,}.`$ (7)
Adjacent regions correspond to values of $`P`$ which differ in one unit and, as a consequence, their spins also differ in 1. Dashed lines in Fig. 1 depict, for two different values of $`\omega _0`$, the evolution of the GS for a GaAs QD within the range of magnetic fields (perpendicular to the QD) given at the edges of the lines.
Fig. 2 shows the evolution of the GS properties from $`\nu =2`$ to $`\nu =1`$ for the system of Fig. 1 with $`\omega _0=5.4`$meV. The GS energy $`E_{GS}`$ has a kink any time a crossing of states occurs. Since these kinks are not obvious in Fig.2(b), we represent, in part (a), the energy splitting $`\mathrm{\Delta }E=E_{exc}E_{GS}`$ between the lowest excited state and the GS. The spin $`S_z`$ and the third component $`M`$ of the total angular momentum of the GS are also given in (c) and (d) respectively.
Independently of the particular details, there is a general trend for any $`N`$: when the magnetic field increases from $`\nu =2`$ to $`\nu =1`$, the GS changes many times (at least $`N/2`$ changes corresponding to electrons $`N/2`$ flipping their spins one by one). In other words, for $`2>\nu >1`$ there are several crossings of the lowest energy states. This means that, for many values of the field, the GS is doubly degenerated (or non degenerated but having an excited state with extremely low excitation energy). Moreover, in almost all the cases, these two states have spins differing in 1 while they have different $`M`$.
## III Kondo Hamiltonian
As a consequence of the above numerical result, the main physics of the problem is captured by a two level system approach . One considers the QD in the regime $`2>\nu >1`$ as having two degenerate (or almost degenerate) GS’s $`|GS_{}^N`$ and $`|GS_{}^N`$ with $`M_{}M_{}`$ and spins differing in 1, i.e. $`GS_{}^N|S_z|GS_{}^N=GS_{}^N|S_z|GS_{}^N+1`$. $`H_{TUN}`$ mixes these GS’s with states $`|N\pm 1`$ in which $`N\pm 1`$ electrons are in the QD. In our two level system description, the tunneling Hamiltonian is projected on the subspaces subtended by the two GS’s for $`N`$ electrons and the connecting $`|N\pm 1`$ electron states. In this process, we consider the connecting $`|N\pm 1`$ states as non-degenerate. Using the notation $`\mathrm{\Sigma }\{,\}`$, we introduce the tunneling spectral amplitudes
$`\mathrm{\Delta }_{,\mathrm{\Sigma }}={\displaystyle \underset{m,\sigma }{}}V_mN1|d_{m,\sigma }|GS_\mathrm{\Sigma }^N`$ (8)
$`\mathrm{\Delta }_{+,\mathrm{\Sigma }}={\displaystyle \underset{m,\sigma }{}}V_mN+1|d_{m,\sigma }^{}|GS_\mathrm{\Sigma }^N`$ (9)
and the spin-flip Hubbard operators
$`X_{\mathrm{\Sigma },\mathrm{\Sigma }^{}}=|GS_\mathrm{\Sigma }^NGS_\mathrm{\Sigma }^{}^N|.`$ (10)
From projected tunneling Hamiltonian $`\overline{H}_{TUN}`$, the interaction between the QD with $`N`$ electrons and the leads is studied by a standard scattering description . An effective coupling is obtained by summing up to all the possible intermediate states $`|I`$, in a second order perturbation approach:
$`H_{eff}={\displaystyle \underset{I}{}}{\displaystyle \frac{\overline{H}_{TUN}|II|\overline{H}_{TUN}}{E_{GS}^NE_I}}`$ (11)
$`={\displaystyle \underset{k,k^{},\mathrm{\Sigma }}{}}\left[{\displaystyle \frac{\mathrm{\Delta }_{,\mathrm{\Sigma }}^2\delta _{k,k^{}}X_{\mathrm{\Sigma },\mathrm{\Sigma }}}{E_{GS}^NE^{N1}\epsilon _F}}+{\displaystyle \frac{\mathrm{\Delta }_{+,\mathrm{\Sigma }}^2c_{k^{},\sigma }^{}c_{k,\sigma }}{E_{GS}^NE^{N+1}+\epsilon _F}}+{\displaystyle \underset{\mathrm{\Sigma }^{}}{}}({\displaystyle \frac{\mathrm{\Delta }_{+,\mathrm{\Sigma }}\mathrm{\Delta }_{+,\mathrm{\Sigma }^{}}}{E^{N+1}\epsilon _FE_{GS}^N}}+{\displaystyle \frac{\mathrm{\Delta }_{,\mathrm{\Sigma }}\mathrm{\Delta }_{,\mathrm{\Sigma }^{}}}{E^{N1}+\epsilon _FE_{GS}^N}})X_{\mathrm{\Sigma },\mathrm{\Sigma }^{}}c_{k^{},\sigma ^{}}^{}c_{k,\sigma }\right],`$ (12)
where, due to spin conservation, $`\mathrm{\Sigma }`$ has the same direction than $`\sigma `$ and $`\mathrm{\Sigma }^{}`$ the same than $`\sigma ^{}`$. The first term is simply a constant. The second term (involving $`c_{k^{},\sigma }^{}c_{k,\sigma }`$) represents a potential scattering which does not involve any spin flip. These two terms are identical to the ones appearing when building up an $`sd`$ Hamiltonian from the Anderson Hamiltonian in the case of $`N=1`$. As it is usually done in that case, one can forget about these two terms which do not contain anything important for the physics we want to address.
The interesting physics is included in the third term of (12) which is a Kondo Hamiltonian
$`H_K={\displaystyle \underset{k,k^{}}{}}`$ $`\left[J\left(X_,c_k^{}^{}c_{k,}+X_,c_{k^{},}^{}c_{k,}\right)+J_{}X_,c_{k^{},}^{}c_{k,}+J_{}X_,c_{k^{},}^{}c_{k,}\right].`$ (13)
with exchange couplings
$`J={\displaystyle \frac{\mathrm{\Delta }_{+,}\mathrm{\Delta }_{+,}}{E^{N+1}\epsilon _FE_{GS}^N}}+{\displaystyle \frac{\mathrm{\Delta }_,\mathrm{\Delta }_,}{E^{N1}+\epsilon _FE_{GS}^N}};J_\mathrm{\Sigma }={\displaystyle \frac{\mathrm{\Delta }_{+,\mathrm{\Sigma }}^2}{E^{N+1}\epsilon _FE_{GS}^N}}+{\displaystyle \frac{\mathrm{\Delta }_{,\mathrm{\Sigma }}^2}{E^{N1}+\epsilon _FE_{GS}^N}}.`$ (14)
$`H_K`$ is a spin-flip scattering Hamiltonian in which the GS of the scatterer flips its spin by means of $`X_,`$ and $`X_,`$ (as shown in Fig. 3) and, at the same time, changes its total angular momentum. Both the difference between $`M_{}`$ and $`M_{}`$, and the correlation effects included in the tunneling spectral amplitudes, are the reason why one must use spin-flip Hubbard operators instead of the usual spin-ladder operators $`S^{(\pm )}`$.
A crucial question is the sign of the exchange couplings (14). With our definitions (9), all the tunneling spectral amplitudes are positive. Therefore, the signs of the exchange couplings are determined by the denominators. As, in the considered situation, the lowest energy corresponds to having $`N`$ electrons in the QD, all the intermediate states $`|I`$, with $`N\pm 1`$ electrons in the QD and $`1`$ electron at the Fermi level of the leads, have higher energy. Therefore, we have a very important result: $`H_K`$ has positive effective exchange couplings. So, a Kondo Hamiltonian (13) with antiferromagnetic couplings (14) has been obtained. The conclusion is: for any values of both $`N`$ and spin, the QD in the regime $`2>\nu >1`$ presents Kondo physics.
## IV Exchange couplings and Kondo temperature
Since the Hamiltonian (13) is formally equivalent to the standard Kondo Hamiltonian for $`N=1`$, the temperature dependence of the conductance (as any other measurable property) presents now characteristics similar to the one found in those experiments which have been interpreted in terms of only one electron. The important issue to be discussed is the characteristic energy scale, $`T_K`$, which is determined by the antiferromagnetic couplings.
$`J`$ and $`J_\mathrm{\Sigma }`$ depend on both the energy difference $`E_{GS}^NE^{N\pm 1}\epsilon _F`$, and tunneling spectral amplitudes $`\mathrm{\Delta }_{,\mathrm{\Sigma }}`$. The former is practically independent on both $`N`$ and $`\mathrm{\Sigma }`$. Therefore, it does not imply any significant difference with respect to the well known case of $`N=1`$. However interesting physical differences appear due to $`\mathrm{\Delta }_{,\mathrm{\Sigma }}`$. Let us analyze the different situations in which the QD can be found:
i) (Upper part of Fig. 1) When the Zeeman interaction is large (for instance having a large inplane component of the magnetic field), the GS’s are always compact states
$`|GS_{}^N=|C_{NK}^K;|GS_{}^N=|C_{NK1}^{K+1}.`$ (15)
This is the simplest case in which
$`X_,=(1)^Kd_{NK1,}^{}d_{K,}`$ (16)
and similarly for the other Hubbard operators. We have also assumed, $`|N1=|C_{NK1}^K`$ and $`|N+1=|C_{NK}^{K+1}`$ (in other case the quantum numbers of the $`d^{}`$ operators must be changed accordingly) and the signs are taken according with definitions (10). Due to the lack of correlation effects in the GS’s, the tunneling spectral amplitudes are
$`\mathrm{\Delta }_,=\mathrm{\Delta }_{+,}=V_{NK1};\mathrm{\Delta }_,=\mathrm{\Delta }_{+,}=V_K.`$ (17)
There is a difference between the antiferromagnetic couplings $`J`$, $`J_{}`$ and $`J_{}`$ due to the fact that $`NK1>K`$. It is easier to tunnel from the leads to spin up states ($`m=NK1`$) within the QD because they are in the outer region of the QD while the first available spin down state ($`m=K`$) is in the inner region of the QD. This tunneling amplitude effect is not described by previous models of Kondo at finite B , but has long been more broadly recognized both experimentally and theoretically for QD with $`2>\nu >1`$.
Although in this regime the QD does not present any correlation effects, the antiferromagnetic couplings are still a function of total energies and tunneling amplitudes. Since the Kondo temperature depends exponentially on $`J`$, it is a very sensitive magnitude with respect to many parameters as $`w`$, $`\omega _0`$, $`B`$, $`V_m`$ etc. However, this is not different from the case of just one electron in the QD in which experiments show $`T_K`$ to be in the range of 1K. In any case, one can predict variations in $`T_K`$ when the regime changes as discussed below.
ii) (Lower part of Fig. 1) When the Zeeman interaction is small enough, the GS’s are skyrmion-like states of topological charge one :
$`|GS_{}^N=|SK_{N,K,\pm }^P;|GS_{}^N=|SK_{N,K,\pm }^{P+1}`$ (18)
The Hubbard operators $`X_{\mathrm{\Sigma },\mathrm{\Sigma }^{}}`$ have complicated, but analytical, expressions in terms of $`\mathrm{\Lambda }_{\pm 1}^{}`$.
When $`N`$ is very large, correlation effects provoke that $`\mathrm{\Delta }_{,\mathrm{\Sigma }}`$ tends to zero (for instance $`\mathrm{\Delta }_{,\mathrm{\Sigma }}[NlnN]^{1/2}`$ for $`P=0`$) implying a quenching of the antiferromagnetic couplings due to the orthogonalization catastrophe. As a consequence, Kondo effect should not be observed when the number of electrons within the QD is large because the Kondo temperature is extremely small in this case.
Even more interesting is the effect for a reduced number of electrons, for instance $`N=8`$ as both in Fig. 1, 2 and in the experiment . In this case, correlation effects do not destroy Kondo effect but reduce $`\mathrm{\Delta }_{,\mathrm{\Sigma }}`$ up to a factor of two . The reduction of $`J`$ and $`J_\mathrm{\Sigma }`$ implies that Kondo temperatures are significantly smaller than for the case i) of compact states. In practice, one can move from the skyrmions regime ii) (lower part of Fig 1) to the compact regime i) (upper part of Fig 1) by increasing an inplane component of the magnetic field. As a consequence of the above analysis, one should detect a clearly measurable increase of $`T_K`$ during this process.
iii) There is a rather unusual case, in which, for intermediate Zeeman interaction and only for a few values of the magnetic field, the degenerate GS’s are one skyrmion and one compact state, i. e.:
$`|GS_{}^N=|SK_{N,K,+}^P;|GS_{}^N=|C_{NK+1}^{K1}.`$ (19)
This happens, for instance, in the last step to the right in Fig. 2. In this case, the two GS have spins differing in more than 1. Therefore, they can not be obtained from the same $`|N1`$ by creating one electron with spin either up or down. In other words, one of the two GS’s has a tunneling spectral amplitude equal to zero so that, for tunneling effects there is only a non degenerate GS and Kondo effect does not occur.
A practical difference between the compact and skyrmion cases resides in the probability of finding two degenerate (or almost) GS’s. As shown in Fig. 1, the crossing regions are much closer to each other for skyrmions (lower part of Fig. 1) than for compact states (upper Fig. 1). To work with a magnetic field in which Kondo effect exists is more probable in the skyrmion regime than in the compact one. However, once Kondo regime is found in the compact regime, the effect is more robust because its $`T_K`$ is higher than the one for skyrmions.
## V Chess board behavior of the conductance
A very characteristic feature of some experiments is an alternating high-low conductance sequence as a function of $`B`$ for a given temperature and number of electrons. When $`N`$ is varied in $`\pm 1`$, the high and low conductance valleys are interchanged. Since a common way of representing the experimental results is to use a color-intensity scale for the conductance as a function of both $`B`$ and a gate potential which varies $`N`$, the data present a chess board aspect. This occurs in a broad range of filling factors and number of electrons.
In order to understand such an experimental behavior, let us start by applying our description for the regime $`1<\nu <2`$. Since the chess board is observed in samples with large $`w`$, skyrmion-like states are high energy excitations. GS’s are restricted to be given just by one configuration, i. e. $`[S^{()}]^j|C_{NK}^K`$ states, where $`j`$ is an integer (for practical purposes one can concentrate in the cases $`j=0`$ and $`1`$). At first sight, the consideration of both $`|C_{NK}^K`$ and $`S^{()}|C_{NK}^K`$ seems unnecessary because these states have an energy difference equal to the Zeeman coupling so that they are never degenerate. However, in a broad range of magnetic fields, these two states are precisely the two lowest energy states (see e.g. Fig. 2).
In the two level system approach in which the two lowest states are separated by an energy splitting $`\mathrm{\Delta }E`$, Kondo-like behavior appears only when the experimental temperature is between $`(\mathrm{\Delta }E)^2/T_K`$ and $`T_K`$ (provided $`\mathrm{\Delta }E<T_K`$) in order to have significant occupation of the two states. The quantitative application of the model would require the computation of both the energy splitting $`\mathrm{\Delta }E`$ and $`T_K`$ which, as mentioned above, are very sensitive to many experimental parameters. Instead, a general understanding of the chess board is obtained from the following qualitative explanation depicted schematically in Fig. 4: when magnetic field is varied, the two lowest levels are in two alternating situations depending whether the system is around or far from a crossing of GS’s. Around a crossing, the two lowest levels correspond to states $`|C_{NK1}^{(K+1)}`$ and $`|C_{NK}^K`$ while far from the crossing they are $`|C_{NK}^K`$ and $`S^{()}|C_{NK}^K`$. In the former case, the two lowest levels can be connected by the tunneling Hamiltonian through low energy states $`|N\pm 1`$ with $`\mathrm{\Delta }_{\pm ,\mathrm{\Sigma }}=V_KV_{NK1}`$. In this case, the system presents high Kondo-like conductance. In the second case of being far from the crossings, the two lowest states are $`|C_{NK}^K`$ and $`S^{()}|C_{NK}^K`$. The coupling $`J`$ between the two states is very weak for two reasons: first, the connecting states $`|N\pm 1`$ are highly excited states implying very large denominators in $`J`$ and, second, since $`S^{()}|C_{NK}^K`$ is a linear combination of different configurations, $`\mathrm{\Delta }_{\pm ,\mathrm{\Sigma }}`$ is significantly reduced by factors proportional to $`1/N`$. These two effects produce, in this region, a very small $`J`$ and, consequently, an exponentially negligible $`T_K`$ so that the measured conductance is not Kondo-like but instead is quenched. This originates the alternating behavior experimentally observed for fix gate voltage (i. e. number of electrons in the QD) and varying magnetic field. The chess board aspect also implies alternating low-high conductance regions for fix $`B`$ and varying gate voltage, which is due to the fact that crossings for $`N\pm 1`$ electrons occur for magnetic fields roughly midway from crossings for $`N`$ electrons.
Since Kondo-like behavior only occurs when the experimental temperature is in the range between $`(\mathrm{\Delta }E)^2/T_K`$ and $`T_K`$, a clear prediction of our scheme is that the highly conducting regions of the chess board would become narrower for decreasing temperature due to the lower limit condition.
The sequence of alternating high-low conductance has been observed also at very low magnetic field ($`B0.5`$T) and large number ($`N50`$) of electrons. So, a complete understanding of the experimental situation requires the extension of our framework to $`\nu >2`$. We are currently involved in this task which implies to take into account GS’s more complicated than simple compact states in the lowest Landau level.
## VI Summary
We have presented a general description of the Kondo effect for any number of electrons in a QD at filling factor $`1<\nu <2`$. Collective spin effects of the $`N`$ electrons are described in terms of a set of spin-flip Hubbard operators which also change the third component of the total angular momentum. Our conclusion is that, for any number of electrons within the QD, the spin-flip scattering of carriers only produces transitions between two (either compact or skyrmion-like) states of the QD with spins differing in 1. This process can be described by a Kondo Hamiltonian with exchange couplings that, we show, are antiferromagnetic and depend on both correlation and tunneling amplitude effects. As a consequence, the system presents a Kondo behavior with an experimentally accessible $`T_K`$. The increase of $`T_K`$ with an inplane component of the magnetic field should allow the detection of the transition from a skyrmion-like regime to a compact state regime. Finally, we present a qualitative explanation for the chess board aspect of the experimental conductance when represented in a grey scale as a function of both the magnetic field and the gate potential affecting the quantum dot.
## VII Acknowledgements
We are grateful to J. Weis and L. P. Kouwenhoven for providing us with experimental information before its publication. This work was supported in part by MEC of Spain under contract No. PB96-0085, Fundacion Ramon Areces and CAM under contract No. 07N/0026/1998.
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# On 𝛼⁺⁺-Stable Graphs
## 1 Introduction
All the graphs considered in this paper have at least two vertices. For such a graph $`G=(V,E)`$ we denote its vertex set by $`V=V(G)`$ and its edge set by $`E=E(G).`$ If $`XV`$, then $`G[X]`$ is the subgraph of $`G`$ spanned by $`X`$. By $`GW`$ we mean the subgraph $`G[VW]`$ , if $`WV(G)`$. By $`GF`$ we denote the partial subgraph of $`G`$ obtained by deleting the edges of $`F`$, for $`FE(G)`$, and we use $`Ge`$, if $`W`$ $`=\{e\}`$. If $`A,B`$ $`V`$ and $`AB=\mathrm{}`$, then $`(A,B)`$ stands for the set $`\{e=ab:aA,bB,eE\}`$. The neighborhood of a vertex $`vV`$ is the set $`N(v)=\{w:wV`$ and $`vwE\}`$, and $`N(A)=\{N(v):vA\}`$, for $`AV`$. If $`\left|N(v)\right|=\left|\{w\}\right|=1`$, then $`v`$ is a pendant vertex and $`vw`$ is a pendant edge of $`G`$. By $`C_n`$, $`K_n`$, $`P_n`$ we denote the chordless cycle on $`n`$ $`4`$ vertices, the complete graph on $`n2`$ vertices, and respectively the chordless path on $`n3`$ vertices.
A stable set of maximum size will be referred as to a maximum stable set of $`G`$. The stability number of $`G`$, denoted by $`\alpha (G)`$, is the cardinality of a maximum stable set in $`G`$. Let $`\mathrm{\Omega }(G)`$ denotes $`\{S:S`$ is a maximum stable set of $`G\}`$ and $`\xi (G)=\left|\{S:S\mathrm{\Omega }(G)\}\right|`$. We call $`\{\mathrm{\Omega }_1,\mathrm{\Omega }_2\}`$ a cover of $`\mathrm{\Omega }(G)`$ if $`\mathrm{\Omega }_1,\mathrm{\Omega }_2\mathrm{\Omega }(G)`$ and $`\mathrm{\Omega }_1\mathrm{\Omega }_2=\mathrm{\Omega }(G)`$; by $`\xi (\mathrm{\Omega }_i),i=1,2`$, we mean the number $`\left|\{S:S\mathrm{\Omega }(G_i)\}\right|`$.
A matching is a set of non-incident edges of $`G`$; a matching of maximum cardinality $`\mu (G)`$ is a maximum matching, and a perfect matching is a matching covering all the vertices of $`G`$. $`G`$ is a König-Egerváry graph provided $`\alpha (G)+\mu (G)=\left|V(G)\right|`$, , .
A graph $`G`$ is $`\alpha ^+`$-stable if $`\alpha (G+e)=\alpha (G)`$, for any edge $`eE(\overline{G})`$, where $`\overline{G}`$ is the complement of $`G`$, . Haynes et al. have characterized the $`\alpha ^+`$-stable as follows:
###### Theorem 1.1
A graph $`G`$ is $`\alpha ^+`$-stable if and only if $`\xi (G)1`$.
Theorem 1.1 implies that for an $`\alpha ^+`$-stable graph either $`\xi (G)=0`$ or $`\xi (G)=1`$. This motivates the following definition. A graph $`G`$ is called ($`i`$) $`\alpha _0^+`$-stable whenever $`\xi (G)=0`$, and ($`\mathrm{𝑖𝑖}`$) $`\alpha _1^+`$-stable provided $`\xi (G)=1`$, . For instance, $`C_4`$ is $`\alpha _0^+`$-stable, while the graphs $`K_3+e,K_4+e`$ in Figure 1 are $`\alpha _1^+`$-stable.
In it was shown that an $`\alpha ^+`$-stable tree can be only $`\alpha _0^+`$-stable, and this is exactly the case of trees possessing a perfect matching. This result was generalized to bipartite graphs in . Nevertheless, there exist both $`\alpha _1^+`$-stable König-Egerváry graphs (e.g., the graph $`K_3+e`$ in Figure 1), and $`\alpha _0^+`$-stable König-Egerváry graphs (e.g., all $`\alpha ^+`$-stable bipartite graphs). A necessary (but not sufficient, e.g., $`K_4e`$) condition for $`\alpha ^+`$-stability is:
###### Proposition 1.2
Any $`\alpha ^+`$-stable König-Egerváry graph has a perfect matching.
Let define a graph $`G`$ as $`\alpha ^{++}`$-stable whenever $`\alpha (G+e_1+e_2)=\alpha (G)`$, for any $`e_1,e_2E(\overline{G})`$, and $`\alpha _{P_3}^+`$-stable provided $`\alpha (G+e_1+e_2)=\alpha (G)`$, for any $`e_1,e_2E(\overline{G})`$ that have a common endpoint. Gunther et al., , studied the structure of $`\alpha ^+`$-stable trees, and in , , some of their results were extended to bipartite graphs and König-Egerváry graphs.
In this paper we characterize $`\alpha ^{++}`$-stable graphs. We settle a number of connections between $`\alpha ^+`$-stable graphs, $`\alpha _{P_3}^+`$-stable and $`\alpha ^{++}`$-stable graphs. In particular, we show that any $`\alpha _0^+`$-stable graph is $`\alpha _{P_3}^+`$-stable. We also give a necessary and sufficient condition for a graph to be $`\alpha _1^+`$-stable and $`\alpha _{P_3}^+`$-stable at the same time.
We prove that a König-Egerváry graph is $`\alpha ^{++}`$-stable if and only if it has a perfect matching consisting of only pendant edges and contains no cycle on $`4`$ vertices. Using this result we describe $`\alpha ^{++}`$-stable bipartite graphs and $`\alpha ^{++}`$-stable trees. For instance, it is shown that a bipartite graph is $`\alpha ^{++}`$-stable if and only if it is well-covered and $`C_4`$-free.
## 2 $`\alpha ^{++}`$-stable graphs
Notice that any $`\alpha ^{++}`$-stable graph is $`\alpha _{P_3}^+`$-stable, but the converse is not generally true. For instance, $`C_4`$ is $`\alpha _{P_3}^+`$-stable but not $`\alpha ^{++}`$-stable. Let us also observe that for $`n2,`$ the graph $`K_ne`$ is $`\alpha _{P_3}^+`$-stable, but it is not $`\alpha ^+`$-stable.
###### Proposition 2.1
If $`GK_ne,n2`$ is $`\alpha _{P_3}^+`$-stable, then it is also $`\alpha ^+`$-stable.
Proof. Assume that $`\alpha (G)=2`$. Since $`GK_ne`$, it follows that $`\left|\mathrm{\Omega }(G)\right|2`$, and consequently $`G`$ is $`\alpha ^+`$-stable, as well. For $`\alpha (G)3`$, suppose, on the contrary, that $`G`$ is not $`\alpha ^+`$-stable, and let $`x,y\{S:S\mathrm{\Omega }(G)\}`$. Hence, for $`S\mathrm{\Omega }(G)`$ and $`zS\{x,y\}`$, we obtain that $`\alpha (G+xy+xz)<\alpha (G)`$, in contradiction with the fact that $`G`$ is $`\alpha _{P_3}^+`$-stable. Therefore, $`G`$ must be $`\alpha ^+`$-stable.
It is worth observing that an $`\alpha ^+`$-stable graph is not necessarily $`\alpha _{P_3}^+`$-stable. For instance, $`K_3+e`$ is $`\alpha ^+`$-stable, in fact it is $`\alpha _1^+`$-stable, but it is not $`\alpha _{P_3}^+`$-stable. However, there exist graphs that are both $`\alpha _1^+`$-stable and $`\alpha _{P_3}^+`$-stable; e.g., the graph $`K_4+e`$ in Figure 1.
###### Proposition 2.2
Any $`\alpha _0^+`$-stable graph is also $`\alpha _{P_3}^+`$-stable.
Proof. Suppose, on the contrary, that some $`\alpha _0^+`$-stable graph $`G`$ is not $`\alpha _{P_3}^+`$-stable, i.e., there are $`x,y,zV`$ such that $`\alpha (G+xy+xz)<\alpha (G)`$. Hence, it follows that $`x\{S:S\mathrm{\Omega }(G)\}`$, in contradiction with the fact that $`G`$ is $`\alpha _0^+`$-stable.
Combining Theorem 1.1 and Propositions 2.1, 2.2, we get:
###### Corollary 2.3
If $`GK_ne`$ has $`\alpha (G)=2`$, then $`G`$ is:
($`i`$) $`\alpha ^+`$-stable if and only if $`\left|\mathrm{\Omega }(G)\right|2`$;
($`\mathrm{𝑖𝑖}`$) $`\alpha _{P_3}^+`$-stable if and only if either it is $`\alpha _0^+`$-stable or it is $`\alpha _1^+`$-stable and $`\left|\mathrm{\Omega }(G)\right|3`$;
($`\mathrm{𝑖𝑖𝑖}`$) $`\alpha ^{++}`$-stable if and only if $`\left|\mathrm{\Omega }(G)\right|3`$;
($`\mathrm{𝑖𝑣}`$) $`\alpha _{P_3}^+`$-stable if and only if it is $`\alpha ^+`$-stable and $`GP_3(K_1,K_m,K_2)`$, where $`P_3(K_1,K_m,K_2)`$ is the graph obtained by substituting the vertices of $`P_3`$ respectively, by $`K_1,K_m,K_2`$, and joining all the vertices of $`K_m`$ with the two vertices of $`K_2`$ and the single vertex of $`K_1`$.
###### Proposition 2.4
If $`G=(V,E)`$ has $`\alpha (G)3`$, then the following assertions are equivalent:
($`i`$) $`G`$ is $`\alpha ^+`$-stable;
($`\mathrm{𝑖𝑖}`$) either $`G`$ is $`\alpha _{P_3}^+`$-stable, or there exist three vertices $`x,y,zV`$ such that $`\left|\{x,y\}S\right|\left|\{x,z\}S\right|2`$ holds for any $`S\mathrm{\Omega }(G)`$, and $`x`$ is the unique vertex of $`G`$ with this property.
Proof. ($`i`$) $``$ ($`\mathrm{𝑖𝑖}`$) Let $`G`$ be $`\alpha ^+`$-stable, but not $`\alpha _{P_3}^+`$-stable. Hence, there are $`x,y,zV`$ such that $`\alpha (G+xy+xz)<\alpha (G)`$. Therefore, we get that $`x\{S:S\mathrm{\Omega }(G)\}`$, because otherwise, any $`S\mathrm{\Omega }(G)`$ not containing $`x`$ is still stable in $`G+xy+xz`$, and consequently, we obtain $`\alpha (G+xy+xz)=\alpha (G)`$, in contradiction with the assumption on $`G`$. In addition, each $`S\mathrm{\Omega }(G)`$ satisfies $`\left|S\{y,z\}\right|1`$, since otherwise, if some $`S_0\mathrm{\Omega }(G)`$ has $`S_0\{y,z\}=\mathrm{}`$, then $`S_0`$ is stable in $`G+xy+xz`$ and this yields $`\alpha (G+xy+xz)=\alpha (G)`$, again in contradiction with the assumption on $`G`$. Finally, $`x`$ is unique, because otherwise $`\xi (G)2`$, which contradicts the $`\alpha ^+`$-stability of $`G`$.
($`\mathrm{𝑖𝑖}`$) $``$ ($`i`$) If $`G`$ is $`\alpha _{P_3}^+`$-stable and $`\alpha (G)3`$, then by Proposition 2.1, $`G`$ is $`\alpha ^+`$-stable. Further, if there are $`x,y,zV`$ such that $`\left|\{x,y\}S\right|\left|\{x,z\}S\right|2`$ holds for any $`S\mathrm{\Omega }(G)`$, and $`x`$ is unique with this property, it follows that $`G`$ is not $`\alpha _{P_3}^+`$-stable, because $`\alpha (G+xy+xz)<\alpha (G)`$, but it is $`\alpha _1^+`$-stable, since $`\{x\}=\{S:S\mathrm{\Omega }(G)\}`$.
###### Lemma 2.5
If for any $`x,yV(G)`$ there exists $`S\mathrm{\Omega }(G)`$ such that $`x,yV(G)S`$, then $`G`$ is both $`\alpha _0^+`$-stable and $`\alpha ^{++}`$-stable.
Proof. Suppose, on the contrary, that there exists $`x\{S:S\mathrm{\Omega }(G)\}`$. Then, for any $`yV(G)\{x\}`$ and $`S\mathrm{\Omega }(G)`$, we get $`\{x,y\}S\mathrm{}`$, in contradiction with the premises on $`G`$. Therefore, $`G`$ is $`\alpha _0^+`$-stable. According to Proposition 2.2, it follows that $`G`$ is $`\alpha _{P_3}^+`$-stable, too. Assume that $`G`$ is not $`\alpha ^{++}`$-stable. Hence, since $`G`$ is $`\alpha _{P_3}^+`$-stable, we infer that there are $`x,y,u,vV(G)`$, pairwise distinct, such that $`\alpha (G+xy+uv)<\alpha (G)`$. Let $`S\mathrm{\Omega }(G)`$ be such that $`x,uV(G)S`$. Then $`S`$ is stable in $`G+xy+uv`$, in contradiction with the assumption on $`G`$. Consequently, $`G`$ is $`\alpha ^{++}`$-stable.
As an example, $`C_{2k+1},k2`$ and $`K_n,n3`$ are both $`\alpha _0^+`$-stable and $`\alpha ^{++}`$-stable, according to Lemma 2.5. Notice that the converse of Lemma 2.5 is not generally true; see, for instance, the graphs $`C_{2k},k3`$. There exist $`\alpha _0^+`$-stable graphs that are not $`\alpha ^{++}`$-stable (e.g., $`C_4`$), and vice-versa, there are $`\alpha ^{++}`$-stable that are not $`\alpha _0^+`$-stable (e.g., $`K_4+e`$).
###### Proposition 2.6
Let $`G`$ be $`\alpha _1^+`$-stable, $`\{v\}=\{S:S\mathrm{\Omega }(G)\}`$ and $`G_0=GN[v]`$. If $`G`$ is not $`\alpha _{P_3}^+`$-stable, then there are $`x`$ and $`y`$ belonging to the same connected component of $`G_0`$, such that $`\alpha (G+xv+yv)<\alpha (G)`$. In other words, there exists a path connecting $`x`$ and $`y`$, which avoid the neighborhood of $`v`$.
Proof. Let $`\{H_k:1ks\},s2`$, be the connected components of $`G_0`$, and suppose that there are $`x`$ and $`y`$ belonging to different connected components of $`G_0`$, (say respectively $`H_i,H_j`$), such that $`\alpha (G+xv+yv)<\alpha (G)`$. Since $`G_0`$ is $`\alpha _0^+`$-stable, it follows that all its connected components are $`\alpha _0^+`$-stable, as well. Let $`S_k`$ $`\mathrm{\Omega }(H_k),1ks`$, and $`S_i`$ $`\mathrm{\Omega }(H_i),S_j`$ $`\mathrm{\Omega }(H_j)`$ be such that $`xS_i,yS_j`$, which exist, because all $`H_k`$ are $`\alpha _0^+`$-stable. Hence, we get that $`\{v\}(\{S_k:1ks\})`$ is stable in $`G+xv+yv`$, in contradiction with $`\alpha (G+xv+yv)<\alpha (G)`$.
###### Theorem 2.7
Let $`G`$ be $`\alpha _1^+`$-stable, $`\{v\}=\{S:S\mathrm{\Omega }(G)\}`$ and $`G_0=GN[v]`$. Then $`G`$ is $`\alpha _{P_3}^+`$-stable if and only if for every pair $`x,yV(G_0)`$ there exists $`S_0\mathrm{\Omega }(G_0)`$ such that $`x,yV(G_0)S_0`$.
Proof. Let $`x,yV(G_0)`$. By definition of $`G_0`$, it follows that $`x,yN[v]`$, and since $`G`$ is $`\alpha _{P_3}^+`$-stable, we get that $`\alpha (G+xv+yv)=\alpha (G)`$. Therefore, there is $`S\mathrm{\Omega }(G)`$ such that $`x,yV(G)S`$. Hence, $`x,yV(G_0)S_0`$, where $`S_0=S\{v\}\mathrm{\Omega }(G_0)`$.
Conversely, $`G`$ is $`\alpha _1^+`$-stable, and for every pair $`x,yV(G_0)`$ there exists $`S\mathrm{\Omega }(G_0)`$ with $`x,yV(G_0)S_0`$. Assume that $`G`$ is not $`\alpha _{P_3}^+`$-stable. Hence, there are $`x,yV(G)`$ such that $`\alpha (G+xv+yv)<\alpha (G)`$, since $`G`$ is $`\alpha _1^+`$-stable. Let $`S_0\mathrm{\Omega }(G_0)`$ be such that $`x,yV(G_0)S_0`$. Then, it follows that $`S_0\{v\}\mathrm{\Omega }(G)\}`$, in contradiction with the assumption on $`G`$. Therefore, $`G`$ is also $`\alpha _{P_3}^+`$-stable.
###### Proposition 2.8
A graph $`G`$ is not $`\alpha _{P_3}^+`$-stable if and only if $`\xi (G)1`$ and there exists a cover $`\{\mathrm{\Omega }_1,\mathrm{\Omega }_2\}`$ of $`\mathrm{\Omega }(G)`$, such that $`\xi (\mathrm{\Omega }_i)2,i=1,2`$.
Proof. If $`G`$ is not $`\alpha _{P_3}^+`$-stable, then $`\alpha (G+e_1+e_2)<\alpha (G)`$ holds for some $`e_1,e_2E(\overline{G})`$ that have a common endpoint. Suppose $`e_1=xy,e_2=yz`$. Let us define
$$\mathrm{\Omega }_1=\{S:x,yS\mathrm{\Omega }(G)\}and\mathrm{\Omega }_2=\{S:y,zS\mathrm{\Omega }(G)\}.$$
Hence, it follows that $`\xi (G)1`$ and $`\xi (\mathrm{\Omega }_i)2,i=1,2`$.
Conversely, assume that $`\xi (G)1`$, i.e., there exists at least one vertex belonging to $`\{S:S\mathrm{\Omega }(G)\}`$, say $`y`$, and that there is some cover $`\{\mathrm{\Omega }_1,\mathrm{\Omega }_2\}`$ of $`\mathrm{\Omega }(G)`$, such that $`\xi (\mathrm{\Omega }_i)2,i=1,2`$. If $`x\{S:S\mathrm{\Omega }_1\}\{y\}`$ and $`v\{S:S\mathrm{\Omega }_2\}\{y\}`$, then $`\alpha (G+xy+uv)<\alpha (G)`$, because any $`S\mathrm{\Omega }(G)`$ contains at least one of the pairs $`\{x,y\}`$ or $`\{y,v\}`$. Therefore, $`G`$ can not be $`\alpha _{P_3}^+`$-stable.
###### Proposition 2.9
A graph $`G`$ is not $`\alpha ^{++}`$-stable if and only if there exists a cover $`\{\mathrm{\Omega }_1,\mathrm{\Omega }_2\}`$ of $`\mathrm{\Omega }(G)`$, such that $`\xi (\mathrm{\Omega }_i)2,i=1,2`$.
Proof. If $`G`$ is not an $`\alpha ^{++}`$-stable graph, then $`\alpha (G+e_1+e_2)<\alpha (G)`$ holds for some $`e_1,e_2E(\overline{G})`$. Suppose $`e_1=xy,e_2=uv`$. Let us define $`\mathrm{\Omega }_1=\{S:x,yS\mathrm{\Omega }(G)\}`$ and $`\mathrm{\Omega }_2=\{S:u,vS\mathrm{\Omega }(G)\}`$. Hence, it follows that $`\xi (\mathrm{\Omega }_i)2,i=1,2`$. Suppose that there exists $`S\mathrm{\Omega }\left(\mathrm{\Omega }_1\mathrm{\Omega }_2\right)`$. Then $`S\mathrm{\Omega }\left(G+e_1+e_2\right)`$, that contradicts the inequality $`\alpha (G+e_1+e_2)<\alpha (G)`$. Hence, $`\mathrm{\Omega }_1\mathrm{\Omega }_2=\mathrm{\Omega }`$, which means that $`\{\mathrm{\Omega }_1,\mathrm{\Omega }_2\}`$ is a cover we were supposed to find.
Conversely, assume that there is a cover $`\{\mathrm{\Omega }_1,\mathrm{\Omega }_2\}`$ of $`\mathrm{\Omega }(G)`$ with $`\xi (\mathrm{\Omega }_i)2,i=1,2`$. If $`x,y\{S:S\mathrm{\Omega }_1\}`$ and $`u,v\{S:S\mathrm{\Omega }_2\}`$, then $`\alpha (G+xy+uv)<\alpha (G)`$, because any $`S\mathrm{\Omega }(G)`$ contains at least one of the pairs $`\{x,y\}`$ or $`\{u,v\}`$. Therefore, $`G`$ is not $`\alpha ^{++}`$-stable.
Combining Propositions 2.4 and 2.9, we deduce the following:
###### Theorem 2.10
For a graph $`G`$ the following assertions are equivalent:
($`i`$) $`G`$ is $`\alpha ^{++}`$-stable;
($`\mathrm{𝑖𝑖}`$) $`G`$ is $`\alpha ^+`$-stable and $`\mathrm{\Omega }(G+e_1)\mathrm{\Omega }(G+e_1)\mathrm{}`$ for any $`e_1,e_2E(\overline{G})`$;
($`\mathrm{𝑖𝑖𝑖}`$) $`\mathrm{\Omega }(G)\mathrm{\Omega }(G+e_1)\mathrm{\Omega }(G+e_1)\mathrm{}`$ for any $`e_1,e_2E(\overline{G})`$;
($`\mathrm{𝑖𝑣}`$) $`G`$ is $`\alpha ^+`$-stable and $`|\{S:S\mathrm{\Omega }(G+e)|1`$ for any $`eE(\overline{G})`$;
($`v`$) $`G`$ is $`\alpha _{P_3}^+`$-stable and there are no $`e_1,e_2E(\overline{G}),e_1=xy,e_2=uv`$, such that:
$$\{x,y\}\{u,v\}=\mathrm{},andforanyS\mathrm{\Omega }(G),\mathrm{max}\{\left|S\{x,y\}\right|,\left|S\{u,v\}\right|\}=2;$$
($`\mathrm{𝑣𝑖}`$) for any cover $`\{\mathrm{\Omega }_1,\mathrm{\Omega }_2\}`$ of $`\mathrm{\Omega }(G)`$ either $`\xi (\mathrm{\Omega }_1)1`$ or $`\xi (\mathrm{\Omega }_2)1`$ holds.
## 3 $`\alpha ^{++}`$-stable König-Egerváry graphs
According to a well-known result of König, , and Egerváry, , any bipartite graph is a König-Egerváry graph. This class includes also non-bipartite graphs (see, for instance, the graph $`K_3+e`$ in Figure 1). If $`G_i=(V_i,E_i),i=1,2`$, are two disjoint graphs, then $`G=G_1G_2`$ is defined as the graph with $`V(G)=V(G_1)V(G_2)`$, and
$$E(G)=E(G_1)E(G_2)\{xy:forsomexV(G_1)andyV(G_2)\}.$$
###### Proposition 3.1
The following assertions are equivalent:
($`i`$) $`G`$ is a König-Egerváry graph;
($`\mathrm{𝑖𝑖}`$) $`G=H_1H_2`$, where $`V(H_1)=S\mathrm{\Omega }(G)`$ and $`\left|V(H_1)\right|\mu (G)=\left|V(H_2)\right|`$;
($`\mathrm{𝑖𝑖𝑖}`$) $`G=H_1H_2`$, where $`V(H_1)=S`$ is a stable set in $`G,`$ $`\left|S\right|\left|V(H_2)\right|`$ and $`(S,V(H_2))`$ contains a matching $`M`$ with $`\left|M\right|=\left|V(H_2)\right|`$.
It it easy to see that a König-Egerváry graph $`G`$ has a perfect matching if and only if $`\alpha (G)=\mu (G)`$. The edges of any maximum matching of a König-Egerváry graph have a specific position with respect to the maximum stable sets.
###### Lemma 3.2
If $`G`$ is a König-Egerváry graph, then for any $`S\mathrm{\Omega }(G)`$ every maximum matching of $`G`$ is contained in $`(S,V(G)S)`$.
###### Proposition 3.3
If $`M`$ is a maximum matching in a graph $`G`$ and $`H`$ is a subgraph of $`G`$ such that $`M=(ME(H))(ME(GH))`$, then $`\mu (G)=\mu (H)+\mu (GH)`$.
Proof. Clearly, $`ME(H)`$ and $`ME(GH)`$ are matchings in $`H`$ and $`GH`$, respectively. Let $`M_1,M_2`$ be maximum matchings in $`H`$ and $`GH`$, respectively. If $`\mu (H)=\left|M_1\right|>\left|ME(H)\right|`$, or $`\mu (GH)=\left|M_2\right|>\left|ME(GH)\right|`$, then $`M_1M_2`$ is a matching in $`G`$ of cardinality larger than $`\left|M\right|`$, in contradiction with $`\left|M\right|=\mu (G)`$. Therefore, $`\mu (G)=\mu (H)+\mu (GH)`$.
###### Proposition 3.4
If $`M`$ is a maximum matching in a König-Egerváry graph $`G`$, and $`H`$ is a subgraph of $`G`$ such that $`M=(ME(H))(ME(GH))`$, then
($`i`$) $`H`$ and $`GH`$ are König-Egerváry graphs;
($`\mathrm{𝑖𝑖}`$) $`\alpha (G)=\alpha (H)+\alpha (GH)`$.
Proof. Let $`S\mathrm{\Omega }(G),S_1=SV(H)`$ and $`S_2=SV(GH)`$. By Lemma 3.2, $`M(S,V(G)S)`$, and according to Proposition 3.1($`\mathrm{𝑖𝑖}`$), $`G=H_1H_2`$, where $`V(H_1)=S\mathrm{\Omega }(G)`$ and $`\left|V(H_1)\right|\mu (G)=\left|M\right|=\left|V(H_2)\right|`$. Hence, we infer that: $`V(H)=S_1(V(H_2)V(GH)),S_1`$ is stable in $`H,ME(H)`$ is a matching in $`H`$ of size $`\left|V(H_2)V(GH)\right|`$, and $`\left|S_1\right|\left|V(H_2)V(GH)\right|`$, i.e., $`H`$ is a König-Egerváry graph, according to Proposition 3.1($`\mathrm{𝑖𝑖𝑖}`$). Similarly, $`GH`$ is also a König-Egerváry graph. Since, by Proposition 3.3, $`\mu (G)=\mu (H)+\mu (GH)`$ and all $`G,H,GH`$ are König-Egerváry graphs, we may conclude that $`\alpha (G)=\alpha (H)+\alpha (GH)`$.
###### Lemma 3.5
If $`H`$ is a subgraph of $`G`$, such that $`\alpha (G)=\alpha (H)+\alpha (GH)`$ and $`G`$ is $`\alpha ^{++}`$-stable, then $`H`$ is $`\alpha ^{++}`$-stable, as well.
Proof. Since $`\alpha (G)=\alpha (H)+\alpha (GH),`$ it follows that any $`S\mathrm{\Omega }(G)`$ satisfies $`\left|SV(H)\right|=\alpha (H)`$. So, if $`\alpha (H+e_1+e_2)<\alpha (H)`$, for some $`e_1,e_2E(\overline{H})`$, it follows that $`\alpha (G+e_1+e_2)<\alpha (G)`$, as well.
###### Lemma 3.6
If $`G`$ is of order $`6`$, has a Hamiltonian path and $`\alpha (G)=3`$, then $`G`$ is not $`\alpha ^{++}`$-stable.
Proof. Suppose that $`V(G)=\{v_i:1i6\}`$ and $`(v_i,v_{i+1})E(G)`$ for any $`i\{1,\mathrm{},5\}`$. Then $`H=G+v_1v_3+v_4v_6`$ has $`\alpha (H)=2`$, i.e., $`G`$ is not $`\alpha ^{++}`$-stable.
###### Proposition 3.7
If $`GK_ne,n=2,3`$ is an $`\alpha ^{++}`$-stable König-Egerváry graph, then $`G`$ has a perfect matching consisting of only pendant edges.
Proof. If $`G`$ is $`\alpha ^{++}`$-stable then, clearly, $`G`$ is $`\alpha _{P_3}^+`$-stable too. Hence, by Proposition 2.1 if $`GK_ne`$ then it is also $`\alpha ^+`$-stable. It is not difficult to check that $`K_ne`$ can be a König-Egerváry graph only for $`n=2,3`$. Therefore, if a König-Egerváry graph $`GK_ne,n=2,3`$ is $`\alpha ^{++}`$-stable, then it is also $`\alpha ^+`$-stable. Now Proposition 1.2 ensures that $`G`$ has a perfect matching, say $`M=\{a_ib_i:1i\alpha (G)\}`$. According to Proposition 3.1 and Lemma 3.2, we may assume that $`S=\{a_i:1i\alpha (G)\}\mathrm{\Omega }(G)`$. We show that $`M`$ consists of only pendant edges. Suppose, on the contrary, that some $`a_kb_kM`$ is not pendant.
Case 1. There exists some $`b_i`$ such that $`a_kb_i,b_ib_kE(G)`$ (see Figure 2(a)). If $`H=G[\{a_k,b_i,a_i,b_k\}]`$, then Proposition 3.4($`\mathrm{𝑖𝑖}`$) implies that $`\alpha (G)=\alpha (H)+\alpha (GH)`$. Since $`H`$ is not $`\alpha ^{++}`$-stable, it follows, by Lemma 3.5, that $`G`$ could not be $`\alpha ^{++}`$-stable, in contradiction with the premises on $`G`$.
Case 2. There exist $`a_ib_iM`$ with $`a_kb_i,a_ib_kE(G)`$ (see Figure 2($`b`$)). If $`H=G[\{a_k,b_i,a_i,b_k\}]`$, then Proposition 3.4($`\mathrm{𝑖𝑖}`$) ensures that $`\alpha (G)=\alpha (H)+\alpha (GH)`$. Since $`H`$ is not $`\alpha ^{++}`$-stable, it follows, by Lemma 3.5, that $`G`$ could not be $`\alpha ^{++}`$-stable, in contradiction with the premises on $`G`$.
Case 3. There exist $`a_i,b_i,a_j,b_j`$, such that $`a_ib_i,a_jb_jM`$ and $`a_kb_i,a_jb_kE(G)`$. In addition, we can assume that $`b_ib_k,b_kb_jE(G)`$, otherwise we return to Case 1. Hence, $`H=G[\{a_i,a_k,a_j,b_i,b_k,b_j\}]`$ contains a path on $`6`$ vertices (see Figure 2($`c`$)). Since $`\alpha (H)=\left|\{a_i,a_k,a_j\}\right|=3`$, Lemma 3.6 implies that $`H`$ is not $`\alpha ^{++}`$-stable, and because $`\alpha (G)=\alpha (H)+\alpha (GH)`$ is true according to Proposition 3.4($`\mathrm{𝑖𝑖}`$), we get, by Lemma 3.5, that $`G`$ cannot be $`\alpha ^{++}`$-stable, in contradiction with the premises on $`G`$.
Thus, $`M`$ must consist of only pendant edges.
It is worth observing that : ($`a`$) Proposition 3.7 fails for non-König-Egerváry graphs; e.g., $`C_5`$ is $`\alpha ^{++}`$-stable and has no perfect matching; ($`b`$) the converse of Proposition 3.7, within the class of König-Egerváry graphs, is not generally true. For instance, the graph $`G_1`$ in Figure 3 has a perfect matching consisting of only pendant edges and it is not $`\alpha ^{++}`$-stable (because $`\alpha (G_1+ad+bc)<\alpha (G_1)`$), while the graph $`G_2`$ in the same figure has a perfect matching consisting of only pendant edges, and it is also $`\alpha ^{++}`$-stable.
However, we can show that:
###### Proposition 3.8
Any graph that has a perfect matching consisting of only pendant edges is $`\alpha _{P_3}^+`$-stable.
Proof. Let $`G`$ be a graph that has a perfect matching $`M=\{a_ib_i:1i\mu (G)\}`$, consisting of only pendant edges, and suppose that $`S_0=\{a_i:1i\mu (G)\}\mathrm{\Omega }(G)`$. Let denote $`H=G+e_1+e_2`$, where $`e_1,e_2E(\overline{G})`$ are such that they have a common endpoint, say $`e_1=uv,e_2=vw`$. We distinguish between the following cases:
Case 1. If $`u,v,wS_0`$, then $`S_0\mathrm{\Omega }(H)`$, and $`\alpha (H)=\alpha (G)`$.
Case 2. If $`u,vS_0`$ or $`u,wS_0`$, then $`S_0\mathrm{\Omega }(H)`$, and $`\alpha (H)=\alpha (G)`$.
Case 3. If $`uS_0`$ and $`v=a_i,w=a_j`$, then $`S_0\{b_j\}\{a_j\}\mathrm{\Omega }(H)`$, and $`\alpha (H)=\alpha (G)`$.
Case 4. If $`vS_0`$ and $`u=a_i,w=a_j`$, then $`S_0\mathrm{\Omega }(H)`$, and $`\alpha (H)=\alpha (G)`$.
Case 5. If $`u=a_i,v=a_j,w=a_k`$, then $`S_0\{b_j\}\{a_j\}\mathrm{\Omega }(H)`$, and $`\alpha (H)=\alpha (G)`$.
Consequently, $`G`$ is $`\alpha _{P_3}^+`$-stable.
###### Theorem 3.9
A graph $`G`$ that has a perfect matching consisting of only pendant edges is $`\alpha ^{++}`$-stable if and only if $`G`$ contains no cycle on $`4`$ vertices.
Proof. Let $`M=\{a_ib_i:1i\mu (G)\}`$ be the perfect matching of $`G`$. Without loss of generality, we can assume that $`S_0=\{a_i:1i\mu (G)\}\mathrm{\Omega }(G)`$.
Suppose, on the contrary, that there is
$$D=\{b_i,b_j,b_k,b_m\}withb_ib_j,b_jb_k,b_kb_m,b_mb_iE(G),$$
i.e., $`G[D]`$ contains a Hamiltonian cycle. If $`H=G[\{a_i,a_j,a_k,a_m\}D]`$, then $`\alpha (G)=\alpha (GH)+\alpha (H)`$, since any $`S\mathrm{\Omega }(G)`$ satisfies $`\left|S\{a_q,b_q\}\right|=1`$ for each $`a_qb_qM`$. On the one hand, by Lemma 3.5, $`H`$ should be $`\alpha ^{++}`$-stable. On the other hand, $`\alpha (H+a_ia_k+a_ja_m)=3<\alpha (H)=4`$, which brings a contradiction. Therefore, $`G`$ has no cycle on $`4`$ vertices.
Conversely, let $`G`$ be such that no $`4`$ vertices span a cycle. Assume, on the contrary, that $`G`$ is not $`\alpha ^{++}`$-stable, i.e., there are $`e_1,e_2E(\overline{G})`$ such that $`H=G+e_1+e_2`$ has $`\alpha (H)<\alpha (G)`$. If at least one of $`e_1,e_2`$ joins two vertices from $`\{b_i:1i\mu (G)\}`$, or one from $`\{a_i:1i\mu (G)\}`$ and the second from $`\{b_i:1i\mu (G)\}`$, then $`\alpha (H)=\alpha (G)`$. According to Proposition 3.8, the same result follows if $`e_1,e_2`$ have a common endpoint. Suppose that $`e_1=a_ia_j,e_2=a_ka_m`$ and $`a_i,a_j,a_k,a_m`$ are pairwise distinct. Hence, we get:
* $`b_ib_kE(G)`$, otherwise any $`S\mathrm{\Omega }(G)`$ containing $`\{a_j,a_m,b_i,b_k\}`$ is stable in $`H`$;
* $`b_jb_mE(G)`$, otherwise any $`S\mathrm{\Omega }(G)`$ containing $`\{a_i,a_k,b_j,b_m\}`$ is stable in $`H`$;
* $`b_ib_mE(G)`$, otherwise any $`S\mathrm{\Omega }(G)`$ containing $`\{a_j,a_k,b_i,b_m\}`$ is stable in $`H`$;
* $`b_jb_kE(G)`$, otherwise any $`S\mathrm{\Omega }(G)`$ containing $`\{a_i,a_m,b_j,b_k\}`$ is stable in $`H`$.
It follows that $`b_ib_k,b_jb_k,b_jb_m,b_ib_mE(G)`$, i.e., $`\{b_i,b_j,b_k,b_m\}`$ spans a $`4`$-cycle in $`G`$, in contradiction with the premises on $`G`$. Consequently, $`G`$ is $`\alpha ^{++}`$-stable.
Combining Proposition 3.7 and Theorem 3.9, we obtain the following characterization of $`\alpha ^{++}`$-stable König-Egerváry graphs.
###### Theorem 3.10
A König-Egerváry graph is $`\alpha ^{++}`$-stable if and only if it has a perfect matching consisting of only pendant edges and contains no cycle on $`4`$ vertices.
Recall that a graph $`G`$ is called: ($`a`$) well-covered if every maximal stable set of $`G`$ is also a maximum stable set, i.e., it is in $`\mathrm{\Omega }(G)`$, ; ($`b`$) very well-covered provided $`G`$ is well-covered and $`\left|V(G)\right|=2\alpha (G)`$, . The following result extends the characterization that Finbow, Hartnell and Nowakowski give in for well-covered graphs having the girth $`6`$.
###### Proposition 3.11
Let $`G`$ be a graph of girth $`6`$, which is isomorphic to neither $`C_7`$ nor $`K_1`$. Then the following assertions are equivalent:
($`i`$) $`G`$ is well-covered;
($`\mathrm{𝑖𝑖}`$) $`G`$ has a perfect matching consisting of pendant edges;
($`\mathrm{𝑖𝑖𝑖}`$) $`G`$ is very well-covered;
($`\mathrm{𝑖𝑣}`$) $`G`$ is a König-Egerváry $`\alpha _0^+`$-stable graph with exactly $`\alpha (G)`$ pendant vertices;
($`v`$) $`G`$ is a König-Egerváry $`\alpha ^{++}`$-stable graph.
Proof. The equivalences ($`i`$) $``$ ($`\mathrm{𝑖𝑖}`$) $``$ ($`\mathrm{𝑖𝑖𝑖}`$) are done in . In it has been proved that ($`\mathrm{𝑖𝑖𝑖}`$) $``$ ($`\mathrm{𝑖𝑣}`$). Finally, ($`\mathrm{𝑖𝑖}`$) $``$ ($`v`$) is true by Theorem 3.10.
###### Corollary 3.12
For a bipartite graph $`G`$ the following assertions are equivalent:
($`i`$) $`G`$ is $`\alpha ^{++}`$-stable;
($`\mathrm{𝑖𝑖}`$) $`G`$ is $`C_4`$-free and has a perfect matching consisting of only pendant edges;
($`\mathrm{𝑖𝑖𝑖}`$) $`G`$ is $`C_4`$-free and well-covered.
###### Corollary 3.13
$`C_n`$ is $`\alpha ^{++}`$-stable if and only if $`n`$ is odd.
Proof. For any $`n2,C_{2n}`$ is not an $`\alpha ^{++}`$-stable graph according to Corollary 3.12.
Assume, on the contrary, that $`C_{2n+1}`$ is not an $`\alpha ^{++}`$-stable graph. Hence, there are $`e_1,e_2E(\overline{C_{2n+1}}),e_1=xy,e_2=uv`$ such that $`\alpha (C_{2n+1}+e_1+e_2)<\alpha (C_{2n+1})`$. We may suppose, without loss of generality, that $`x=v_1u`$. Now, if $`u=v_{2i+1}`$ (for some $`i0`$), then $`x,uS=\{v_{2i}:1in\}\mathrm{\Omega }(C_{2n+1})`$, and if $`u=v_{2i}`$ (for some $`i0`$), then $`x,yS=\{v_2,v_4,\mathrm{},v_{2i2}\}\{v_{2i+1},v_{2i+3},\mathrm{},v_{2n+1}\}\mathrm{\Omega }(C_{2n+1})`$. Hence, we infer that $`S`$ is stable in $`C_{2n+1}+e_1+e_2`$, as well, in contradiction with $`\alpha (C_{2n+1}+e_1+e_2)<\alpha (C_{2n+1})`$. Therefore, $`C_{2n+1}`$ is $`\alpha ^{++}`$-stable.
Combining Corollary 3.12 and Proposition 3.11 we get the following extension of one Ravindra’s theorem, , where he proved the first three equivalences.
###### Corollary 3.14
For a tree $`T`$ the following assertions are equivalent:
($`i`$) $`T`$ is well-covered;
($`\mathrm{𝑖𝑖}`$) $`T`$ has a perfect matching consisting of pendant edges;
($`\mathrm{𝑖𝑖𝑖}`$) $`T`$ is very well-covered;
($`\mathrm{𝑖𝑣}`$) $`T`$ is $`\alpha ^{++}`$-stable.
## 4 Conclusions
In this paper we keep investigating graphs whose stability number is invariant with respect to some natural operations on graphs. While in , we were interested in measuring the influence of adding one edge to a graph, here we define a class of graphs whose stability number is unaffected by two edges addition.
Further we concentrate on König-Egerváry graphs, which is one of the most attractive generalizations of bipartite graphs. One the one hand, Proposition 3.11 claims that for girth $`6`$, $`\alpha ^{++}`$-stable König-Egerváry graphs and well-covered graphs are the same. On the other hand, Theorem 3.10 shows that an $`\alpha ^{++}`$-stable König-Egerváry graph contains no cycle on $`4`$ vertices. It leaves an interesting open question concerning interconnections between well-covered graphs and $`\alpha ^{++}`$-stable König-Egerváry graphs of girth $`3`$ or $`5`$.
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# About EAS size spectra and primary energy spectra in the knee region
## 1 Introduction
High statistical accuracy in modern EAS experiments in the knee region encouraged the investigation of the fine structure of EAS size spectra. Although the origin of the knee is still a matter of debate, recently the series of publications appeared, where the sharpness and the spectral structure in the knee region were interpreted by the contribution of heavy nuclei from a single local supernova. Along with this, the absolute differential EAS size spectra measured at different atmosphere depths and different zenith angles are not explained yet from the point of view of a single $`AA_{Air}`$ interaction model and a single model of primary energy spectra and elemental composition. Such an attempt has been made in based on the QGS interaction model and rigidity-dependent steepening primary energy spectra for a description of vertical MSU, AKENO and Tien-Shan EAS size spectra. In the present work we worked out a formalism of the inverse problem solution - reconstruction of the primary energy spectrum and elemental composition based on the known EAS size spectra of KASCADE , AKENO , EAS-TOP and ANI measured at different zenith angles. The calculations were done in the frames of QGSJET and SIBYLL interaction models. As a primary spectrum we have tested the modified rigidity-dependent steepening primary energy spectra and the hypothesis of the additional component in the knee region . In this case the type of nucleus of the additional component was considering as unknown and determining by the best fit of the fine structure of EAS size spectra in the knee region.
## 2 EAS inverse problem
In general, the energy spectra ($`\mathrm{}_A/E_0`$) of primary nuclei ($`A`$) and detectable EAS size spectra ($`I/N_e^{}`$) are related by the integral equation -
$$\frac{I(E_e,\overline{\theta },t)}{N_e^{}}=\underset{A}{}_{E_{min}}^{\mathrm{}}\frac{\mathrm{}_A}{E_0}W_\theta (E_0,A,N_e^{},\overline{\theta },t)𝑑E_0$$
(1)
where $`E_0,A,\overline{\theta }`$ are energy, nucleon number (1-59) and average zenith angle of primary nuclei, $`E_e`$ is an energy threshold of detected EAS electrons, $`N_e^{}(E>E_e)`$ is the estimation value of EAS size obtained by the electron lateral distribution function. Here, by the EAS size ($`N_e(E>0)`$) we mean the total number of EAS electrons at given observation level $`(t)`$. The kern ($`W_\theta `$) of integral equation (1) is determined as
$$W_\theta \frac{1}{\mathrm{\Delta }_\theta }_{\theta _1}^{\theta _2}_0^{\mathrm{}}\frac{\mathrm{\Omega }(E_0,A,\theta ,t)}{N_e}P_\theta \frac{\mathrm{\Psi }(N_e)}{N_e^{}}\mathrm{sin}\theta d\theta dN_e$$
where $`\mathrm{\Omega }/N_e`$ is an EAS size spectrum at the observation level $`(t)`$ for given $`E_0,A,\theta `$ parameters of a primary nucleus and depends on $`AA_{Air}`$ interaction model; $`\mathrm{\Delta }_\theta =\mathrm{cos}\theta _1\mathrm{cos}\theta _2`$;
$$P_\theta P(Ne,E_0,A,\theta )=\frac{1}{XY}D(N_e,E_0,A,\theta ,x,y)𝑑x𝑑y$$
is a probability to detect an EAS by scintillation array at EAS core coordinates $`|x|<X/2`$, $`|y|<Y/2`$ and to obtain the estimations of EAS parameters ($`N_e^{}`$, $`s`$ \- shower age, $`x^{},y^{}`$ \- shower core location) with given accuracies; $`\mathrm{\Psi }/N_e^{}`$ is a distribution of $`N_e^{}(N_e,s,x,y)`$ for given EAS size ($`N_e`$).
One may achieve significant simplification of equation (1) providing the following conditions during experiments:
a) selection of EAS cores in a range where $`P_\theta 1`$,
b) the log-Gaussian form of the measuring error ($`\mathrm{\Psi }/N_e^{}`$) with an average value $`\mathrm{ln}(N_e\delta )`$ and a RMSD $`\sigma _N`$, where $`\delta `$ involves all transfer factors (an energy threshold of detected EAS electrons, $`\gamma `$ and $`\mu `$ contributions) and slightly depends on $`E_0`$ and $`A`$,
c) transformation (standardization) of the measured EAS size spectra to the EAS size spectra at observation level
$$\frac{I(0,\overline{\theta },t)}{N_e}\eta \frac{I(E_e,\overline{\theta },t)}{N_e^{}},$$
where $`\eta =\delta ^{(\gamma _e1)}\mathrm{exp}\{(\gamma _e1)^2\sigma _N^2/2\}`$ and $`\gamma _e`$ is the EAS size power index,
d) consideration of either all-particle primary energy spectrum $`\mathrm{}_\mathrm{\Sigma }/E_0`$ with effective nucleus $`A_{eff}(E_0)`$ or energy spectra of primary nuclei ($`\mathrm{}_{A_\xi }/E_0`$) gathered in a limited number of groups ($`\xi =1,\mathrm{}\xi _{\mathrm{max}}`$) as unknown functions.
(a-d) conditions make EAS data of different experiments more comparable and equation (1) converts to the form
$$\frac{I(0,\overline{\theta },t)}{N_e}=\eta _{E_{min}}^{\mathrm{}}\frac{\mathrm{}_\mathrm{\Sigma }}{E_0}\frac{\mathrm{\Omega }(E_0,A_{eff}(E_0),\overline{\theta },t)}{N_e}𝑑E_0$$
(2)
or
$$\frac{I(0,\overline{\theta },t)}{N_e}=\eta \underset{\xi =1}{\overset{\xi _{\mathrm{max}}}{}}_{E_{min}}^{\mathrm{}}\frac{\mathrm{}_{A_\xi }}{E_0}\frac{\mathrm{\Omega }(E_0,\overline{A}_\xi ,\overline{\theta },t)}{N_e}𝑑E_0$$
(3)
However, even in this form the determination of primary energy spectra by measured EAS size spectra and solution of integral equations (2,3) in general is unsolvable problem. At the same time, using the a priori information about energy spectra of primary nuclei ($`\mathrm{}_{A_\xi }/E_0`$) and the EAS size spectra $`I/N_e^{}f_{i,j}(N_{e,i}^{},\overline{\theta }_j,t)`$ measured in $`i=1,\mathrm{}m`$ size intervals and $`j=1,\mathrm{}n`$ zenith angular intervals, one may transform the inverse problem into $`\chi ^2`$-minimization problem
$$\mathrm{min}\{\chi ^2\}\mathrm{min}\left\{\underset{i}{\overset{m}{}}\underset{j}{\overset{n}{}}\frac{(f_{i,j}F_{i,j})^2}{\sigma _f^2+\sigma _F^2}\right\}$$
(4)
with unknown (free) spectral parameters. Here $`F_{i,j}F(N_{e,i}^{},\overline{\theta }_j,t)`$ are the expected EAS size spectra determined at the right hands of equations (1-3) and $`\sigma _f,\sigma _F`$ are the uncertainties (RMSD) of measured ($`f_{i,j}`$) and expected ($`F_{i,j}`$) shower size spectra.
One may also unify the data of different experiments applying minimization $`\chi _U^2`$ with re-normalized EAS size spectra
$$\mathrm{min}\{\chi _U^2\}\mathrm{min}\left\{\chi ^2(\frac{f_{i,j,k}}{_i_jf_{i,j,k}},\frac{F_{i,j,k}}{_i_jF_{i,j,k}})\right\}$$
(5)
where index $`k=1,\mathrm{}l`$ determines the observation levels ($`t`$) of experiments. Expression (5) offers an advantage for experiments where the values of methodical shift ($`\delta `$) and measuring error ($`\sigma _N`$) are unknown or known with insufficient accuracy.
The energy spectra of primary nuclei are preferable to determine (a priori) in the following generalized form
$$\frac{\mathrm{}_A}{E_0}\beta \mathrm{\Phi }_AE_0^{\gamma _1(A)}\left(1+\left(\frac{E_0}{E_{knee}(A)}\right)^ϵ\right)^{(\gamma _1(A)\gamma _2)/ϵ}$$
(6)
Unknown (free) spectral parameters in approximation (6) are $`\beta `$, $`E_{knee}(A)`$ (so called ”knee” of energy spectrum of $`A`$ nucleus), $`\gamma _1`$ and $`\gamma _2`$ (spectral asymptotic slopes before and after knee), $`ϵ`$ (sharpness parameter of knee, $`1ϵ10`$). The values of $`\mathrm{\Phi }_A`$ and $`\gamma _1(A)`$ parameters are known from approximations of balloon and satellite data at $`A1,4,\mathrm{}59`$ and $`E_0110^3`$ TeV. Parameter $`\beta 1`$ determines the normalization of spectra (6) in $`10^210^5`$ TeV energy range.
Thus, minimizing $`\chi ^2`$\- functions (4,5) on the basis of measured values of $`I(\overline{\theta }_{i,k})/N_{ej,k}`$ and corresponding expected EAS size spectra (2,3) at given $`m`$ zenith angular intervals, $`n`$ EAS size intervals and $`l`$ experiments one may evaluate the parameters of the primary spectrum (6). Evidently, the accuracies of solutions for spectral parameters strongly depend on the number of measured intervals ($`mnl`$), statistical errors and correctness of $`\mathrm{\Omega }(E_0,A,\theta ,t)/N_e`$ determination in the framework of a given interaction model.
## 3 Results
Here, the parametric solutions of the EAS inverse problem are obtained on the basis of KASCADE ($`t=`$1020 g/cm<sup>2</sup>), AKENO (910 g/cm<sup>2</sup>), EAS-TOP (810 g/cm<sup>2</sup>) and ANI (700 g/cm<sup>2</sup>) published EAS size spectra. These experiments were carried out at different observation levels and were chosen for two reasons: satisfaction of (a-c) conditions from the section (2) and a high statistical accuracy of presented data (especially KASCADE experiment). Unfortunately, during the standardization of EAS size spectra (condition (c)) the value of $`\eta `$ parameter is not always known with proper accuracy, which is the main reason of discrepancy in the results of different experiments. In our calculations we included $`\eta `$ in the list of unknown spectral parameters and determined by the minimization of functional (4). The problem does not exist if the minimization of re-normalized EAS size spectra is unified, since the linear parameters ($`\eta _k\beta `$) are canceled out from the functional (5).
The differential EAS size spectra $`\mathrm{\Omega }(E_0,A,\theta ,t)/N_e`$ for given $`E_00.032,0.1,\mathrm{},100`$ PeV, $`A1,4,12,16,28,56`$, $`t0.5,0.6,\mathrm{},1`$ Kg/cm<sup>2</sup>, $`\mathrm{cos}\theta 0.8,0.9,1`$ were calculated using CORSIKA562(NKG) EAS simulation code at QGSJET and SIBYLL interaction models. Intermediate values are calculated using 4-dimensional log-linear interpolations. The estimations of errors of the expected EAS size spectra $`\mathrm{\Omega }/N_e`$ at fixed $`E_0,A,\theta ,t`$ parameters did not exceed $`35\%`$.
The basic results of minimizations (4,5) at a given number ($`\nu `$) of unknown spectral parameters and the values of $`\chi ^2/q`$ (or $`\chi _U^2/q_u`$ for unified data), are presented in Tables 1-4, where $`q=mn\nu 1`$ and $`q_u=_k(mn)_k\nu 1`$ are corresponding degrees of freedom. The upper (lower) rows of each experiment in Tables 1 and 4 correspond to parameters obtained by the QGSJET (SIBYLL) interaction model.
### 3.1 Test of rigidity-dependent energy spectra
Table 1 contains the approximation values of spectral parameters at rigidity-dependent approach
$$E_{knee}(A)=RZ$$
(7)
where $`Z`$ is a charge of $`A`$ nucleus and R is a parameter of magnetic rigidity (or a critical (cutoff) energy in more modern model ). The results were obtained by minimizations (4,5) applying (6,7). The magnitudes of $`\gamma _1(A)`$ for all nuclei were taken from . The number of unknown (free) parameters is equal to $`\nu =4`$ and corresponding solutions at two interaction models are presented in Table 1.
The stability of solutions (or the steep of $`\chi ^2`$ global minimum) of minimizations (4,5) is seen from obtained errors. So, it is seen that the $`\chi ^2`$ minimum for AKENO data does not depend on the $`\epsilon `$ sharpness parameter at SIBYLL and partly at QGSJET interaction models.
The obtained slopes of primary energy spectra after knee agree with the same calculations performed by QGS model and exceed well known expected values ($`33.1`$) in the $`10^{17}`$ eV energy range . Such a steep of primary spectra after the knee is a result of $`(\gamma _2,R)`$ correlations in (4-7). In case of $`R2000`$ TeV the $`E_k(Fe)5.210^4`$ TeV and the primary energy spectra in the large interval ($`E_k(H)E_k(Fe)`$) at fixed $`\gamma _1(A_\xi )`$ can be conformed with corresponding EAS size spectra provided abnormally steep slope ($`\gamma _23.43.5`$) after the knee.
The results of expected EAS size spectra in comparison with corresponding experimental data are shown in Fig. 1 (the thin solid lines by QGSJET model, the thin dashed lines by SIBYLL model). Despite the satisfactory agreement ($`\chi ^21`$) of EAS size spectra with predictions of rigidity-dependent steepening spectra (6,7) and QGSJET interaction model the form (fine structure) of the measured EAS size spectra in the knee region differs from the form of corresponding expected spectra. It is worth mentioning that the difference is formally small and does not exceed several percents in precise KASCADE data.
As a next step, we attempted to test the rigidity-dependent approximation (7) directly. The knee-energies $`E_k(A_\xi )`$ were chosen as unknown (free) parameters in the approximations of EAS data (4-7), where $`\xi =1,\mathrm{}\xi _{\mathrm{max}}`$. However, the stable solutions for free parameters were obtained only at unified minimization (5), fixed values $`\gamma _1(A_\xi ),\gamma _2,\epsilon `$ and number of nuclear groups $`\xi _{\mathrm{max}}5`$. In Table 2 the values of spectral parameters $`E_{knee}(A_\xi )`$ obtained by minimization of $`\chi _U^2`$ (5) for 5 groups of primary nuclei ($`H`$), ($`He,Li`$), ($`BeNa`$), ($`MgCl`$), ($`ArNi`$) at given values of $`\gamma _2=3.1`$ , $`ϵ=4.0`$ and $`\nu =5`$ are presented. It is seen, that approach (7) is performed only for nuclei with $`A>1`$ and $`R400`$ TV or (7) is valid for all nuclei at $`R400`$ TV but there is an additional proton flux with $`E_{knee}^{(p)}3000`$ TeV which shifts the knee value of the total proton energy spectrum.
The following testing of the rigidity-dependent approach (7) is based on the investigation of the all-particle energy spectrum. Toward this end the fits of primary energy spectra $`d\mathrm{}_A/dE_0`$ for different nuclei ($`A=1,\mathrm{}59`$) known from were extrapolated up to $`E_A=10^5`$ TeV energies taking into account (6,7) at $`R600`$ TeV . The obtained expected all-particle energy spectrum $`dI_\mathrm{\Sigma }/dE_0=_Ad\mathrm{}_A/dE_0`$ was approximated by expression similar to (6) at five free parameters ($`\mathrm{\Phi }_\mathrm{\Sigma },\gamma _1,\gamma _2,E_k,\epsilon `$). The average values of primary nuclei ($`\mathrm{exp}(\overline{\mathrm{ln}A})`$) obtained from extrapolations of spectra were approximated by step function
$$A_{eff}(E_0)=a+b\mathrm{ln}\left(\frac{E_0}{E_k}\right)$$
(8)
where $`b=b_1`$ at $`E_0<E_k`$ and $`b=b_2`$ at $`E_0>E_k`$. The results of these approximations are presented in a last row of Table 3 . The value of sharpness parameter was equal to $`\epsilon =1\pm 0.1`$.
The first and second rows of Table 3 content the parameters of the all-particle energy spectra ($`\mathrm{}_\mathrm{\Sigma }/E_0`$), which were obtained by minimization $`\chi _U^2`$ (expressions 2,5-7) of unified EAS size data at $`\nu =6`$ and $`\epsilon =1`$. The approximation (8) has been used for $`A_{eff}(E_0)`$ which is at the right hand of expression (2). It is necessary to note that the solutions for $`a,b_1,b_2`$ parameters can be obtained only by re-normalized EAS size spectra (5) because of the strong correlation between a linear parameter $`\eta `$ and an effective nucleus $`A_{eff}`$.
It is seen from Table 3 that the model of rigidity-dependent energy spectra predicts the increase ($`b_1>0`$, third row of Table 3) of the mean nucleus with energy, whereas the presented analysis of EAS data points out a decrease of $`A_{eff}`$ with energy ($`b_1<0`$, first two rows of Table 3) in the energy range of $`E_0<E_k`$. It is obvious, that despite $`\overline{A}`$ could not be exactly equal to $`A_{eff}`$, at least their dependence on energy must be the same.
The results of recent precise experiments DICE and CASA-BLANKA also point out to the decrease of $`\overline{\mathrm{ln}A}`$ with energy at $`E_0<E_k`$. This dependence of $`A_{eff}`$ and $`\overline{\mathrm{ln}A}`$ on energy might be explained by the contribution of an additional light component in a primary nuclei flux.
From the above analyses follows that rigidity-dependent steepening energy spectra in combination with QGSJET or SIBYLL interaction model can not explain the obtained results of the fine structure of EAS size spectra in the knee region (Table 1 and Fig. 1), the large values of knee $`E_k(H)`$ for Hydrogen component (Table 2) and dependence of the effective nucleus $`A_{eff}(E_0)`$ on primary energy before the knee (Table 3). In this connection we have carried on the search of more adequate model of primary energy spectra and elemental composition.
### 3.2 Test of additional component
Based on predictions the primary energy spectra in approximation (6) have been added by a new (polar cap ) component $`\mathrm{}_{Add}/E_0`$ with power energy spectrum
$$\frac{\mathrm{}_{Add}}{E_0}=\mathrm{\Phi }^{(p)}\left(E_k^{(p)}\right)^{\gamma _1^{(p)}}\left(\frac{E_0}{E_k^{(p)}}\right)^{\gamma ^{(p)}}$$
(9)
where $`\gamma ^{(p)}=\gamma _1^{(p)}`$ at $`E_0<E_k^{(p)}`$ and $`\gamma ^{(p)}=\gamma _2`$ at $`E_0>E_k^{(p)}`$.
New spectral parameters $`\mathrm{\Phi }^{(p)},\gamma ^{(p)},E_k^{(p)}`$ and nucleon number $`A^{(p)}`$ of the additional component are considered as unknown and determined together with parameters of rigidity-dependent energy spectra (6,7) by minimization of $`\chi ^2`$ and $`\chi _U^2`$ (4,5) at $`\nu =7`$. The results of expected EAS size spectra for each experiment (KASCADE, AKENO, EAS-TOP and ANI) taking into account contribution of additional component (9) are shown in Fig. 1 (the thick solid lines by QGSJET model and the thick dashed lines by SIBYLL model). It is seen that the additional component with high accuracy (2-5% for KASCADE and ANI data) describes the fine structure of EAS size spectra in the knee region. The values of slopes of additional component before the knee turned out to be $`\gamma _1^{(p)}1.5\pm _{0.2}^{0.6}`$. The nucleon number ($`A^{(p)}`$) of this component with high reliability did not exceed of $`A^{(p)}=1`$ for most of experiments especially at QGSJET interaction model (except from AKENO ($`A^{(p)}56`$)). The unified analyses of all experiments at QGSJET and SIBYLL interaction models also gave a proton or neutron ($`A^{(p)}=1`$) composition of the additional component. The values of other spectral parameters at $`\gamma _1^{(p)}=1.5`$ and $`A^{(p)}=1`$ are included in Table 4.
The obtained result disagrees with where the alike component consists of several heavy nuclei. However, our result is based on the high accuracy of the coincidence ($`23\%`$ for KASCADE data) of expected and measured EAS size spectra by both $`\chi ^2`$ criterion and the overlapping of fine structure of spectra.
The comparison of parameters $`\gamma _2`$ from Tables 1,4 shows that spectral slopes after the knee from Table 4 roughly overlap with the expected slope ($`33.1`$) well known from $`N_e10^7`$ EAS data. This can be explained the fact that the spectral brake of summary proton component ($`E_k^{(p)}210^3`$ TeV) is closer to the iron component ($`E_k(Fe)10^4`$ TeV) at $`R400`$ TV.
The coexistence of rigidity-dependent primary energy spectra and additional proton flux with spectral parameters $`E_k^{(p)}2000`$ TeV, $`\gamma _1^{(p)}1.5`$, $`\gamma _2^{(p)}=\gamma _23.1`$ explains also the shifted value of the knee $`E_k(H)`$ in Table 2 and the decrease of $`A_{eff}(E_0)`$ at $`E_0<E_k`$ (see section 3.1).
Although the additional component does not contribute to shower sizes below or above the knee significantly , it is essential for the reproducing of the sharp knee of EAS size spectra. This idea is taken from and is confirmed here.
The final all-particle energy spectrum ($`I/E_0`$) obtained by unified EAS size data at QGSJET interaction model
$$\frac{I}{E_0}=\beta \left(\underset{A}{}\frac{\mathrm{}_A}{E_0}+\frac{\mathrm{}_{Add}}{E_0}\right)$$
and corresponding energy spectra ($`\mathrm{}_A/E_0`$) of 6 nuclear groups with additional component ($`\mathrm{}_{Add}/E_0`$) at normalization $`\beta =1`$ are presented in Fig. 2. The solid (dashed) line is the all-particle energy spectrum obtained by unified (only KASCADE) EAS size spectra. Dotted lines are the energy spectra of different primary components obtained by unified EAS data. Symbols in Fig. 2 are the data from DICE , CASA-BLANKA arrays and reviews .
The numerical values of all-particle energy spectrum (solid line in Fig. 2), corresponding error, spectrum of additional proton component and average nucleus versus primary energy are presented in Table 5.
## Conclusion
High statistical accuracy of experiments KASCADE, EAS-TOP, AKENO and ANI allowed to obtain approximations of primary energy spectra and elemental composition with accuracy $`15\%`$ in the knee region. Along with this, KASCADE and ANI EAS size spectra at $`\theta <37^0`$ are described with the accuracy of $`25\%`$ in a whole measurement interval.
Obtained results show the evidence of QGSJET interaction model at least in $`10^510^7`$ TeV energy range, rigidity-dependent steepening primary energy spectra at $`R200400`$ TV and existence of the additional proton (or neutron) component with spectral power index $`\gamma _1^{(p)}1.5\pm _{0.2}^{0.6}`$ before the knee $`E_k^{(p)}2030`$ TeV. The contribution of the additional proton (neutron) component in all-particle energy spectrum turned out to be $`20\pm 5\%`$ at primary energy $`E_0=E_k^{(p)}`$.
## Acknowledgements
We thank Peter Biermann for extensive discussions and Heinigerd Rebel, Johannes Knapp and Dieter Heck for providing the CORSIKA code.
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# A two-box-shift morphism between Specht modules
## 0 Introduction
### 0.1 Known results
The irreducible modules of the complex group algebra of the symmetric group on $`n`$ letters possess the Specht lattices $`S^\lambda `$ as combinatorially described integral models, indexed by partitions $`\lambda `$ of $`n`$. Let $`d1`$ be an integer. Suppose given a partition $`\lambda `$ of $`n`$ and integers $`1gk\lambda _11`$ such that the shift of $`d`$ boxes situated at the bottom of column $`k+1`$ of $`\lambda `$ downwards to column $`g`$ yields a partition again, which we denote by $`\mu `$. Let $`\lambda ^{}`$ be the partition transposed to $`\lambda `$. The number
$$m_0:=\lambda _g^{}\lambda _{k+1}^{}+(k+1)g+d$$
of steps by which we shift leftwards and downwards is called the box shift length.
For instance, we may shift $`d=2`$ boxes of $`\lambda =(5,4,4,2,1,1)`$ from column $`k+1=4`$ downwards to column $`g=2`$ to obtain the partition $`\mu =(5,3,3,2,2,2)`$. We have $`\lambda ^{}=(6,4,3,3,1)`$, hence the box shift length equals $`m_0=5`$.
Theorem. Based on \[CL 74\], Carter and Payne \[CP 80\] have shown that
$$\text{Hom}_{K𝒮_n}(K_\text{Z}S^\lambda ,K_\text{Z}S^\mu )0,$$
$`K`$ being an infinite field of characteristic $`p`$ such that $`d<p^{v_p(m_0)}`$, i.e. such that the number of shifted boxes is smaller than the $`p`$-part of the box shift length. In case $`d=1`$, this condition on $`p`$ translates to $`p|m_0`$. In case $`d=2`$, it translates for $`p3`$ to $`p|m_0`$, and for $`p=2`$ to $`4|m_0`$.
Theorem. Suppose $`d=1`$. In \[K 99, 4.3.31, cf. 0.7.1\], the existence of an element of order $`m_0`$ in
$$\text{Hom}_{(\text{Z}/(m_0))𝒮_n}(\text{Z}/(m_0)_\text{Z}S^\lambda ,\text{Z}/(m_0)_\text{Z}S^\mu )$$
has been shown via construction of such a morphism.
### 0.2 Results
We maintain the situation from Section 0.1.
Suppose $`d=2`$. Let $`m`$ be either $`m_0`$ or $`m_0/2`$, the latter value being taken in case $`m_0`$ is even, a similar combinatorial integer is even and the diagram of $`\lambda `$ takes a certain shape near the column $`g`$.
Theorem (2.37). There exists a morphism of order $`m`$ in
$$\text{Hom}_{(\text{Z}/(m))𝒮_n}(\text{Z}/(m)_\text{Z}S^\lambda ,\text{Z}/(m)_\text{Z}S^\mu ).$$
As a consequence, we recover the according part of the result of Carter and Payne \[CP 80\]. An account of the construction of this morphism, including the definition of the modulus $`m`$, is given at the end of Section 1.1. In Section 0.3, we attempt to give an informal description of that construction. Some examples may be found in Section 2.4.
The proof by means of our techniques requires consideration of (essentially) $`9^2`$ cases - the base $`9`$ resulting from the structure of the one-step Garnir relations, the exponent $`2`$ being the number of shifted boxes (cf. A.1). Moreover, the connection between the modulus and the box shift length results from the calculation only.
A morphism modulo the box shift length $`m_0`$ itself need not exist. For instance, we obtain
$$\text{Hom}_{\text{Z}𝒮_5}(S^{(2,2,1)},\text{Z}/(5!)_\text{Z}S^{(1,1,1,1,1)})\text{}\text{Hom}_{\text{Z}𝒮_5}(S^{(2,2,1)},\text{Z}/(2)_\text{Z}S^{(1,1,1,1,1)})\text{Z}/2.$$
Let $`d`$ be arbitrary, but suppose $`g=k`$. This is to say, we shift $`d`$ boxes from column $`g+1`$ to column $`g`$.
Theorem (4.9). Let $`\text{ilog}_pd:=\mathrm{max}\{i\text{Z}_0|p^id\}`$. There exists a morphism of order $`m:=m_0{\displaystyle \underset{p\text{ prime, }p|m_0}{}}p^{\mathrm{min}(v_p(m_0),\text{ilog}_p(d))}`$ in
$$\text{Hom}_{(\text{Z}/(m))𝒮_n}(\text{Z}/(m)_\text{Z}S^\lambda ,\text{Z}/(m)_\text{Z}S^\mu ).$$
Again, the according part of the result of Carter and Payne \[CP 80\] is recovered. Roughly speaking, this morphism is given by picking $`d`$ entries from column $`g+1`$ of a given $`\lambda `$-polytabloid, by placing them at the bottom of column $`g`$, and by mapping this $`\lambda `$-polytabloid to an alternating sum of the resulting $`\mu `$-polytabloids over the various possibilities to pick those entries.
Remark. For an abelian group $`A`$ and an integer $`m2`$, we write $`A[m]:=\text{Kern}(A\text{}A)`$. There is an isomorphism
$$\text{Hom}_{(\text{Z}/(m))𝒮_n}(\text{Z}/(m)_\text{Z}S^\lambda ,\text{Z}/(m)_\text{Z}S^\mu )\text{}\text{Ext}_{\text{Z}𝒮_n}^1(S^\lambda ,S^\mu )[m]$$
which translates both results on the existence of a morphism of order $`m`$ into assertions on the existence of an extension of order $`m`$ in $`\text{Ext}^1`$. Cf. Section 0.4.
Remark. Let $`m2`$ be an integer. For arbitrary partitions $`\lambda `$, $`\mu `$ of $`n`$, we dispose of the transposition isomorphism
$$\text{Hom}_{(\text{Z}/(m))𝒮_n}(\text{Z}/(m)_\text{Z}S^\lambda ,\text{Z}/(m)_\text{Z}S^\mu )\text{}\text{Hom}_{(\text{Z}/(m))𝒮_n}(\text{Z}/(m)_\text{Z}S^\mu ^{},\text{Z}/(m)_\text{Z}S^\lambda ^{}),$$
given by dualization, followed by alternation and isomorphic substitution (3.16). Based on the simple assertion (3.2), we give in (3.27) an explicit, but nevertheless not quite satisfactory formula for the transpose of the two-box-shift morphism in (2.37) in terms of tabloids. Whereas the image of a polytabloid under that transpose is known to be an element of the target Specht module, hence known to be expressible as a linear combination of polytabloids, a formula for such an expression is lacking. Some examples are given in Section 3.4, cf. also (4.15).
### 0.3 Construction (informal description)
The following description might serve as an intuitive thread through our formalism concerning the two-box-shift morphism in (2.37). We maintain the situation of the beginning of Section 0.2 and refer to the usual way to depict partitions and tableaux \[J 78\].
We write the Specht lattice $`S^\lambda `$ as a quotient of the free $`\text{Z}𝒮_n`$-module $`F^\lambda `$ on one generator, displayed as having as Z-linear basis the set of $`\lambda `$-tableaux equipped with the action of the $`𝒮_n`$ on the tableau entries. The kernel of the canonical epimorphism $`F^\lambda \text{}S^\lambda `$ is generated, over $`\text{Z}𝒮_n`$, by signed column transpositions and one-step Garnir relations. Given a tableau, a signed column transposition is the sum of the tableau and the tableau having two entries of a column interchanged. A one-step Garnir relation depends on subsets of two subsequent columns of a tableau $`[a]`$ the sizes of which add up to (the length of the left column)$`+1`$. After factoring out the signed column transpositions, such a Garnir relation can be written, up to a scalar, as the alternating sum of the tableaux arising from $`[a]`$ by permutation inside the union of the chosen subsets of the columns, the sign being that of the afforded permutation. Our morphism is constructed as a factorization of a morphism $`(F^\lambda \text{}S^\mu \text{}S^\mu /mS^\mu )`$ over the canonical epimorphism $`(F^\lambda \text{}S^\lambda \text{}S^\lambda /mS^\lambda )`$. To a $`\lambda `$-tableau $`[a]`$, we may apply certain place operations, indexed by double paths, which produce $`\mu `$-polytabloids. A double path is a pair of sequences of positions in the union of the diagrams of $`\lambda `$ and $`\mu `$ that do not intersect, that run strictly from right to left, that start arbitrarily in column $`k+1`$ and that end at the the position of the shifted boxes in column $`g`$. (So at least $`4`$ and at most $`4+2(kg)`$ positions are occupied.) Given such a double path, the place operation it gives rise to is obtained by inserting entries of column $`k+1`$ into the double path from the right and by subsequently pushing entries through along each path separately. In particular, the last two entries forced to move are installed at the positions occupied by the shifted boxes in the diagram of $`\mu `$. Cf. (2.1). Each double path has a weight attached, i.e. a tuple recording the number of positions it occupies in each column, and a sign, depending on the positions it occupies in column $`k+1`$. Sending $`[a]`$ to the accordingly signed sum of the resulting $`\mu `$-polytabloids indexed by double paths of a fixed weight yields a $`\text{Z}𝒮_n`$-linear map from $`F^\lambda `$ to $`S^\mu `$ that annihilates the signed column transpositions. Forming a linear combination of these maps, indexed over the weights, equipped with certain coefficients which are polynomial in the combinatorial data, and dividing by an essentially polynomial factor of redundancy yields a $`\text{Z}𝒮_n`$-linear map $`F^\lambda \text{}S^\mu `$ that annihilates all one-step Garnir relations except for those involving column $`g`$ and column $`g+1`$, which vanish only modulo $`mS^\mu `$. Thus $`(F^\lambda \text{}S^\mu \text{}S^\mu /mS^\mu )`$ factors over a $`\text{Z}𝒮_n`$-linear map $`S^\lambda /mS^\lambda \text{}S^\mu /mS^\mu `$. Since the image of $`F^\lambda \text{}S^\mu `$ has not been contained in a strict multiple $`tS^\mu S^\mu `$, $`t2`$, the resulting morphism $`S^\lambda /mS^\lambda \text{}S^\mu /mS^\mu `$ is of order $`m`$ in its Hom-group.
### 0.4 Motivation
We consider the integral group ring $`\text{Z}𝒮_n`$ as a subring of a product of integral matrix rings via Wedderburn’s embedding, sending a group element $`g𝒮_n`$ to the tuple $`(\rho ^\lambda (g))_\lambda `$ of its operating matrices on the Specht lattices with respect to a choice of integral bases. Let $`m2`$ be an integer. A Z-linear map $`S^\lambda \text{}S^\mu `$ that becomes $`\text{Z}𝒮_n`$-linear modulo $`m2`$, i.e. that satisfies $`f\rho ^\lambda (g)_m\rho ^\mu (g)f`$ for all $`g𝒮_n`$, imposes this congruence as a necessary condition on an element of our product of integral matrix rings to lie in the image of that embedding. More generally, unscrewing the regular lattice into simple lattices via a binary tree of short exact sequences, we obtain a necessary and sufficient system of modular morphisms by means of the correspondence
$$\text{Hom}_{\text{Z}𝒮_n}(X,Y/mY)\text{}\text{Ext}_{\text{Z}𝒮_n}^1(X,Y)[m],$$
where $`X`$ and $`Y`$ are rationally disjoint $`\text{Z}𝒮_n`$-lattices. This correspondence attaches to a morphism $`X\text{}Y/mY`$ the extension obtained as pullback of
$$0\text{}Y\text{}Y\text{}Y/mY\text{}0.$$
Conversely, a retraction up to $`m`$ to the inclusion of an extension of $`X`$ by $`Y`$ yields the corresponding modular morphism as being induced on the cokernels. Since the $`\text{Ext}^1`$-order of a short exact sequence that occurs in that binary tree is not necessarily square-free, modular means modulo prime powers. See \[K 99\] for an elaboration on this theme.
### 0.5 Acknowledgements
I’d like to thank Prof. Rentschler for kind hospitality during the time of the preparation of this article. Large parts of the work have been supported by the EU TMR-network ‘Algebraic Lie Representations’, grant no. ERB FMRX-CT97-0100. I’d like to thank my fellow postdocs in Paris for help of various kinds. I’d like to thank Prof. Steffen König for encouragement, advice and support during and after the time of the preparation of my thesis, upon the methods of which the present result relies.
### 0.6 Leitfaden
## 1 Garnir
### 1.1 Specht lattices
For $`a,b\text{Z}`$, we let $`[a,b]:=\{i\text{Z}|ai\text{ and }ib\}`$.
Let $`n1`$. The symmetric group on the set $`[1,n]`$ is denoted by $`𝒮_n:=\text{Aut}_{\text{Sets}}[1,n]`$. The sign of a permutation $`\sigma 𝒮_n`$ is denoted by $`\epsilon _\sigma `$. Maps are written in various manners, on the right, on the left, using indices etc. The symmetric group acts on the right. Occasionally, we shall allow ourselves to treat tuples consisting of pairwise different entries and single elements as sets. Intervals are to be read as subsets of Z. Conjugation in a group is also written as $`h^g:=g^1hg`$. Given a module $`X`$ over a commutative ring $`R`$ and an element $`rR`$, we shall usually write $`X/r:=X/rX`$ (cf. e.g. 3.14). If $`X`$ and $`Y`$ are modules over an $`R`$-algebra $`A`$, we identify $`\text{Hom}_A(X/r,Y/r)=\text{Hom}_A(X,Y/r)`$ via composition with the residue class map $`X\text{}X/r`$.
Let
$$\begin{array}{ccc}\hfill \text{N}& \text{}& \text{N}_0\hfill \\ \hfill i& \text{}& \lambda _i\hfill \end{array}$$
be a partition of $`n`$, i.e. assume $`_i\lambda _i=n`$ and $`\lambda _i\lambda _{i+1}`$ for $`i\text{N}`$. Usually, we write a partition as a tuple, i.e. $`\lambda =(\lambda _1,\lambda _2,\mathrm{})`$. Sometimes, we abbreviate repeated entries by exponents that indicate their multiplicity, such as e.g. $`(9,2,2,2,1,1)=(9,2^3,1^2)(9,8,1)`$. Let $`[\lambda ]:=\{i\times j\text{N}\times \text{N}|\lambda _ij\}`$ be the diagram of $`\lambda `$. A $`\lambda `$-tableau is a bijection
$$\begin{array}{ccc}\hfill [\lambda ]& \text{}& [1,n]\hfill \\ \hfill i\times j& \text{}& a_{j,i}\hfill \end{array}.$$
The transposed partition of $`\lambda `$ is denoted by $`\lambda ^{}`$, i.e. $`i\times j[\lambda ]j\times i[\lambda ^{}]`$. For a $`\lambda `$-tableau $`[a]`$, we denote by $`[a^{}]`$ the $`\lambda ^{}`$-tableau obtained by composition with $`i\times j\text{}j\times i`$, mapping $`[\lambda ^{}]`$ onto $`[\lambda ]`$. A permutation $`\sigma 𝒮_n`$ acts on the set of $`\lambda `$-tableaux $`T^\lambda `$ via composition $`[a]\text{}[a]\sigma `$. Let $`F^\lambda `$ be the free Z-module on $`T^\lambda `$, endowed with the induced action of the $`𝒮_n`$. Let
$$\begin{array}{ccccccc}\hfill [\lambda ]& \text{}& \text{N}\hfill & ,& \hfill [\lambda ]& \text{}& \text{N}\hfill \\ \hfill i\times j& \text{}& i\hfill & & \hfill i\times j& \text{}& j\hfill \end{array}$$
denote the projections. To a $`\lambda `$-tableau $`[a]`$ we attach a $`\lambda `$-tabloid
$$\{a\}:=\{[a]\}:=[a]^1\pi _R^\lambda \text{N}^{[1,n]}.$$
In general, $`X`$ being a set, an element $`\sigma 𝒮_n`$ acts on $`fX^{[1,n]}`$ via $`(i)(f\sigma ):=(i\sigma ^1)f`$. In particular, the free Z-module on the set of $`\lambda `$-tabloids, denoted by $`M^\lambda `$, carries a structure as a $`\text{Z}𝒮_n`$-lattice, viz. $`\{a\}\sigma =\{[a]\sigma \}`$. Let
$$\begin{array}{ccc}\hfill C_{[a]}& :=& \{\kappa 𝒮_n|[a]^1\pi _C^\lambda =([a]\kappa )^1\pi _C^\lambda \}\hfill \\ \hfill R_{[a]}& :=& \{\rho 𝒮_n|[a]^1\pi _R^\lambda =([a]\rho )^1\pi _R^\lambda \}\hfill \end{array}$$
be the column stabilizer and the row stabilizer of $`[a]`$, respectively. Note that $`C_{[a]\sigma }=(C_{[a]})^\sigma `$ and $`R_{[a]\sigma }=(R_{[a]})^\sigma `$. Let the Specht lattice $`S^\lambda `$ be the $`\text{Z}𝒮_n`$-sublattice of $`M^\lambda `$ generated Z-linearly by the $`\lambda `$-polytabloids
$$a:=[a]:=\underset{\kappa C_{[a]}}{}\{a\}\kappa \epsilon _\kappa M^\lambda .$$
Let $`\zeta [1,n]`$. We denote by $`𝒮_\zeta :=C_{𝒮_n}([1,n]\backslash \zeta )`$ the subgroup of $`𝒮_n`$ consisting of permutations that act merely on $`\zeta `$. For an element $`x`$ of a $`\text{Z}𝒮_n`$-module $`X`$ we denote
$$x\zeta :=\underset{\sigma S_\zeta }{}x\sigma \epsilon _\sigma .$$
Given $`p1`$ and $`1jk\lambda _p^{}`$, we write $`a_{p,[j,k]}:=(a_{p,i})_{i[j,k]}`$ and abbreviate $`a_p:=a_{p,[1,\lambda _p^{}]}`$. Let $`\xi a_p`$, $`\eta a_{p+1}`$ be such that $`\mathrm{\#}\xi +\mathrm{\#}\eta =\lambda _p^{}+1`$. Letting $`𝒮_\xi \times 𝒮_\eta \backslash 𝒮_{\xi \eta }`$ denote a chosen set of representatives of right cosets, the expression
$$G_{[a],\xi ,\eta }^{\prime \prime }:=\underset{\sigma 𝒮_\xi \times 𝒮_\eta \backslash 𝒮_{\xi \eta }}{}[a]\sigma \epsilon _\sigma $$
is called a one-step Garnir relation. For $`p1`$, $`u,va_p`$, $`uv`$, the expression
$$[a]+[a](u,v)$$
is called a signed column transposition. From \[J 78, 7.2\] and from the proof of \[J 78, 8.4\] we take that the (finite) set
$$\begin{array}{cc}& \{[a]+[a](u,v)|[a]T^\lambda ,p1,u,va_p,uv\}\hfill \\ \hfill & \{G_{[a],\xi ,\eta }^{\prime \prime }|[a]T^\lambda ,p1,s,t1,s+t=\lambda _p^{}+1,\xi =a_{p,[\lambda _p^{}s+1,\lambda _p^{}]},\eta =a_{p+1,[1,t]}\}\hfill \end{array}$$
$`\text{Z}𝒮_n`$-linearly generates the kernel of the epimorphism
$$\begin{array}{ccc}\hfill F^\lambda & \text{}& S^\lambda \hfill \\ \hfill [a]& \text{}& a.\hfill \end{array}$$
Given a $`\text{Z}𝒮_n`$-module $`X`$, its alternated module is written $`X^{}:=X_\text{Z}\text{Z}^{}`$, $`\text{Z}^{}`$ being the abelian group Z, equipped with $`1\sigma :=\epsilon _\sigma `$, $`\sigma 𝒮_n`$. So $`\text{Z}^{}S^{(1^n)}`$. For each partition $`\lambda `$ of $`n`$, we fix a $`\lambda `$-tableau $`[a_\lambda ]`$ and let
$$\{[a_\lambda ^{}]\sigma \}^{}:=(\{a_\lambda ^{}\}1)\sigma M^{\lambda ^{},}$$
for $`\sigma 𝒮_n`$. Note that $`\{a^{}\}^{}\sigma =\{[a^{}]\sigma \}^{}`$ and that $`\{a^{}\}^{}\kappa =\{a^{}\}\epsilon _\kappa `$ for $`[a]T^\lambda `$, $`\sigma 𝒮_n`$, $`\kappa C_{[a]}`$. We obtain a factorization into epimorphisms
$$\left(\begin{array}{ccc}\hfill F^\lambda & \text{}& S^\lambda \hfill \\ \hfill [a]& \text{}& a\hfill \end{array}\right)=\left(\begin{array}{ccccc}\hfill F^\lambda & \text{}& M^{\lambda ^{},}\hfill & \text{}& S^\lambda \hfill \\ \hfill [a]& \text{}& \{a^{}\}^{}\hfill & \text{}& a\hfill \end{array}\right)$$
by first factoring out the signed column transpositions. Using this factorization, we may write
$$G_{[a],\xi ,\eta }^{}:=(G_{[a],\xi ,\eta }^{\prime \prime })\nu _M^\lambda =\frac{1}{\mathrm{\#}\xi !\mathrm{\#}\eta !}\{a^{}\}^{}(\xi \eta ).$$
A $`\lambda `$-tableau $`[a]`$ is called standard if $`a_{j,i}a_{j^{},i^{}}`$ for $`i\times j,i^{}\times j^{}[\lambda ]`$ such that $`ii^{}`$ and $`jj^{}`$. The tuple $`(a|[a]\text{ is a standard }\lambda \text{-tableau})`$ is a Z-linear basis of $`S^\lambda `$ \[J, 8.4\]. Its elements are called standard $`\lambda `$-polytabloids. Let $`[\stackrel{ˇ}{a}_\lambda ]`$ be the standard $`\lambda `$-tableau determined by $`\stackrel{ˇ}{a}_{\lambda ,j,i+1}=\stackrel{ˇ}{a}_{\lambda ,j,i}+1`$ whenever $`i\times j,(i+1)\times j[\lambda ]`$. (We do not require a priori that $`[a_\lambda ]=[\stackrel{ˇ}{a}_\lambda ]`$, a requirement which is convenient for practical purposes, however.)
Account of the result (2.37). The notation deviates from our working notation further down. In particular, the following definitions are valid only in the remainder of this subsection (and will be repeated further down in case they coincide nonetheless).
Let $`1g<k\lambda _11`$ such that
$$\mu _i^{}:=\{\begin{array}{cc}\lambda _i^{}+2\hfill & \text{for }i=g\hfill \\ \lambda _i^{}2\hfill & \text{for }i=k+1\hfill \\ \lambda _i^{}\hfill & \text{else}\hfill \end{array}$$
defines a partition $`\mu `$. (For the case of $`g=k`$, see (2.17)).
A double path $`\gamma `$ is a pair of integers $`l(\delta )1`$ together with a pair of maps
$$\begin{array}{ccc}\hfill \gamma _\delta :[0,l\left(\delta \right)]& \text{}& \left[\lambda \right]\left[\mu \right]\hfill \\ \hfill i& \text{}& \alpha (\delta ,i)\times \beta (\delta ,i),\hfill \end{array}$$
where $`\delta [1,2]`$, subject to the following conditions.
* $`i<i^{}`$ implies $`\beta (\delta ,i)<\beta (\delta ,i^{})`$ for $`\delta [1,2]`$ and $`i,i^{}[0,l\left(\delta \right)]`$.
* $`\gamma _1\left([0,l\left(1\right)]\right)\gamma _2\left([0,l\left(2\right)]\right)=\mathrm{}`$.
* $`\alpha (\delta ,0)\times \beta (\delta ,0)=\left(\lambda _g^{}+\delta \right)\times g`$ for $`\delta [1,2]`$.
* $`\beta (\delta ,l\left(\delta \right))=k+1`$ for $`\delta [1,2]`$.
* $`\alpha (1,l\left(1\right))<\alpha (2,l\left(2\right))`$.
The set of double paths is denoted by $`\mathrm{\Gamma }`$. Suppose given $`\gamma \mathrm{\Gamma }`$. Let $`\epsilon _\gamma :=(1)^{\alpha (1,l(1))+\alpha (2,l(2))}`$. For $`j[g+1,k]`$, we let
$$e(\gamma ,j):=\underset{\delta [1,2]}{}\mathrm{\#}\gamma _\delta ^1\left([1,\lambda _j^{}]\times \left\{j\right\}\right).$$
Let
$$\begin{array}{ccc}\hfill [0,\lambda _{k+1}^{}]& \text{}& [0,\mu _{k+1}^{}]\hfill \\ \hfill i& \text{}& \mathrm{\#}\left([1,i]\backslash \{\alpha (1,l\left(1\right)),\alpha (2,l\left(2\right))\}\right)\hfill \\ \hfill \mathrm{min}\left(\phi ^1\left(\left\{j\right\}\right)\right)& \text{}& j.\hfill \end{array}$$
Given a $`\lambda `$-tableau $`[a]`$, we let the $`\mu `$-tableau $`[a^\gamma ]`$ be defined by
$$\begin{array}{cccc}\hfill a_{j,i}^\gamma & :=& a_{j,i}\hfill & \text{ for }i\times j\left[\mu \right]\backslash \left(\gamma _1\left([0,l\left(1\right)1]\right)\gamma _2\left([0,l\left(2\right)1]\right)\text{N}\times \left\{k+1\right\}\right)\hfill \\ \hfill a_{\beta (\delta ,i),\alpha (\delta ,i)}^\gamma & :=& a_{\beta (\delta ,i+1),\alpha (\delta ,i+1)}\hfill & \text{ for }\delta [1,2],i[0,l\left(\delta \right)1]\hfill \\ \hfill a_{k+1,i}^\gamma & :=& a_{k+1,\psi \left(i\right)}\hfill & \text{ for }i[1,\mu _{k+1}^{}].\hfill \end{array}$$
For $`j[g,k+1]`$, we let
$$X_j:=\left(\lambda _j^{}j\right)\left(\lambda _{k+1}^{}\left(k+1\right)\right).$$
For $`i[0,1]`$, we denote
$$L\left(i\right):=\left\{j[g+1,k1]\right|\lambda _{j+1}^{}=\lambda _j^{}i\}.$$
We write $`[g+1,k]=_{\kappa [1,K]}[p(\kappa ),q(\kappa )]`$ such that $`p(1)=g+1`$, such that $`p(\kappa )q(\kappa )`$, $`[p(\kappa ),q(\kappa )1]L(0)L(1)`$ and $`q(\kappa )L(0)L(1)`$ for $`\kappa [1,K]`$, such that $`q(\kappa )+1=p(\kappa +1)`$ for $`\kappa [1,K1]`$, and such that $`q(K)=k`$. For $`Z\text{Z}`$ and $`i0`$, we write $`Z^{(i)}:=(Z+i1)!/(Z1)!`$. We let
$$R:=\left(\underset{i[0,1]}{}\underset{jL\left(i\right)}{}X_j^{\left(2i\right)}\right)\left(\underset{\stackrel{\kappa [1,K],}{[p\left(\kappa \right),q\left(\kappa \right)1]L\left(1\right)}}{}\mathrm{gcd}(2,X_{q\left(\kappa \right)})\right)$$
and
$$m:=\{\begin{array}{cc}X_g+2\hfill & \text{ in case }[p\left(1\right),q\left(1\right)1]L\left(1\right)\hfill \\ \left(X_g+2\right)/\mathrm{gcd}(2,X_g,X_{g+1})\hfill & \text{ in case }[p\left(1\right),q\left(1\right)1]L\left(1\right)\hfill \end{array}.$$
The $`\text{Z}𝒮_n`$-linear map
$$\begin{array}{ccc}\hfill F^\lambda & \text{}& S^\mu \hfill \\ \hfill \left[a\right]& \text{}& \underset{\gamma \mathrm{\Gamma }}{}\left(\underset{j[g+1,k]}{}X_j^{\left(2e(\gamma ,j)\right)}\right)a^\gamma \epsilon _\gamma \hfill \end{array}$$
is divisible by $`R`$, and its quotient by $`R`$ factors $`\text{Z}𝒮_n`$-linearly as
$$\left(\begin{array}{ccccc}\hfill F^\lambda & \text{}& S^\mu \hfill & \text{}& S^\mu /m\hfill \\ & & b\hfill & \text{}& b+mS^\mu \hfill \end{array}\right)=\left(\begin{array}{ccccc}\hfill F^\lambda & \text{}& S^\lambda \hfill & \text{}& S^\mu /m\hfill \\ \hfill \left[a\right]& \text{}& a\hfill & & \end{array}\right).$$
The resulting morphism $`S^\lambda \text{}S^\mu /m`$ is of order $`m`$ as an element of $`\text{Hom}_{\text{Z}𝒮_n}(S^\lambda ,S^\mu /m)`$.
### 1.2 A Garnir formula
Let $`\lambda `$ be a partition of $`n`$, let $`[a]`$ be a $`\lambda `$-tableau, let $`1p<q\lambda _1`$, let $`d1`$. Let $`\xi a_p`$, $`s:=\mathrm{\#}\xi `$, $`\overline{\xi }:=a_p\backslash \xi `$ and $`\eta a_q`$, $`t:=\mathrm{\#}\eta `$, $`\overline{\eta }:=a_q\backslash \eta `$, $`u:=\mathrm{\#}\overline{\eta }`$, be such that $`s+t=\lambda _p^{}+1d`$. We choose a disjoint decomposition $`\overline{\xi }=\overline{\xi }_0\overline{\xi }_1`$ such that $`\mathrm{\#}\overline{\xi }_0=d1`$ and $`\mathrm{\#}\overline{\xi }_1=t`$.
###### Lemma 1.1 (cf. \[K 99, 4.3.4\])
Assume given $`x\phi a_p`$, $`y\psi a_q`$. Then
$$\begin{array}{ccc}\hfill a(\phi \psi )& =& \mathrm{\#}\phi a(\phi \psi \backslash x)\mathrm{\#}\psi a(x,y)(\phi \psi \backslash x)\hfill \\ & =& \mathrm{\#}\psi a(\phi \psi \backslash y)\mathrm{\#}\phi a(x,y)(\phi \psi \backslash y).\hfill \end{array}$$
In fact,
$$\begin{array}{ccc}\hfill a(\phi \psi )& =& \underset{z\phi \psi }{}\underset{\sigma 𝒮_{\phi \psi },z\sigma =x}{}a\sigma \epsilon _\sigma \hfill \\ & =& \left(\underset{z\phi \backslash x}{}\underset{\sigma 𝒮_{\phi \psi },z\sigma =x}{}a\sigma \epsilon _\sigma \right)\hfill \\ & +& \left(a(\phi \psi \backslash x)\right)\hfill \\ & +& \left(\underset{z\psi \backslash y}{}\underset{\sigma 𝒮_{\phi \psi },z\sigma =x}{}a\sigma \epsilon _\sigma \right)\hfill \\ & +& \left(\underset{\sigma 𝒮_{\phi \psi },y\sigma =x}{}a\sigma \epsilon _\sigma \right)\hfill \\ & =& \left((\mathrm{\#}\phi 1)a(\phi \psi \backslash x)\right)\hfill \\ & +& \left(a(\phi \psi \backslash x)\right)\hfill \\ & & \left((\mathrm{\#}\psi 1)a(x,y)(\phi \psi \backslash x)\right)\hfill \\ & & \left(a(x,y)(\phi \psi \backslash x)\right).\hfill \end{array}$$
The second equality follows by symmetry.
Suppose given disjoint subsets $`\phi ,\psi [1,n]`$ and a bijection $`\phi \text{}\psi `$. We denote
$$(\phi ,\psi )_\alpha :=\underset{x\phi }{}(x,(x)\alpha )𝒮_n.$$
In case $`\phi a_p`$, $`\psi a_q`$, $`\mathrm{\#}\phi =\mathrm{\#}\psi `$, the element $`a(\phi ,\psi ):=a(\phi ,\psi )_\alpha `$ is independent of the choice of a bijection $`\phi \text{}\psi `$.
###### Lemma 1.2 (\[K 99, 4.3.5\])
In case $`d=1`$ we have
$$a(\xi \eta )=s!t!a(\overline{\xi },\eta ).$$
We calculate
$$\begin{array}{ccc}\hfill a(\xi \eta )& \stackrel{\text{(}\text{1.1}\text{)}}{=}& \frac{t}{s+1}a(x,y)(\xi \eta )\hfill \\ & \stackrel{\text{(}\text{1.1}\text{)}}{=}& \frac{t(t1)}{(s+1)(s+2)}a(x,y)(x^{},y^{})(\xi \eta )\hfill \\ & \stackrel{\text{(}\text{1.1}\text{)}}{=}& \mathrm{}\hfill \end{array}$$
$$\begin{array}{ccc}& \stackrel{\text{(}\text{1.1}\text{)}}{=}& \frac{s!t!}{(s+t)!}a(\overline{\xi },\eta )(\xi \eta )\hfill \\ & =& s!t!a(\overline{\xi },\eta ),\hfill \end{array}$$
where $`\overline{\xi }=\{x,x^{},\mathrm{}\}`$, $`\eta =\{y,y^{},\mathrm{}\}`$.
###### Lemma 1.3
Let $`\phi a_p`$ and $`\psi a_q`$ be given such that $`\mathrm{\#}\phi =\mathrm{\#}\psi `$ and let $`\sigma 𝒮_{a_p}`$. We obtain
$$a(\phi ,\psi )\sigma =a((\phi )\sigma ,\psi )\epsilon _\sigma .$$
Factorization reduces us to the consideration of a transposition $`\sigma =(x,y)`$, where, moreover, we may assume $`x\phi `$ and $`ya_p\backslash \phi `$. Choosing $`z\psi `$, we conclude
$$\begin{array}{ccc}\hfill a(\phi ,\psi )(x,y)& =& a((\phi \backslash x),(\psi \backslash z))(x,z)(x,y)\hfill \\ & =& a((\phi \backslash x),(\psi \backslash z))(y,z)\hfill \\ & =& a((\phi )(x,y),\psi ).\hfill \end{array}$$
###### Lemma 1.4
The formula
$$a(\xi \eta )=\frac{s!}{(d1)!}a(\overline{\xi }_1,\eta )\overline{\xi }$$
holds.
The claimed equality may be reformulated to
$$a(\xi \eta )\overline{\xi }=\frac{s!(t+d1)!}{(s+t)!(d1)!}a(\overline{\xi }_1,\eta )(\xi \eta )\overline{\xi }.$$
We perform an induction over $`d`$, the case $`d=1`$ being treated in (1.2). We assume $`d2`$ and perform an induction over $`t`$, the case $`t=0`$ being trivial. We assume $`t1`$, choose $`x_0\overline{\xi }_0`$, $`x_1\overline{\xi }_1`$, $`y\eta `$ and denote $`\overline{\xi }_2:=(\overline{\xi }_1)(x_0,x_1)`$. Evaluating in two ways, we obtain
$$\begin{array}{cc}& a(\xi x_1\eta )(\overline{\xi }\backslash x_1)\hfill \\ \stackrel{\text{1., case }d1\text{}t}{=}& \frac{(s+1)!(t+d2)!}{(s+t+1)!(d2)!}a(\overline{\xi }_2,\eta )(\xi x_1\eta )(\overline{\xi }\backslash x_1)\hfill \\ =& \frac{(s+1)!(t+d2)!}{(s+t)!(d2)!}a(\overline{\xi }_2,\eta )(\xi \eta )(\overline{\xi }\backslash x_1)\hfill \\ \stackrel{\text{2., (}\text{1.1}\text{)}}{=}& (s+1)a(\xi \eta )(\overline{\xi }\backslash x_1)\hfill \\ & t\left[a(x_1,y)\left((\xi y)(\eta \backslash y)\right)(\overline{\xi }\backslash x_1)\right]\hfill \\ \stackrel{\text{case }d\text{}t1}{=}& \frac{s+1}{t+d1}a(\xi \eta )\overline{\xi }\hfill \\ & t\left[\frac{(s+1)!(t+d2)!}{(s+t)!(d1)!}a(x_1,y)(\overline{\xi }_1\backslash x_1,\eta \backslash y)\left((\xi y)(\eta \backslash y)\right)(\overline{\xi }\backslash x_1)\right].\hfill \end{array}$$
Therefore,
$$\begin{array}{ccc}\hfill a(\xi \eta )\overline{\xi }& =& \frac{s!(t+d1)!}{(s+t)!(d2)!}a(\overline{\xi }_2,\eta )(\xi \eta )(\overline{\xi }\backslash x_1)\hfill \\ & +& t\frac{s!(t+d1)!}{(s+t)!(d1)!}a(\overline{\xi }_1,\eta )(\xi \eta )(\overline{\xi }\backslash x_1)\hfill \\ & \stackrel{\text{(}\text{1.3}\text{)}}{=}& (d1)\frac{s!(t+d1)!}{(s+t)!(d1)!}a(\overline{\xi }_1,\eta )(\xi \eta )(x_0,x_1)(\overline{\xi }_0\overline{\xi }_1\backslash x_1)\hfill \\ & +& t\frac{s!(t+d1)!}{(s+t)!(d1)!}a(\overline{\xi }_1,\eta )(\xi \eta )(\overline{\xi }_0\overline{\xi }_1\backslash x_1)\hfill \\ & \stackrel{\text{(}\text{1.1}\text{)}}{=}& \frac{s!(t+d1)!}{(s+t)!(d1)!}a(\overline{\xi }_1,\eta )(\xi \eta )\overline{\xi }.\hfill \end{array}$$
###### Proposition 1.5 (Garnir formula)
Assume given a subset $`\phi _{j[1,p1]}a_j`$. We have
$$a(\phi \xi \eta )=\frac{s!}{(s+t)!(d1)!}a(\overline{\xi }_1,\eta )\overline{\xi }(\phi \xi \eta ).$$
Let $`A`$ be the set of injections of $`\phi `$ into $`\phi \xi \eta `$. We write
$$\begin{array}{ccc}\hfill a(\phi \xi \eta )& =& \underset{\alpha A}{}\underset{\sigma 𝒮_{\phi \xi \eta },\sigma |_\phi =\alpha }{}a\sigma \epsilon _\sigma \hfill \\ \hfill a(\overline{\xi }_1,\eta )\overline{\xi }(\phi \xi \eta )& =& \underset{\alpha A}{}\underset{\sigma 𝒮_{\phi \xi \eta },\sigma |_\phi =\alpha }{}a(\overline{\xi }_1,\eta )\overline{\xi }\sigma \epsilon _\sigma \hfill \end{array}$$
and fix an embedding $`\alpha A`$ in order to compare the summands. We choose a permutation $`\rho 𝒮_{\phi \xi \eta }`$ that restricts to $`\rho |_\phi =\alpha `$. On the one hand, we obtain
$$\begin{array}{ccc}\hfill \underset{\sigma 𝒮_{\phi \xi \eta },\sigma |_\phi =\alpha }{}a\sigma \epsilon _\sigma & \stackrel{\sigma =\sigma ^{}\rho }{=}& \underset{\sigma ^{}𝒮_{\xi \eta }}{}a\sigma ^{}\epsilon _\sigma ^{}\rho \epsilon _\rho \hfill \\ & =& a(\xi \eta )\rho \epsilon _\rho \hfill \\ & \stackrel{\text{(}\text{1.4}\text{)}}{=}& \frac{s!}{(d1)!}a(\overline{\xi }_1,\eta )\overline{\xi }\rho \epsilon _\rho ,\hfill \end{array}$$
on the other hand, we have
$$\begin{array}{ccc}\hfill \underset{\sigma 𝒮_{\phi \xi \eta },\sigma |_\phi =\alpha }{}a(\overline{\xi }_1,\eta )\overline{\xi }\sigma \epsilon _\sigma & \stackrel{\sigma =\sigma ^{}\rho }{=}& \underset{\sigma ^{}𝒮_{\xi \eta }}{}a(\overline{\xi }_1,\eta )\overline{\xi }\sigma ^{}\epsilon _\sigma ^{}\rho \epsilon _\rho \hfill \\ & =& (s+t)!a(\overline{\xi }_1,\eta )\overline{\xi }\rho \epsilon _\rho ,\hfill \end{array}$$
whence the result follows by comparison. Independence of the right hand side of the choice of $`\overline{\xi }_1`$ can also be seen by application of (1.3) to an appropriate permutation $`\tau 𝒮_{\overline{\xi }}`$, for we have $`()\tau \overline{\xi }=()\overline{\xi }\epsilon _\tau `$.
###### Lemma 1.6
For $`i_0[0,s]`$, we write
$$B_{[a],\xi ,\eta }(i_0):=\underset{s_0[i_0,s]}{}(1)^{s_0}\underset{\stackrel{\xi _0\xi ,}{\mathrm{\#}\xi _0=s_0}}{}\underset{\stackrel{\eta _0\eta ,}{\mathrm{\#}\eta _0=s_0i_0}}{}\underset{\stackrel{\phi _0\overline{\eta },}{\mathrm{\#}\phi _0=i_0}}{}a(\xi _0,\eta _0\phi _0).$$
Let $`i[1,u]`$. We define and obtain
$$\begin{array}{ccc}\hfill C_{[a],\xi ,\eta }(i)& :=& \underset{\stackrel{\phi \overline{\eta },}{\mathrm{\#}\phi =i}}{}a(\xi \eta \phi )\hfill \\ & =& s!(t+i)!\underset{i_0[0,\mathrm{min}(i,s)]}{}\left(\begin{array}{c}ui_0\\ ii_0\end{array}\right)B_{[a],\xi ,\eta }(i_0).\hfill \end{array}$$
We fix a subset $`\phi \overline{\eta }`$ of cardinality $`\mathrm{\#}\phi =i`$. Indexing by subsets that become interchanged by a certain element of $`𝒮_{\xi \eta \phi }`$, we obtain
$$\begin{array}{c}a(\xi \eta \phi )\hfill \\ =s!(t+i)!\underset{i_0[0,\mathrm{min}(i,s)]}{}\underset{s_0[i_0,s]}{}\underset{\stackrel{\xi _0\xi ,}{\mathrm{\#}\xi _0=s_0}}{}\underset{\stackrel{\eta _0\eta ,}{\mathrm{\#}\eta _0=s_0i_0}}{}\underset{\stackrel{\phi _0\phi ,}{\mathrm{\#}\phi _0=i_0}}{}a(\xi _0,\eta _0\phi _0)\epsilon _{(\xi _0,\eta _0\phi _0)}.\hfill \end{array}$$
The number of subsets $`\phi `$ of cardinality $`\mathrm{\#}\phi =i`$ that lie in between $`\phi _0\phi \overline{\eta }`$, where $`\phi _0`$ is a given subset of $`\overline{\eta }`$ of cardinality $`\mathrm{\#}\phi _0=i_0`$, amounts to $`\left(\begin{array}{c}ui_0\\ ii_0\end{array}\right)`$. Whence the result is obtained by collecting
$$\underset{\stackrel{\phi \overline{\eta },}{\mathrm{\#}\phi =i}}{}\underset{\stackrel{\phi _0\phi ,}{\mathrm{\#}\phi _0=i_0}}{}U(\phi _0)=\left(\begin{array}{c}ui_0\\ ii_0\end{array}\right)\underset{\stackrel{\phi _0\overline{\eta },}{\mathrm{\#}\phi _0=i_0}}{}U(\phi _0),$$
$`U(\phi _0)`$ being some expression independent of $`\phi `$.
###### Lemma 1.7
Suppose given in addition a subset $`\psi \xi `$, of cardinality $`\mathrm{\#}\psi =:v`$. For $`i_0[0,sv]`$, we write
$$B_{[a],\xi ,\psi ,\eta }^{}(i_0):=\frac{(1)^{i_0}}{v!t!}\underset{\stackrel{\xi _0\xi \backslash \psi ,}{\mathrm{\#}\xi _0=i_0}}{}\underset{\stackrel{\phi _0\eta ,}{\mathrm{\#}\phi _0=i_0}}{}a(\xi _0,\phi _0)(\psi \eta ).$$
Let $`i[1,u]`$. We define and obtain
$$\begin{array}{ccc}\hfill C_{[a],\xi ,\psi ,\eta }^{}(i)& :=& \underset{\stackrel{\phi \eta ,}{\mathrm{\#}\phi =i}}{}a(\xi \phi )(\psi \eta )\hfill \\ & =& t!s!(i+v)!\underset{i_0[0,\mathrm{min}(i,sv)]}{}\left(\begin{array}{c}ti_0\\ ii_0\end{array}\right)B_{[a],\xi ,\psi ,\eta }^{}(i_0).\hfill \end{array}$$
We fix a subset $`\phi \eta `$ of cardinality $`\mathrm{\#}\phi =i`$ and obtain
$$\begin{array}{cc}& \frac{1}{s!i!}a(\xi \phi )(\psi \eta )\hfill \\ =& \frac{1}{(sv)!i!}\left(\begin{array}{c}i+v\\ v\end{array}\right)a((\xi \backslash \psi )\phi )(\psi \eta )\hfill \\ =& \left(\begin{array}{c}i+v\\ v\end{array}\right)\underset{i_0[0,\mathrm{min}(i,sv)]}{}\underset{\stackrel{\xi _0\xi \backslash \psi ,}{\mathrm{\#}\xi _0=i_0}}{}\underset{\stackrel{\phi _0\phi ,}{\mathrm{\#}\phi _0=i_0}}{}(1)^{i_0}a(\xi _0,\phi _0)(\psi \eta ),\hfill \end{array}$$
whence we conclude as in (1.6) that
$$\begin{array}{cc}& \underset{\stackrel{\phi \eta ,}{\mathrm{\#}\phi =i}}{}a(\xi \phi )(\psi \eta )\hfill \\ =& \underset{\stackrel{\phi \eta ,}{\mathrm{\#}\phi =i}}{}\frac{s!(i+v)!}{v!}\underset{i_0[0,\mathrm{min}(i,sv)]}{}\underset{\stackrel{\xi _0\xi \backslash \psi ,}{\mathrm{\#}\xi _0=i_0}}{}\underset{\stackrel{\phi _0\phi ,}{\mathrm{\#}\phi _0=i_0}}{}(1)^{i_0}a(\xi _0,\phi _0)(\psi \eta )\hfill \\ =& \frac{s!(i+v)!}{v!}\underset{i_0[0,\mathrm{min}(i,sv)]}{}\underset{\stackrel{\xi _0\xi \backslash \psi ,}{\mathrm{\#}\xi _0=i_0}}{}\left(\begin{array}{c}ti_0\\ ii_0\end{array}\right)\underset{\stackrel{\phi _0\eta ,}{\mathrm{\#}\phi _0=i_0}}{}(1)^{i_0}a(\xi _0,\phi _0)(\psi \eta ).\hfill \end{array}$$
## 2 A two-box-shift morphism
### 2.1 Double paths
Let $`\lambda `$ be a partition of $`n`$. We assume given integers $`g,k`$ such that $`1gk\lambda _11`$ and such that
$$\mu _i^{}:=\{\begin{array}{cc}\lambda _i^{}+2\hfill & \text{for }i=g\hfill \\ \lambda _i^{}2\hfill & \text{for }i=k+1\hfill \\ \lambda _i^{}\hfill & \text{else}\hfill \end{array}$$
defines a partition $`\mu `$. A weight $`e`$ is a map
$$\begin{array}{ccc}\hfill [1,\lambda _1]& \text{}& [0,2]\hfill \\ \hfill j& \text{}& e_j\hfill \end{array}$$
that maps $`g`$ and $`k+1`$ to $`e_g=e_{k+1}=2`$ and that maps each $`j[1,\lambda _1]\backslash [g,k+1]`$ to $`e_j=0`$. The set of weights is denoted by $`E`$. A pattern $`\mathrm{\Xi }`$ of weight $`e`$ is a subset $`\mathrm{\Xi }[1,2]\times [g,k+1]`$ such that
$$\mathrm{\#}(\mathrm{\Xi }([1,2]\times \{j\}))=e_j$$
for $`j[g,k+1]`$. A double path $`\gamma `$ of weight $`e`$ is an injection from a pattern $`\mathrm{\Xi }`$ of weight $`e`$ to $`[\lambda ][\mu ]`$ of the form
$$\begin{array}{ccc}\hfill \mathrm{\Xi }& \text{}& [\lambda ][\mu ]\hfill \\ \hfill i\times j& \text{}& \overline{\gamma }(j,i)\times j\hfill \end{array}$$
such that
$$\begin{array}{cc}\overline{\gamma }(g,i)=\lambda _g^{}+i\hfill & \text{for }i[1,2]\hfill \\ \overline{\gamma }(k+1,1)<\overline{\gamma }(k+1,2).\hfill & \end{array}$$
Sometimes, we denote its pattern by $`\mathrm{\Xi }_\gamma :=\mathrm{\Xi }`$. Its sign is given by $`\epsilon _\gamma :=(1)^{\overline{\gamma }(k+1,1)+\overline{\gamma }(k+1,2)}`$, not to be confused with the sign of a permutation. The set of double paths of weight $`e`$ is denoted by $`\dot{\mathrm{\Gamma }}(e)`$.
A double path $`\gamma `$ gives rise to a place operation in the following manner. Let
$$\begin{array}{ccc}\hfill [0,\lambda _{k+1}^{}]& \text{}& [0,\mu _{k+1}^{}]\hfill \\ \hfill i& \text{}& \mathrm{\#}\left([1,i]\backslash \{\overline{\gamma }(k+1,1),\overline{\gamma }(k+1,2)\}\right)\hfill \\ \hfill \mathrm{min}(\phi ^1(\{j\}))& \text{}& j.\hfill \end{array}$$
Given a weight $`eE`$, we define an operation of sets
$$\begin{array}{ccccc}\hfill \dot{\mathrm{\Gamma }}(e)& \times & T^\lambda \hfill & \text{}& T^\mu \hfill \\ \hfill \gamma & \times & [a]\hfill & \text{}& [a^\gamma ]\hfill \end{array}$$
by
$$\begin{array}{cccc}\hfill a_{j,i}^\gamma & :=& a_{j,i}\text{ for }i\times j[\mu ]\backslash (\mathrm{\Xi })\gamma \text{ and }jk+1,\hfill & \\ \hfill a_{j,\overline{\gamma }(j,i)}^\gamma & :=& a_{j^{},\overline{\gamma }(j^{},i)}\text{ for }i[1,2]\text{ and for }j[g,k]\text{,}\hfill & \\ & & \text{where }j^{}\text{ is minimal in }[j+1,k+1]\text{ with }i\times j^{}\mathrm{\Xi },\hfill & \\ \hfill a_{k+1,i}^\gamma & :=& a_{k+1,\psi (i)}\text{ for }i[1,\mu _{k+1}^{}].\hfill & \end{array}$$
We define a $`\text{Z}𝒮_n`$-linear map
$$\begin{array}{ccc}\hfill F^\lambda & \text{}& S^\mu \hfill \\ \hfill [a]& \text{}& \underset{\gamma \dot{\mathrm{\Gamma }}(e)}{}a^\gamma \epsilon _\gamma .\hfill \end{array}$$
###### Example 2.1
For example, consider the case $`\lambda =(7,7,7,4)`$, $`\mu =(7,6,6,4,1,1)`$, $`g=1`$, $`k=6`$. Let $`e=(2,1,1,1,2,0,2)`$, and let $`\gamma `$ be the (non-ordered) double path of weight $`e`$ given pictorially by
a ‘u’ indicating the $`\gamma `$-image of some $`1\times j`$, a ‘v’ indicating the $`\gamma `$-image of some $`2\times j`$. The operation of $`\gamma `$ yields e.g.
I.e. the double path $`\gamma `$ subsequently pushes $`23\text{}19\text{}15\text{}6\text{}`$ and $`25\text{}18\text{}9\text{}`$.
###### Lemma 2.2 (path switch)
For $`q[g+1,k+1]`$, we dispose of an involution
$$\begin{array}{ccccc}\hfill [1,2]& \times & [g,k+1]\hfill & \text{}& [1,2]\times [g,k+1]\hfill \\ \hfill i& \times & j\hfill & \text{}& \{\begin{array}{cccc}\hfill i(1,2)& \times & j\hfill & \text{for }j[g+1,q1]\hfill \\ \hfill i& \times & j\hfill & \text{else}\hfill \end{array}.\hfill \end{array}$$
The composition $`\iota _p\gamma `$ of this involution, restricted to $`\iota _p^1(\mathrm{\Xi }_\gamma )`$, with a double path $`\gamma `$ furnishes a double path which we denote by $`\iota _p\gamma `$, being a slight abuse of notation. So $`\mathrm{\Xi }_{\iota _p\gamma }=\iota _p^1(\mathrm{\Xi }_\gamma )`$.
Let $`[a]T^\lambda `$, $`eE`$, let $`p[g+1,k]`$ such that $`e_p=2`$, let $`\gamma \dot{\mathrm{\Gamma }}(e)`$, $`u:=a_{p,\overline{\gamma }(p,1)}`$ and $`v:=a_{p,\overline{\gamma }(p,2)}`$. We assert the following equalities.
$$\begin{array}{cccc}\hfill \text{(i)}& \hfill a^\gamma (u,v)& =& a^{\iota _p\gamma }\hfill \\ \hfill \text{(ii)}& \hfill a^{\iota _{p+1}\gamma }& =& a^\gamma \hfill \\ \hfill \text{(iii)}& \hfill ([a]f_e^{\prime \prime })(u,v)& =& [a]f_e^{\prime \prime }\hfill \end{array}$$
The operation of the double path $`\iota _p\gamma `$ on $`[a]`$ can be performed using the operation of $`\gamma `$, but with entries $`u`$ and $`v`$ interchanged, and, moreover, with entries placed by $`\gamma `$ in column $`g`$ interchanged, yielding the sign in (i). Analoguously (ii). Since composition with $`\iota _p`$ is an involution on $`\dot{\mathrm{\Gamma }}(e)`$, multiplication by $`(u,v)`$ attaches a sign to $`[a]f_e^{\prime \prime }`$ by (i), whence (iii).
###### Proposition 2.3
Let $`eE`$. There is a factorization
$$(F^\lambda \text{}S^\mu )=(F^\lambda \text{}M^{\lambda ^{},}\text{}S^\mu ).$$
We need to show that a signed column transposition $`[a]+[a](u,v)`$, $`[a]`$ a $`\lambda `$-tableau, $`p1`$, $`u,va_p`$, $`uv`$, vanishes under $`f_e^{\prime \prime }`$, for which we may assume $`p`$ to lie in $`[g+1,k+1]`$. We abbreviate the double path image of $`\gamma `$ in $`[a]`$ by $`I_\gamma :=((\mathrm{\Xi })\gamma )[a]`$.
Case $`p[g+1,k]`$.
Subcase $`u,vI_\gamma `$. See (2.2, iii).
Subcase $`u=:a_{p,\overline{\gamma }(p,i)}I_\gamma `$, $`v=:a_{p,l}I_\gamma `$. Abbreviating by $`\tau _\gamma `$ the transposition that interchanges the elements $`\overline{\gamma }(p,i)\times p`$ and $`l\times p`$ of $`[\lambda ][\mu ]`$, we conclude from
$$a^\gamma (u,v)=a^{\gamma \tau _\gamma }$$
that the sum over these summands, multiplied by $`(u,v)`$, yields minus the sum over the summands of the subcase $`uI_\gamma ,vI_\gamma `$. In fact, $`\gamma \text{}\gamma \tau _\gamma `$ is a bijection between the sets of double paths occurring in these subcases which preserves the sign.
Subcase $`u,vI_\gamma `$. $`a^\gamma (u,v)=a^\gamma `$ shows these summands to yield zero.
Case $`p=k+1`$. The arguments of the case $`p[g+1,k]`$ work in the subcases $`u,vI_\gamma `$ and $`u,vI_\gamma `$ as well.
Subcase $`u=:a_{k+1,\overline{\gamma }(k+1,i)}I_\gamma `$, $`v=:a_{k+1,l}I_\gamma `$. Writing the transposition $`\tau _\gamma `$ as above, we conclude from
$$a^\gamma (u,v)=(1)^{\overline{\gamma }(p,i)l+1}a^{\gamma \tau _\gamma }$$
that the sum over these summands yields minus the result of the sum over the summands of the subcase $`uI_\gamma ,vI_\gamma `$. In fact, $`\gamma \text{}\gamma \tau _\gamma `$ is a bijection between the sets of double paths occurring in these subcases which changes the sign by $`\epsilon _{\gamma \tau _\gamma }/\epsilon _\gamma =(1)^{l\overline{\gamma }(p,i)}`$.
### 2.2 Morphism, unreduced version
We shall show the vanishing, modulo our modulus, of the required Garnir relations. Rather than giving short versions of all calculations, we give full versions, but only of the two ‘large’ cases (A.1, A.2), and only in Appendix A.
For $`i[g,k+1]`$, we let
$$X_i:=(\lambda _i^{}i)(\lambda _{k+1}^{}(k+1))$$
and denote $`X_i^{(j)}:=(X_i+j1)!/(X_i1)!`$, $`j[0,2]`$. Let $`M^{\lambda ^{},}\text{}S^\mu `$ be the $`\text{Z}𝒮_n`$-linear map defined by
$$f^{}:=\underset{eE}{}\left(\underset{i[g+1,k]}{}X_i^{(2e_i)}\right)f_e^{}.$$
Let $`[a]`$ be a $`\lambda `$-tableau, $`p[g,k]`$, $`\xi a_p`$, $`s:=\mathrm{\#}\xi `$, $`\overline{\xi }:=a_p\backslash \xi `$ and $`\eta a_{p+1}`$, $`t:=\mathrm{\#}\eta `$, $`\overline{\eta }:=a_{p+1}\backslash \eta `$ such that $`s+t=\lambda _p^{}+1`$. Denote $`u:=\mathrm{\#}\overline{\eta }=\lambda _{p+1}^{}t`$.
###### Lemma 2.4 (cf. (A.1))
Suppose $`g<p<k`$ and $`s,t2`$. The map $`f^{}`$ annihilates $`G_{[a],\xi ,\eta }^{}`$.
We fix a map
$$\begin{array}{ccc}\hfill [1,\lambda _1]\backslash \{p,p+1\}& \text{}& [0,2]\hfill \\ \hfill j& \text{}& \stackrel{~}{e}_j\hfill \end{array}$$
that maps $`g`$ and $`k+1`$ to $`e_g=e_{k+1}=2`$, and that maps $`j[1,\lambda _1]\backslash [g,k+1]`$ to $`e_j=0`$. For $`\alpha ,\beta [0,2]`$, we denote by $`\stackrel{~}{e}\alpha \beta `$ be the prolongation of $`\stackrel{~}{e}`$ to $`[1,\lambda _1]`$ defined by $`\stackrel{~}{e}\alpha \beta |_{[1,\lambda _1]\backslash \{p,p+1\}}:=\stackrel{~}{e}`$, $`(\stackrel{~}{e}\alpha \beta )_p:=\alpha `$ and $`(\stackrel{~}{e}\alpha \beta )_{p+1}:=\beta `$. We contend that
$$\underset{\alpha ,\beta [0,2]}{}X_p^{(2\alpha )}X_{p+1}^{(2\beta )}G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}\alpha \beta }^{}=0,$$
from which the lemma ensues.
There exist elements $`x,y\xi `$, $`xy`$, $`x^{},y^{}\eta `$, $`x^{}y^{}`$, $`z\overline{\xi }`$, which we choose and fix. For $`\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e}00)`$, we let $`x_\gamma :=a_{j,\overline{\gamma }(j,1)}`$, where $`j[p+2,k+1]`$ is minimal with $`1\times j\mathrm{\Xi }_\gamma `$, and $`y_\gamma :=a_{j,\overline{\gamma }(j,2)}`$, where $`j[p+2,k+1]`$ is minimal with $`2\times j\mathrm{\Xi }_\gamma `$. I.e. we pick the entries $`x_\gamma ,y_\gamma `$ that ‘cross the columns’ $`p`$ and $`p+1`$ under the operation of $`\gamma `$. We write
$$\begin{array}{ccc}\hfill U_\gamma & :=& (s+t2)!^1a^\gamma \epsilon _\gamma (\overline{\xi }\backslash z,\eta \backslash \{x^{},y^{}\})(x_\gamma ,x,x^{})(y_\gamma ,y,y^{})(\xi \eta )\hfill \\ \hfill V_{1,\gamma }& :=& (s+t1)!^1a^\gamma \epsilon _\gamma (\overline{\xi },\eta \backslash x^{})(x_\gamma ,x,x^{})(\xi \eta )\hfill \\ \hfill V_{2,\gamma }& :=& (s+t1)!^1a^\gamma \epsilon _\gamma (\overline{\xi },\eta \backslash y^{})(y_\gamma ,y,y^{})(\xi \eta ),\hfill \end{array}$$
and let
$$\begin{array}{ccc}A\hfill & :=& \underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e}00)}{}U_\gamma \underset{w^{}\overline{\eta }}{}(w^{},z)\hfill \\ B\hfill & :=& \underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e}00)}{}U_\gamma (1\underset{w\overline{\xi }\backslash z}{}(w,z))\hfill \\ C_1\hfill & :=& \underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e}00)}{}U_\gamma (z,x_\gamma )\hfill \\ C_2\hfill & :=& \underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e}00)}{}U_\gamma (z,y_\gamma )\hfill \\ D\hfill & :=& \underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e}00)}{}V_{1,\gamma }\underset{w\overline{\xi }}{}(w,y_\gamma )\hfill \\ & =& \underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e}00)}{}V_{2,\gamma }\underset{w\overline{\xi }}{}(w,x_\gamma )\hfill \\ H\hfill & :=& \underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e}00)}{}V_{1,\gamma }\underset{w^{}\overline{\eta }}{}(w^{},y_\gamma )\hfill \\ & =& \underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e}00)}{}V_{2,\gamma }\underset{w^{}\overline{\eta }}{}(w^{},x_\gamma )\hfill \\ F_1\hfill & :=& \underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e}00)}{}V_{1,\gamma }\hfill \\ F_2\hfill & :=& \underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e}00)}{}V_{2,\gamma }.\hfill \end{array}$$
Calculations carried out in (A.1) yield the following table.
$$\begin{array}{ccc}\hfill G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}22}^{}& =& 2(su)A+(su)(su1)B+2(su)(su+1)D\hfill \\ \hfill G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}12}^{}& =& 2A2(su)B2(su+1)D+2(su+1)H\hfill \\ \hfill G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}02}^{}& =& B2H\hfill \\ \hfill G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}21}^{}& =& 2A+2(su1)B(su1)(C_1+C_2)+4(su)D\hfill \\ \hfill G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}11}^{}& =& 2B+(C_1+C_2)2D+2H+(su)(F_1+F_2)\hfill \\ \hfill G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}01}^{}& =& (F_1+F_2)\hfill \\ \hfill G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}20}^{}& =& B(C_1+C_2)+2D\hfill \\ \hfill G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}10}^{}& =& (F_1+F_2)\hfill \\ \hfill G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}00}^{}& =& 0\hfill \end{array}$$
We note that $`su=X_pX_{p+1}`$ and evaluate the linear combination
$$\begin{array}{c}\underset{\alpha ,\beta [0,2]}{}X_p^{(2\alpha )}X_{p+1}^{(2\beta )}G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}\alpha \beta }^{}\hfill \\ \begin{array}{cccc}=& 1\hfill & 1\hfill & (2(X_pX_{p+1})A+(X_pX_{p+1})(X_pX_{p+1}1)B\hfill \\ & & & +2(X_pX_{p+1})(X_pX_{p+1}+1)D)\hfill \\ +& X_p\hfill & 1\hfill & (2A2(X_pX_{p+1})B2(X_pX_{p+1}+1)D\hfill \\ & & & +2(X_pX_{p+1}+1)H)\hfill \\ +& X_p(X_p+1)\hfill & 1\hfill & \left(B2H\right)\hfill \\ +& 1\hfill & X_{p+1}\hfill & (2A+2(X_pX_{p+1}1)B\hfill \\ & & & (X_pX_{p+1}1)(C_1+C_2)+4(X_pX_{p+1})D)\hfill \\ +& X_p\hfill & X_{p+1}\hfill & (2B+(C_1+C_2)2D+2H\hfill \\ & & & +(X_pX_{p+1})(F_1+F_2))\hfill \\ +& X_p(X_p+1)\hfill & X_{p+1}\hfill & \left((F_1+F_2)\right)\hfill \\ +& 1\hfill & X_{p+1}(X_{p+1}+1)\hfill & \left(B(C_1+C_2)+2D\right)\hfill \\ +& X_p\hfill & X_{p+1}(X_{p+1}+1)\hfill & \left((F_1+F_2)\right)\hfill \\ +& X_p(X_p+1)\hfill & X_{p+1}(X_{p+1}+1)\hfill & \left(0\right)\hfill \\ =& 0.\hfill & & \end{array}\hfill \end{array}$$
###### Lemma 2.5
Suppose $`g<p<k`$ and $`s=\lambda _p^{}`$, $`t=1`$. The map $`f^{}`$ annihilates $`G_{[a],\xi ,\eta }^{}`$.
Concerning the map $`\stackrel{~}{e}`$ and its prolongations, as well as concerning the integers $`x_\gamma `$, $`y_\gamma `$, we continue to use the notation of (2.4). We fix such a map $`\stackrel{~}{e}`$ and need to show that
$$\underset{\alpha ,\beta [0,2]}{}X_p^{(2\alpha )}X_{p+1}^{(2\beta )}G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}\alpha \beta }^{}=0.$$
There exist elements $`x,y\xi `$, $`xy`$, which we choose and fix. For $`\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e}00)`$, we write
$$\begin{array}{ccc}\hfill V_{1,\gamma }& :=& \frac{1}{(s1)!}a^\gamma \epsilon _\gamma (y_\gamma ,y,\eta )(\xi \eta )\hfill \\ \hfill V_{2,\gamma }& :=& \frac{1}{(s1)!}a^\gamma \epsilon _\gamma (x_\gamma ,x,\eta )(\xi \eta )\hfill \end{array}$$
and let
$$\begin{array}{ccc}A\hfill & :=& \underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e}00)}{}V_{1,\gamma }\underset{w^{}\overline{\eta }}{}(x_\gamma ,x,w^{})\hfill \\ & =& \underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e}00)}{}V_{2,\gamma }\underset{w^{}\overline{\eta }}{}(y_\gamma ,y,w^{})\hfill \\ C_1\hfill & :=& \underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e}00)}{}V_{1,\gamma }(x_\gamma ,x)\hfill \\ C_2\hfill & :=& \underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e}00)}{}V_{2,\gamma }(y_\gamma ,y)\hfill \end{array}$$
$$\begin{array}{ccc}H\hfill & :=& s^1\underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e}00)}{}V_{1,\gamma }\underset{w^{}\overline{\eta }}{}(x_\gamma ,w^{})\hfill \\ & =& s^1\underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e}00)}{}V_{2,\gamma }\underset{w^{}\overline{\eta }}{}(y_\gamma ,w^{}).\hfill \\ F_1\hfill & :=& s^1\underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e}00)}{}V_{1,\gamma }\hfill \\ F_2\hfill & :=& s^1\underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e}00)}{}V_{2,\gamma }\hfill \end{array}$$
Calculations similar to those of (A.1) yield the following table, to be compared to the table of (2.4).
$$\begin{array}{ccc}\hfill G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}22}^{}& =& 2(su)A\hfill \\ \hfill G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}12}^{}& =& 2A+2(su+1)H\hfill \\ \hfill G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}02}^{}& =& 2H\hfill \\ \hfill G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}21}^{}& =& 2A(su1)(C_1+C_2)\hfill \\ \hfill G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}11}^{}& =& (C_1+C_2)+2H+(su)(F_1+F_2)\hfill \\ \hfill G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}01}^{}& =& (F_1+F_2)\hfill \\ \hfill G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}20}^{}& =& (C_1+C_2)\hfill \\ \hfill G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}10}^{}& =& (F_1+F_2)\hfill \\ \hfill G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}00}^{}& =& 0\hfill \end{array}$$
###### Lemma 2.6
Suppose $`g<p<k`$ and $`s=1`$, $`t=\lambda _p^{}=\lambda _{p+1}^{}`$. The map $`f^{}`$ annihilates $`G_{[a],\xi ,\eta }^{}`$.
Concerning $`\stackrel{~}{e}`$, $`x_\gamma `$, $`y_\gamma `$, we continue to use the notation of (2.4). We fix such a map $`\stackrel{~}{e}`$ and need to show that
$$\underset{\alpha ,\beta [0,2]}{}X_p^{(2\alpha )}X_{p+1}^{(2\beta )}G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}\alpha \beta }^{}=0.$$
There exist elements $`x^{},y^{}\eta `$, $`x^{}y^{}`$, $`z\overline{\xi }`$, which we choose and fix. For $`\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e}00)`$, we write
$$\begin{array}{ccc}\hfill U_{1,\gamma }& :=& \frac{1}{(t1)!}a^\gamma \epsilon _\gamma (\overline{\xi }\backslash z,\eta \backslash \{x^{},y^{}\})(x_\gamma ,\xi ,x^{})(y_\gamma ,y^{})(\xi \eta )\hfill \\ \hfill U_{2,\gamma }& :=& \frac{1}{(t1)!}a^\gamma \epsilon _\gamma (\overline{\xi }\backslash z,\eta \backslash \{x^{},y^{}\})(y_\gamma ,\xi ,y^{})(x_\gamma ,x^{})(\xi \eta )\hfill \\ \hfill V_{1,\gamma }& :=& \frac{1}{t!}a^\gamma \epsilon _\gamma (\overline{\xi },\eta \backslash x^{})(x_\gamma ,\xi ,x^{})(\xi \eta )\hfill \\ \hfill V_{2,\gamma }& :=& \frac{1}{t!}a^\gamma \epsilon _\gamma (\overline{\xi },\eta \backslash y^{})(y_\gamma ,\xi ,y^{})(\xi \eta )\hfill \end{array}$$
and let
$$\begin{array}{ccc}B\hfill & :=& \underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e}00)}{}U_{1,\gamma }(1\underset{w\overline{\xi }\backslash z}{}(w,z))\hfill \\ & =& \underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e}00)}{}U_{2,\gamma }(1\underset{w\overline{\xi }\backslash z}{}(w,z))\hfill \\ C_1\hfill & :=& \underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e}00)}{}U_{1,\gamma }(z,x_\gamma )\hfill \end{array}$$
$$\begin{array}{ccc}C_2\hfill & :=& \underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e}00)}{}U_{2,\gamma }(z,y_\gamma )\hfill \\ D\hfill & :=& \underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e}00)}{}V_{1,\gamma }\underset{w\overline{\xi }}{}(w,y_\gamma )\hfill \\ & =& \underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e}00)}{}V_{2,\gamma }\underset{w\overline{\xi }}{}(w,x_\gamma )\hfill \\ F_1\hfill & :=& \underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e}00)}{}V_{1,\gamma }\hfill \\ F_2\hfill & :=& \underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e}00)}{}V_{2,\gamma }\hfill \end{array}$$
Calculations similar to those of (A.1), using the Garnir relation $`B=C_1+C_2`$, yield the following table, to be compared to the table of (2.4).
$$\begin{array}{ccc}\hfill G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}22}^{}& =& 4D\hfill \\ \hfill G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}12}^{}& =& 2B4D\hfill \\ \hfill G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}02}^{}& =& B\hfill \\ \hfill G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}21}^{}& =& 4D\hfill \\ \hfill G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}11}^{}& =& B2D+(F_1+F_2)\hfill \\ \hfill G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}01}^{}& =& (F_1+F_2)\hfill \\ \hfill G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}20}^{}& =& 2D\hfill \\ \hfill G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}10}^{}& =& (F_1+F_2)\hfill \\ \hfill G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}00}^{}& =& 0\hfill \end{array}$$
Note that $`su=X_pX_{p+1}=1`$.
###### Lemma 2.7
Suppose $`g<p=k`$ and $`s,t2`$, $`\eta =a_{k+1,[1,t]}`$. The map $`f^{}`$ annihilates $`G_{[a],\xi ,\eta }^{}`$.
We fix a map
$$\begin{array}{ccc}\hfill [1,\lambda _1]\backslash \{k\}& \text{}& [0,2]\hfill \\ \hfill j& \text{}& \stackrel{~}{e}_j\hfill \end{array}$$
that maps $`g`$ and $`k+1`$ to $`e_g=e_{k+1}=2`$, and that maps $`j[1,\lambda _1]\backslash [g,k+1]`$ to $`e_j=0`$. For $`\alpha [0,2]`$, $`\stackrel{~}{e}\alpha `$ denotes the map which prolongs $`\stackrel{~}{e}`$ to $`[1,\lambda _1]`$ via $`(\stackrel{~}{e}\alpha )_k=\alpha `$. We need to show that
$$\underset{\alpha [0,2]}{}X_k^{(2\alpha )}G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}\alpha }^{}=0.$$
There exist elements $`x,y\xi `$, $`xy`$, $`z\overline{\xi }`$, which we choose and fix. Let $`x^{}:=a_{g+1,1}`$ and $`y^{}:=a_{g+1,2}`$, so that $`x^{},y^{}\eta `$. For $`v,wa_{k+1}`$, $`vw`$, we denote
$$\dot{\mathrm{\Gamma }}(\stackrel{~}{e},a,v,w):=\{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e}0)|a_{k+1,\overline{\gamma }(k+1,1)}=v,a_{k+1,\overline{\gamma }(k+1,2)}=w\}$$
and let
$$\begin{array}{ccc}A\hfill & :=& \frac{1}{(s+t2)!}\underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e},a,x^{},y^{})}{}a^\gamma \epsilon _\gamma (\overline{\xi }\backslash z,\eta \backslash \{x^{},y^{}\})(x,x^{})(y,y^{})(\xi \eta )\underset{w^{}\overline{\eta }}{}(w^{},z)\hfill \\ B\hfill & :=& \frac{1}{(s+t2)!}\underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e},a,x^{},y^{})}{}a^\gamma \epsilon _\gamma \hfill \\ & & (\overline{\xi }\backslash z,\eta \backslash \{x^{},y^{}\})(x,x^{})(y,y^{})(\xi \eta )(1\underset{w\overline{\xi }\backslash z}{}(w,z))\hfill \\ D\hfill & :=& \frac{1}{(s+t1)!}\underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e},a,x^{},y^{})}{}\underset{w\overline{\xi }}{}a^\gamma \epsilon _\gamma (\overline{\xi }\backslash w,\eta \backslash \{x^{},y^{}\})(x,x^{})(w,y^{})(\xi \eta )\hfill \\ & =& \frac{1}{(s+t1)!}\underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e},a,x^{},y^{})}{}\underset{w\overline{\xi }}{}a^\gamma \epsilon _\gamma (\overline{\xi }\backslash w,\eta \backslash \{x^{},y^{}\})(y,y^{})(w,x^{})(\xi \eta )\hfill \\ H\hfill & :=& \frac{1}{(s+t1)!}\underset{w^{}\overline{\eta }}{}\underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e},a,x^{},w^{})}{}a^\gamma \epsilon _\gamma (\overline{\xi },\eta \backslash x^{})(x,x^{})(\xi \eta )\hfill \\ & =& \frac{1}{(s+t1)!}\underset{w^{}\overline{\eta }}{}\underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e},a,y^{},w^{})}{}a^\gamma \epsilon _\gamma (\overline{\xi },\eta \backslash y^{})(y,y^{})(\xi \eta ).\hfill \end{array}$$
Calculations similar to those of (A.1) yield the following table.
$$\begin{array}{ccc}\hfill G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}2}^{}& =& \frac{1}{2}\left(2(su)A+(su)(su1)B+2(su)(su+1)D\right)\hfill \\ \hfill G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}1}^{}& =& \frac{1}{2}\left(2A2(su)B2(su+1)D+2(su+1)H\right)\hfill \\ \hfill G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}0}^{}& =& \frac{1}{2}\left(B2H\right)\hfill \end{array}$$
Comparison with the table in (2.4) is possible since $`X_{k+1}=0`$.
###### Lemma 2.8
Suppose $`g<p=k`$ and $`s=\lambda _k^{}`$, $`t=1`$, $`\eta =a_{k+1,1}`$. The map $`f^{}`$ annihilates $`G_{[a],\xi ,\eta }^{}`$.
Concerning $`\stackrel{~}{e}`$ and $`\dot{\mathrm{\Gamma }}`$, we continue to use the notation of (2.7). We fix such a map $`\stackrel{~}{e}`$. There exist elements $`x,y\xi `$, $`xy`$, which we choose and fix. We let
$$\begin{array}{ccc}A\hfill & :=& \frac{1}{(s1)!}\underset{w^{}\overline{\eta }}{}\underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e},a,\eta ,w^{})}{}a^\gamma \epsilon _\gamma (x,\eta )(y,w^{})(\xi \eta )\hfill \\ H\hfill & :=& \frac{1}{s!}\underset{w^{}\overline{\eta }}{}\underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e},a,\eta ,w^{})}{}a^\gamma \epsilon _\gamma (x,\eta )(\xi \eta )\hfill \end{array}$$
Calculations similar to those of (A.1) yield the following table, to be compared to the table of (2.7).
$$\begin{array}{ccc}\hfill G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}2}^{}& =& \frac{1}{2}\left(2(su)A\right)\hfill \\ \hfill G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}1}^{}& =& \frac{1}{2}\left(2A+2(su+1)H\right)\hfill \\ \hfill G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}0}^{}& =& \frac{1}{2}\left(2H\right)\hfill \end{array}$$
###### Lemma 2.9
Suppose $`g<p=k`$ and $`s=1`$, $`t=\lambda _k^{}=\lambda _{k+1}^{}`$. The map $`f^{}`$ annihilates $`G_{[a],\xi ,\eta }^{}`$.
Concerning $`\stackrel{~}{e}`$ and $`\dot{\mathrm{\Gamma }}`$, we continue to use the notation of (2.7). We fix such a map $`\stackrel{~}{e}`$. There exist elements $`x^{}y^{}\eta `$, $`x^{}y^{}`$, $`z\overline{\xi }`$, which we choose and fix. We let
$$\begin{array}{ccc}B\hfill & :=& \frac{1}{(t1)!}\underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e},a,x^{},y^{})}{}a^\gamma \epsilon _\gamma (\overline{\xi }\backslash z,\eta \backslash \{x^{},y^{}\})(\xi \eta )(1\underset{w\overline{\xi }\backslash z}{}(w,z))\hfill \\ D\hfill & :=& \frac{1}{t!}\underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e},a,x^{},y^{})}{}\underset{w\overline{\xi }}{}a^\gamma \epsilon _\gamma (\overline{\xi }\backslash w,\eta \backslash \{x^{},y^{}\})(\xi ,x^{})(w,y^{})(\xi \eta )\hfill \\ & =& \frac{1}{t!}\underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e},a,x^{},y^{})}{}\underset{w\overline{\xi }}{}a^\gamma \epsilon _\gamma (\overline{\xi }\backslash w,\eta \backslash \{x^{},y^{}\})(\xi ,y^{})(w,x^{})(\xi \eta ).\hfill \end{array}$$
Calculations similar to those of (A.1) yield the following table.
$$\begin{array}{ccc}\hfill G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}2}^{}& =& \frac{1}{2}\left(4D\right)\hfill \\ \hfill G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}1}^{}& =& \frac{1}{2}\left(2B4D\right)\hfill \\ \hfill G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}0}^{}& =& \frac{1}{2}\left(B\right)\hfill \end{array}$$
In comparison with the table of (2.7), we note that $`su=1`$.
###### Lemma 2.10 (cf. (A.2))
Suppose $`g=p<k`$ and $`s,t2`$. There exist elements $`x^{},y^{}\eta `$, $`x^{}y^{}`$, $`z\overline{\xi }`$, which we choose and fix. For $`\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e}0)`$, we let $`x_\gamma :=a_{j,\overline{\gamma }(j,1)}`$, where $`j[g+2,k+1]`$ is minimal with $`1\times j\mathrm{\Xi }_\gamma `$, and $`y_\gamma :=a_{j,\overline{\gamma }(j,2)}`$, where $`j[g+2,k+1]`$ is minimal with $`2\times j\mathrm{\Xi }_\gamma `$. I.e. we pick the entries $`x_\gamma ,y_\gamma `$ that ‘cross the column’ $`g+1`$ under the operation of $`\gamma `$. The set of maps
$$\begin{array}{ccc}\hfill [1,\lambda _1]\backslash \{g+1\}& \text{}& [0,2]\hfill \\ \hfill j& \text{}& \stackrel{~}{e}_j\hfill \end{array}$$
that send $`g`$ and $`k+1`$ to $`\stackrel{~}{e}_g=\stackrel{~}{e}_{k+1}=2`$, and that map $`j[1,\lambda _1]\backslash [g,k+1]`$ to $`e_j=0`$, is denoted by $`\stackrel{~}{E}`$. For $`\stackrel{~}{e}\stackrel{~}{E}`$ and $`\beta [0,2]`$, we denote by $`\stackrel{~}{e}\beta `$ the prolongation of $`\stackrel{~}{e}`$ to $`[g,k+1]`$ by $`(\stackrel{~}{e}\beta )_{g+1}=\beta `$. For $`\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e}0)`$, we write
$$U_\gamma :=a^\gamma \epsilon _\gamma (\overline{\xi }\backslash z,\eta \backslash \{x^{},y^{}\})(x_\gamma ,x^{})(y_\gamma ,y^{})$$
and let
$$\begin{array}{ccc}A_{[a],\xi ,\eta ,\stackrel{~}{e}}\hfill & :=& \underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e}0)}{}U_\gamma \underset{w^{}\overline{\eta }}{}(w^{},z)\hfill \\ B_{[a],\xi ,\eta ,\stackrel{~}{e}}\hfill & :=& \underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e}0)}{}U_\gamma (1\underset{w\overline{\xi }\backslash z}{}(w,z))\hfill \\ C_{1,[a],\xi ,\eta ,\stackrel{~}{e}}\hfill & :=& \underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e}0)}{}U_\gamma (z,x_\gamma )\hfill \\ C_{2,[a],\xi ,\eta ,\stackrel{~}{e}}\hfill & :=& \underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e}0)}{}U_\gamma (z,y_\gamma ).\hfill \end{array}$$
We obtain
$$\begin{array}{c}\hfill G_{[a],\xi ,\eta }^{}f^{}=(X_g+2)\underset{\stackrel{~}{e}\stackrel{~}{E}}{}\left(\underset{j[g+2,k]}{}X_j^{(2\stackrel{~}{e}_j)}\right)(2A_{[a],\xi ,\eta ,\stackrel{~}{e}}+(X_g+1)B_{[a],\xi ,\eta ,\stackrel{~}{e}}\\ \hfill X_{g+1}(C_{1,[a],\xi ,\eta ,\stackrel{~}{e}}+C_{2,[a],\xi ,\eta ,\stackrel{~}{e}})).\end{array}$$
Calculations carried out in (A.2) yield the following table.
$$\begin{array}{ccc}\hfill G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}2}^{}& =& 2(su+2)A_{[a],\xi ,\eta ,\stackrel{~}{e}}+(su+2)(su+1)B_{[a],\xi ,\eta ,\stackrel{~}{e}}\hfill \\ \hfill G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}1}^{}& =& 2A_{[a],\xi ,\eta ,\stackrel{~}{e}}+2(su+1)B_{[a],\xi ,\eta ,\stackrel{~}{e}}(su+1)(C_{1,[a],\xi ,\eta ,\stackrel{~}{e}}+C_{2,[a],\xi ,\eta ,\stackrel{~}{e}})\hfill \\ \hfill G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}0}^{}& =& B_{[a],\xi ,\eta ,\stackrel{~}{e}}(C_{1,[a],\xi ,\eta ,\stackrel{~}{e}}+C_{2,[a],\xi ,\eta ,\stackrel{~}{e}})\hfill \end{array}$$
We note that $`su=X_gX_{g+1}`$ and evaluate the linear combination
$$\begin{array}{c}\underset{\beta [0,2]}{}X_{g+1}^{(2\beta )}G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}\beta }^{}\hfill \\ \begin{array}{ccc}=& 1\hfill & (2(X_gX_{g+1}+2)A_{[a],\xi ,\eta ,\stackrel{~}{e}}\hfill \\ & & +(X_gX_{g+1}+2)(X_gX_{g+1}+1)B_{[a],\xi ,\eta ,\stackrel{~}{e}})\hfill \\ +& X_{g+1}\hfill & (2A_{[a],\xi ,\eta ,\stackrel{~}{e}}+2(X_gX_{g+1}+1)B_{[a],\xi ,\eta ,\stackrel{~}{e}}\hfill \\ & & (X_gX_{g+1}+1)(C_{1,[a],\xi ,\eta ,\stackrel{~}{e}}+C_{2,[a],\xi ,\eta ,\stackrel{~}{e}}))\hfill \\ +& X_{g+1}(X_{g+1}+1)\hfill & \left(B_{[a],\xi ,\eta ,\stackrel{~}{e}}(C_{1,[a],\xi ,\eta ,\stackrel{~}{e}}+C_{2,[a],\xi ,\eta ,\stackrel{~}{e}})\right)\hfill \end{array}\hfill \\ =2(X_g+2)A_{[a],\xi ,\eta ,\stackrel{~}{e}}+(X_g+2)(X_g+1)B_{[a],\xi ,\eta ,\stackrel{~}{e}}X_{g+1}(X_g+2)(C_{1,[a],\xi ,\eta ,\stackrel{~}{e}}+C_{2,[a],\xi ,\eta ,\stackrel{~}{e}}).\hfill \end{array}$$
###### Lemma 2.11
Suppose $`g=p<k`$ and $`s=\lambda _g^{}`$, $`t=1`$. Concerning $`\stackrel{~}{e}`$, $`x_\gamma `$, $`y_\gamma `$, we let the notation be as in (2.10). For $`\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e}0)`$, we write
$$\begin{array}{ccc}\hfill V_{1,\gamma }& :=& a^\gamma \epsilon _\gamma (y_\gamma ,\eta )\hfill \\ \hfill V_{2,\gamma }& :=& a^\gamma \epsilon _\gamma (x_\gamma ,\eta )\hfill \end{array}$$
and let
$$\begin{array}{ccc}A_{[a],\xi ,\eta ,\stackrel{~}{e}}\hfill & :=& \underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e}0)}{}V_{1,\gamma }\underset{w^{}\overline{\eta }}{}(x_\gamma ,w^{})\hfill \\ & =& \underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e}0)}{}V_{2,\gamma }\underset{w^{}\overline{\eta }}{}(y_\gamma ,w^{})\hfill \\ B_{[a],\xi ,\eta ,\stackrel{~}{e}}\hfill & :=& 0\hfill \\ C_{1,[a],\xi ,\eta ,\stackrel{~}{e}}\hfill & :=& \underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e}0)}{}V_{1,\gamma }\hfill \\ C_{2,[a],\xi ,\eta ,\stackrel{~}{e}}\hfill & :=& \underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e}0)}{}V_{2,\gamma }.\hfill \end{array}$$
We obtain
$$\begin{array}{c}\hfill G_{[a],\xi ,\eta }^{}f^{}=(X_g+2)\underset{\stackrel{~}{e}\stackrel{~}{E}}{}\left(\underset{j[g+2,k]}{}X_j^{(2\stackrel{~}{e}_j)}\right)(2A_{[a],\xi ,\eta ,\stackrel{~}{e}}+(X_g+1)B_{[a],\xi ,\eta ,\stackrel{~}{e}}\\ \hfill X_{g+1}(C_{1,[a],\xi ,\eta ,\stackrel{~}{e}}+C_{2,[a],\xi ,\eta ,\stackrel{~}{e}})).\end{array}$$
Calculations similar to those of (A.2) yield the following table, to be compared to the table of (2.10).
$$\begin{array}{ccc}\hfill G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}2}^{}& =& 2(su+2)A_{[a],\xi ,\eta ,\stackrel{~}{e}}\hfill \\ \hfill G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}1}^{}& =& 2A_{[a],\xi ,\eta ,\stackrel{~}{e}}(su+1)(C_{1,[a],\xi ,\eta ,\stackrel{~}{e}}+C_{2,[a],\xi ,\eta ,\stackrel{~}{e}})\hfill \\ \hfill G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}0}^{}& =& (C_{1,[a],\xi ,\eta ,\stackrel{~}{e}}+C_{2,[a],\xi ,\eta ,\stackrel{~}{e}})\hfill \end{array}$$
###### Lemma 2.12
Suppose $`g=p<k`$ and $`s=1`$, $`t=\lambda _g^{}=\lambda _{g+1}^{}`$. Concerning $`\stackrel{~}{e}`$, $`x_\gamma `$, $`y_\gamma `$, we let the notation be as in (2.10). There exist elements $`x^{},y^{}\eta `$, $`x^{}y^{}`$, $`z\overline{\xi }`$, which we choose and fix. For $`\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e}0)`$, we write
$$U_\gamma :=a^\gamma \epsilon _\gamma (\overline{\xi }\backslash z,\eta \backslash \{x^{},y^{}\})(x_\gamma ,x^{})(y_\gamma ,y^{})$$
and let
$$\begin{array}{ccc}A_{[a],\xi ,\eta ,\stackrel{~}{e}}\hfill & :=& 0\hfill \\ B_{[a],\xi ,\eta ,\stackrel{~}{e}}\hfill & :=& \underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e}0)}{}U_\gamma (1\underset{w\overline{\xi }\backslash z}{}(w,z))\hfill \\ C_{1,[a],\xi ,\eta ,\stackrel{~}{e}}\hfill & :=& \underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e}0)}{}U_\gamma (x_\gamma ,z)\hfill \\ C_{2,[a],\xi ,\eta ,\stackrel{~}{e}}\hfill & :=& \underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e}0)}{}U_\gamma (y_\gamma ,z).\hfill \end{array}$$
We obtain
$$\begin{array}{c}\hfill G_{[a],\xi ,\eta }^{}f^{}=(X_g+2)\underset{\stackrel{~}{e}\stackrel{~}{E}}{}\left(\underset{j[g+2,k]}{}X_j^{(2\stackrel{~}{e}_j)}\right)(2A_{[a],\xi ,\eta ,\stackrel{~}{e}}+(X_g+1)B_{[a],\xi ,\eta ,\stackrel{~}{e}}\\ \hfill X_{g+1}(C_{1,[a],\xi ,\eta ,\stackrel{~}{e}}+C_{2,[a],\xi ,\eta ,\stackrel{~}{e}})).\end{array}$$
Calculations similar to those of (A.2) yield the following table.
$$\begin{array}{ccc}\hfill G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}2}^{}& =& 6B_{[a],\xi ,\eta ,\stackrel{~}{e}}\hfill \\ \hfill G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}1}^{}& =& 4B_{[a],\xi ,\eta ,\stackrel{~}{e}}2(C_{1,[a],\xi ,\eta ,\stackrel{~}{e}}+C_{2,[a],\xi ,\eta ,\stackrel{~}{e}}))\hfill \\ \hfill G_{[a],\xi ,\eta }^{}f_{\stackrel{~}{e}0}^{}& =& B_{[a],\xi ,\eta ,\stackrel{~}{e}}(C_{1,[a],\xi ,\eta ,\stackrel{~}{e}}+C_{2,[a],\xi ,\eta ,\stackrel{~}{e}})\hfill \end{array}$$
In comparison to the table of (2.10), we note that $`X_gX_{g+1}=1`$.
###### Lemma 2.13
Suppose $`g=p=k`$ and $`s1`$, $`t2`$, $`\eta =a_{g+1,[1,t]}`$. Let $`\gamma ^1`$ be the double path defined by $`\gamma _{g+1,1}^1:=1`$, $`\gamma _{g+1,2}^1:=2`$. Let $`x^{}:=a_{g+1,1}`$ and $`y^{}:=a_{g+1,2}`$, so that $`x^{},y^{}\eta `$. There exists an element $`z\overline{\xi }`$, which we choose and fix. We write
$$U:=a^{\gamma ^1}\epsilon _{\gamma ^1}(\overline{\xi }\backslash z,\eta \backslash \{x,y\})$$
and denote
$$\begin{array}{ccc}\hfill A& :=& U\underset{w^{}\overline{\eta }}{}(w^{},z)\hfill \\ \hfill B& :=& U(1\underset{w\overline{\xi }\backslash z}{}(w,z))\hfill \end{array}$$
We obtain
$$G_{[a],\xi ,\eta }^{}f^{}=\frac{1}{2}\left((X_g+1)(X_g+2)B+2(X_g+2)A\right),$$
whence the map $`f^{}`$ annihilates $`G_{[a],\xi ,\eta }^{}`$ modulo $`X_g+2`$ in case $`X_g+2`$ is odd, modulo $`(X_g+2)/2`$ in case $`X_g+2`$ is even.
Calculations similar to those of (A.2) yield
$$G_{[a],\xi ,\eta }^{}f^{}=\frac{1}{2}\left((su+1)(su+2)B+2(su+2)A\right).$$
Note that $`su=X_g`$.
###### Lemma 2.14
Suppose $`g=p=k`$ and $`s=\lambda _g^{}`$, $`t=1`$, $`\eta =a_{g+1,1}`$. For $`i[2,\lambda _{g+1}^{}]`$, we let $`\gamma ^i`$ be the double path defined by $`\gamma _{g+1,1}^i:=1`$, $`\gamma _{g+1,2}^i:=i`$. Let
$$A:=\underset{i[2,\lambda _{g+1}^{}]}{}a^{\gamma ^i}\epsilon _{\gamma ^i}$$
We obtain
$$G_{[a],\xi ,\eta }^{}f^{}=(X_g+2)A,$$
whence the map $`f^{}`$ annihilates $`G_{[a],\xi ,\eta }^{}`$ modulo $`X_g+2`$.
Calculations similar to those of (A.2) yield
$$G_{[a],\xi ,\eta }^{}f^{}=(su+2)A.$$
###### Remark 2.15
Suppose $`g=p=k`$. The application $`M^{\lambda ^{},}\text{}S^\mu `$ maps $`\{(\stackrel{ˇ}{a}_\lambda )^{}\}^{}`$ to a linear combination with coefficients $`\pm 1`$ of standard $`\mu `$-polytabloids.
See Subsection 1.1 for the definition of the $`\lambda `$-tableau $`[\stackrel{ˇ}{a}_\lambda ]`$.
We summarize.
###### Proposition 2.16 (provisional version of (2.37))
Suppose $`g<k`$. The $`\text{Z}𝒮_n`$-linear map $`M^{\lambda ^{},}\text{}S^\mu `$ factors over
where $`m^{}`$ is an integer that divides the element
$$\begin{array}{c}\hfill G_{[a],\xi ,\eta }^{}f^{}=(X_g+2)\underset{\stackrel{~}{e}\stackrel{~}{E}}{}\left(\underset{j[g+2,k]}{}X_j^{(2\stackrel{~}{e}_j)}\right)(2A_{[a],\xi ,\eta ,\stackrel{~}{e}}+(X_g+1)B_{[a],\xi ,\eta ,\stackrel{~}{e}}\\ \hfill X_{g+1}(C_{1,[a],\xi ,\eta ,\stackrel{~}{e}}+C_{2,[a],\xi ,\eta ,\stackrel{~}{e}}))\end{array}$$
of $`S^\mu `$ for any $`\lambda `$-tableau $`[a]`$ and any pair of subsets $`\xi a_g`$, $`\eta a_{g+1}`$ such that $`\mathrm{\#}\xi +\mathrm{\#}\eta =\lambda _p^{}+1`$. The set $`\stackrel{~}{E}`$ is defined in (2.10). The elements $`A_{[a],\xi ,\eta ,\stackrel{~}{e}}`$, $`B_{[a],\xi ,\eta ,\stackrel{~}{e}}`$, $`C_{1,[a],\xi ,\eta ,\stackrel{~}{e}}`$, $`C_{2,[a],\xi ,\eta ,\stackrel{~}{e}}`$ of $`S^\mu `$ are defined in (2.10, 2.11, 2.12), the definition depending on whether $`\xi `$ resp. $`\eta `$ contains $`2`$ elements.
This is the conclusion we draw from (2.10, 2.11, 2.12) concerning the modulus $`m^{}`$ and from (2.4, 2.5, 2.6, 2.7, 2.8, 2.9) concerning the vanishing of the remaining Garnir relations under $`f^{}`$. The assertion is provisional in that we have neither specified $`m^{}`$ nor commented on the possible divisibility of $`M^{\lambda ^{},}\text{}S^\mu `$ yet.
###### Proposition 2.17 (to be included in (2.37), cf. (4.9), cf. \[K 99, 4.4.3\])
Suppose $`g=k`$. The $`\text{Z}𝒮_n`$-linear map $`M^{\lambda ^{},}\text{}S^\mu `$ factors over
where $`m:=X_g+2`$ in case $`X_g`$ is odd and $`m:=(X_g+2)/2`$ in case $`X_g`$ is even. The resulting morphism $`S^\lambda \text{}S^\mu /m`$ is of order $`m`$ in $`\text{Hom}_{\text{Z}𝒮_n}(S^\lambda ,S^\mu /m)`$.
The assertions follow by (2.13, 2.14, 2.15).
We recall that in this case the double path formalism yields
$$\begin{array}{ccc}\hfill M^{\lambda ^{},}& \text{}& S^\mu \hfill \\ \hfill \{a^{}\}^{}& \text{}& \underset{1v<w\lambda _{g+1}^{}}{}(1)^{v+w}a^{v,w},\hfill \end{array}$$
where we let
$$\begin{array}{ccc}\hfill [0,\lambda _{g+1}^{}]& \text{}& [0,\mu _{g+1}^{}]\hfill \\ \hfill i& \text{}& \mathrm{\#}\left([1,i]\backslash \{v,w\}\right)\hfill \\ \hfill \mathrm{min}(\phi ^1(\{j\}))& \text{}& j\hfill \end{array}$$
in order to define the $`\mu `$-tableau $`[a^{v,w}]`$ by
$$\begin{array}{cccc}\hfill a_{j,i}^{v,w}& :=& a_{j,i}\hfill & \text{for }(j[1,\lambda _1]\backslash \{g,g+1\}\text{ and }i[1,\mu _j^{}])\text{ or }(j=g\text{ and }i[1,\lambda _g^{}])\hfill \\ \hfill a_{g,\lambda _g^{}+1}^{v,w}& :=& a_{g+1,v}\hfill & \\ \hfill a_{g,\lambda _g^{}+2}^{v,w}& :=& a_{g+1,w}\hfill & \\ \hfill a_{g+1,i}^{v,w}& :=& a_{g+1,\psi (i)}\hfill & \text{for }i[1,\mu _{g+1}^{}].\hfill \end{array}$$
### 2.3 Morphism, reduced version
To achieve that the image of our morphism from $`M^{\lambda ^{},}`$ to $`S^\mu `$ gets is not contained in some $`tS^\mu S^\mu `$, $`t2`$, we have to divide $`f^{}`$ by a factor of redundancy.
We assume $`g<k`$ throughout this section. For $`gpqk+1`$, we let $`𝒳[p,q]`$ denote the set of partial patterns, i.e. the set of intermediate sets
$$[1,2]\times (\{g,k+1\}[p,q])\mathrm{\Xi }[1,2]\times [p,q].$$
For $`p>q`$ we let $`𝒳[p,q]:=\mathrm{}`$. For $`prsq`$ and $`\mathrm{\Xi }𝒳[p,q]`$, we let $`\mathrm{\Xi }_{[r,s]}:=\mathrm{\Xi }([1,2]\times [r,s])𝒳[r,s]`$. The partial weight of $`\mathrm{\Xi }`$ is defined to be the map
$$\begin{array}{ccc}\hfill [p,q]& \text{}& [0,2]\hfill \\ \hfill r& \text{}& e_{\mathrm{\Xi },r}=\mathrm{\#}\mathrm{\Xi }_{[r,r]}.\hfill \end{array}$$
Sometimes, we shall write this map as a tuple $`e_\mathrm{\Xi }=(e_{\mathrm{\Xi },p},\mathrm{},e_{\mathrm{\Xi },q})`$. An ordered double path of weight $`e`$ is a double path $`\gamma `$ that in addition satisfies the condition
$$\overline{\gamma }(p,1)<\overline{\gamma }(p,2)\text{ for }p[g+1,k].$$
Given a pattern $`\mathrm{\Xi }𝒳[g,k+1]`$, we let $`\dot{\mathrm{\Gamma }}(\mathrm{\Xi })`$ be the set of double paths $`\gamma `$ of partial pattern $`\mathrm{\Xi }_\gamma =\mathrm{\Xi }`$, and we let $`\stackrel{}{\mathrm{\Gamma }}(\mathrm{\Xi })`$ be the set of ordered double paths $`\gamma `$ of partial pattern $`\mathrm{\Xi }_\gamma =\mathrm{\Xi }`$. We define $`\text{Z}𝒮_n`$-linear maps
$$\begin{array}{ccc}\hfill F^\lambda & \text{}& S^\mu \hfill \\ \hfill [a]& \text{}& \underset{\gamma \dot{\mathrm{\Gamma }}(\mathrm{\Xi })}{}a^\gamma \epsilon _\gamma \hfill \\ \hfill [a]& \text{}& \underset{\gamma \stackrel{}{\mathrm{\Gamma }}(\mathrm{\Xi })}{}a^\gamma \epsilon _\gamma .\hfill \end{array}$$
We allow ourselves to write e.g. for the partial pattern $`\{1\times u,\mathrm{\hspace{0.33em}2}\times u,\mathrm{\hspace{0.33em}2}\times (u+1)\}𝒳[u,u+1]`$ etc. as long as no confusion concerning $`u`$ can arise. So ‘$`+`$’ stands for ‘coordinate contained’ and ‘$``$’ stands for ‘coordinate not contained’.
Suppose given a $`\lambda `$-tableau $`[a]`$ and a double path $`\gamma `$. We let $`x_\gamma :=a_{j,\overline{\gamma }(j,1)}`$, $`j[g+2,k+1]`$ being minimal with $`1\times j\mathrm{\Xi }_\gamma `$, and $`y_\gamma :=a_{j,\overline{\gamma }(j,2)}`$, $`j[g+2,k+1]`$ being minimal with $`2\times j\mathrm{\Xi }_\gamma `$. We fix positions $`uv[1,\lambda _{g+1}^{}]`$. For a $`\lambda `$-tableau $`[a]`$ and a partial pattern $`\mathrm{\Xi }𝒳[g+2,k+1]`$, we write $`x^{}:=a_{g+1,u}`$, $`y^{}:=a_{g+1,v}`$ and let
$$\begin{array}{ccc}\hfill F^\lambda & \text{}& S^\mu \hfill \\ \hfill [a]& \text{}& \underset{\gamma \dot{\mathrm{\Gamma }}(\text{}\mathrm{\Xi })}{}a^\gamma \epsilon _\gamma (x_\gamma ,x^{})(y_\gamma ,y^{})\hfill \\ \hfill [a]& \text{}& \underset{\gamma \dot{\mathrm{\Gamma }}(\text{}\mathrm{\Xi })}{}a^\gamma \epsilon _\gamma (x_\gamma ,x^{})\hfill \\ \hfill [a]& \text{}& \underset{\gamma \dot{\mathrm{\Gamma }}(\text{}\mathrm{\Xi })}{}a^\gamma \epsilon _\gamma (y_\gamma ,y^{})\hfill \\ \hfill [a]& \text{}& \underset{\gamma \stackrel{}{\mathrm{\Gamma }}(\text{}\mathrm{\Xi })}{}a^\gamma \epsilon _\gamma (x_\gamma ,x^{})(y_\gamma ,y^{})\hfill \end{array}$$
$$\begin{array}{ccc}\hfill [a]& \text{}& \underset{\gamma \stackrel{}{\mathrm{\Gamma }}(\text{}\mathrm{\Xi })}{}a^\gamma \epsilon _\gamma (x_\gamma ,x^{})\hfill \\ \hfill [a]& \text{}& \underset{\gamma \stackrel{}{\mathrm{\Gamma }}(\text{}\mathrm{\Xi })}{}a^\gamma \epsilon _\gamma (y_\gamma ,y^{}).\hfill \end{array}$$
For instance, in this definition, it is understood that $`\text{}=\{1\times g,\mathrm{\hspace{0.33em}2}\times g\}[1,2]\times [g,g+1]`$. We note that $`A_\mathrm{\Xi }^{u,v}=A_\mathrm{\Xi }^{v,u}`$.
For $`i[0,1]`$, let $`L(i):=\{j[g+1,k1]|\lambda _{j+1}^{}=\lambda _j^{}i\}`$ be the set of $`i`$-steps. Let
$$[g,k]=L(0)L(1)\overline{L}$$
be a disjoint union. A partial pattern $`\mathrm{\Xi }𝒳[p,q]`$, $`gpqk+1`$, is called bulky in case there is some $`u[p,q1]L(0)`$ such that
$$\mathrm{\Xi }_{[u,u+1]}\{\text{},\text{},\text{},\text{},\text{},\text{},\text{}\},$$
or in case there is some $`u[p,q1]L(1)`$ such that
$$\mathrm{\Xi }_{[u,u+1]}=\text{}.$$
Let $`𝒳_{\text{nb}}[p,q]𝒳[p,q]`$ be the subset of nonbulky partial patterns in $`𝒳[p,q]`$.
###### Lemma 2.18
Let $`m^{}`$ be a number dividing the elements
$$\begin{array}{cc}\hfill (X_g+2)2& \underset{\mathrm{\Xi }𝒳[g+2,k+1]}{}(\underset{j[g+2,k]}{}X_j^{(2e_{\mathrm{\Xi },j})})A_\mathrm{\Xi }^{u,v}\hfill \\ \hfill (X_g+2)(X_g+1)& \underset{\mathrm{\Xi }𝒳[g+2,k+1]}{}(\underset{j[g+2,k]}{}X_j^{(2e_{\mathrm{\Xi },j})})A_\mathrm{\Xi }^{u,v}\hfill \\ \hfill (X_g+2)X_{g+1}& \underset{\mathrm{\Xi }𝒳[g+2,k+1]}{}(\underset{j[g+2,k]}{}X_j^{(2e_{\mathrm{\Xi },j})})(A_\mathrm{\Xi }^{u,}+A_\mathrm{\Xi }^{,u})\hfill \end{array}$$
of $`\text{Hom}_{\text{Z}𝒮_n}(F^\lambda ,S^\mu )`$ for any choice of $`uv[1,\lambda _{g+1}^{}]`$ resp. of $`u[1,\lambda _{g+1}^{}]`$. Then there is a factorization
$$(M^{\lambda ^{},}\text{}S^\mu \text{}S^\mu /m^{})=(M^{\lambda ^{},}\text{}S^\lambda \text{}S^\mu /m^{}).$$
Let $`[a]`$ be a $`\lambda `$-tableau, let $`\xi a_g`$, $`\eta a_{g+1}`$ such that $`\mathrm{\#}\xi +\mathrm{\#}\eta =\lambda _p^{}+1`$. Let $`\stackrel{~}{e}`$ be a map from $`[1,\lambda _1]\backslash \{g+1\}`$ to $`[0,2]`$ that maps $`g`$ and $`k+1`$ to $`e_g=e_{k+1}=2`$, and that maps $`j[1,\lambda _1]\backslash [g,k+1]`$ to $`e_j=0`$.
Case $`\mathrm{\#}\xi 2`$, $`\mathrm{\#}\eta 2`$. Writing $`a_{g+1,u}:=x^{}`$ and $`a_{g+1,v}:=y^{}`$, we obtain, in the notation of (2.10), including choices invoved,
$$\begin{array}{ccc}\hfill A_{[a],\xi ,\eta ,\stackrel{~}{e}}& =& \underset{\mathrm{\Xi }𝒳[g+2,k+1],e_\mathrm{\Xi }=\stackrel{~}{e}|_{[g+2,k+1]}}{}\left([a](\overline{\xi }\backslash z,\eta \backslash \{x,y\})(\underset{w^{}\overline{\eta }}{}(w^{},z))\right)A_\mathrm{\Xi }^{u,v}\hfill \\ \hfill B_{[a],\xi ,\eta ,\stackrel{~}{e}}& =& \underset{\mathrm{\Xi }𝒳[g+2,k+1],e_\mathrm{\Xi }=\stackrel{~}{e}|_{[g+2,k+1]}}{}\left([a](\overline{\xi }\backslash z,\eta \backslash \{x,y\})(1\underset{w\overline{\xi }\backslash z}{}(w,z))\right)A_\mathrm{\Xi }^{u,v}\hfill \\ \hfill C_{1,[a],\xi ,\eta ,\stackrel{~}{e}}& =& \underset{\mathrm{\Xi }𝒳[g+2,k+1],e_\mathrm{\Xi }=\stackrel{~}{e}|_{[g+2,k+1]}}{}\left([a](\overline{\xi }\backslash z,\eta \backslash \{x,y\})(x^{},z)\right)A_\mathrm{\Xi }^{,v}\hfill \\ & =& \underset{\mathrm{\Xi }𝒳[g+2,k+1],e_\mathrm{\Xi }=\stackrel{~}{e}|_{[g+2,k+1]}}{}\left([a](\overline{\xi }\backslash z,\eta \backslash \{x,y\})(y^{},z)\right)A_\mathrm{\Xi }^{,u}\hfill \\ \hfill C_{2,[a],\xi ,\eta ,\stackrel{~}{e}}& =& \underset{\mathrm{\Xi }𝒳[g+2,k+1],e_\mathrm{\Xi }=\stackrel{~}{e}|_{[g+2,k+1]}}{}\left([a](\overline{\xi }\backslash z,\eta \backslash \{x,y\})(y^{},z)\right)A_\mathrm{\Xi }^{u,},\hfill \end{array}$$
where the multitransposition $`(\overline{\xi }\backslash z,\eta \backslash \{x,y\})`$ is to be read as with respect to and independent of a choice of a bijection (cf. Section 1.2).
Case $`\mathrm{\#}\xi =\lambda _g^{}`$, $`\mathrm{\#}\eta =1`$. Writing $`\{a_{g+1,u}\}:=\eta `$, we obtain, in the notation of (2.11),
$$\begin{array}{ccc}\hfill A_{[a],\xi ,\eta ,\stackrel{~}{e}}& =& \underset{\mathrm{\Xi }𝒳[g+2,k+1],e_\mathrm{\Xi }=\stackrel{~}{e}|_{[g+2,k+1]}}{}\underset{v[1,\lambda _{g+1}^{}],a_{g+2,v}\overline{\eta }}{}\left([a]\right)A_\mathrm{\Xi }^{u,v}\hfill \\ \hfill B_{[a],\xi ,\eta ,\stackrel{~}{e}}& =& 0\hfill \\ \hfill C_{1,[a],\xi ,\eta ,\stackrel{~}{e}}& =& \underset{\mathrm{\Xi }𝒳[g+2,k+1],e_\mathrm{\Xi }=\stackrel{~}{e}|_{[g+2,k+1]}}{}\left([a]\right)A_\mathrm{\Xi }^{,u}\hfill \\ \hfill C_{2,[a],\xi ,\eta ,\stackrel{~}{e}}& =& \underset{\mathrm{\Xi }𝒳[g+2,k+1],e_\mathrm{\Xi }=\stackrel{~}{e}|_{[g+2,k+1]}}{}\left([a]\right)A_\mathrm{\Xi }^{u,}.\hfill \end{array}$$
Case $`\mathrm{\#}\xi =1`$, $`\mathrm{\#}\eta =\lambda _g^{}=\lambda _{g+1}^{}`$. Writing $`a_{g+1,u}:=x^{}`$ and $`a_{g+1,v}:=y^{}`$, we obtain, in the notation of (2.12),
$$\begin{array}{ccc}\hfill A_{[a],\xi ,\eta ,\stackrel{~}{e}}& =& 0\hfill \\ \hfill B_{[a],\xi ,\eta ,\stackrel{~}{e}}& =& \underset{\mathrm{\Xi }𝒳[g+2,k+1],e_\mathrm{\Xi }=\stackrel{~}{e}|_{[g+2,k+1]}}{}\left([a](\overline{\xi }\backslash z,\eta \backslash \{x,y\})(1\underset{w\overline{\xi }\backslash z}{}(w,z))\right)A_\mathrm{\Xi }^{u,v}\hfill \\ \hfill C_{1,[a],\xi ,\eta ,\stackrel{~}{e}}& =& \underset{\mathrm{\Xi }𝒳[g+2,k+1],e_\mathrm{\Xi }=\stackrel{~}{e}|_{[g+2,k+1]}}{}\left([a](\overline{\xi }\backslash z,\eta \backslash \{x,y\})(x^{},z)\right)A_\mathrm{\Xi }^{,v}\hfill \\ & =& \underset{\mathrm{\Xi }𝒳[g+2,k+1],e_\mathrm{\Xi }=\stackrel{~}{e}|_{[g+2,k+1]}}{}\left([a](\overline{\xi }\backslash z,\eta \backslash \{x,y\})(y^{},z)\right)A_\mathrm{\Xi }^{,u}\hfill \\ \hfill C_{2,[a],\xi ,\eta ,\stackrel{~}{e}}& =& \underset{\mathrm{\Xi }𝒳[g+2,k+1],e_\mathrm{\Xi }=\stackrel{~}{e}|_{[g+2,k+1]}}{}\left([a](\overline{\xi }\backslash z,\eta \backslash \{x,y\})(y^{},z)\right)A_\mathrm{\Xi }^{u,},\hfill \end{array}$$
where $`(\overline{\xi }\backslash z,\eta \backslash \{x,y\})`$ is to be read as with respect to and independent of a choice of a bijection.
The lemma now follows from (2.16).
###### Lemma 2.19
Let $`\mathrm{\Xi }𝒳_{\text{nb}}[g,k+1]`$. The image
$$[\stackrel{ˇ}{a}_\lambda ]\stackrel{}{f}_\mathrm{\Xi }^{\prime \prime }$$
is a linear combination of standard $`\mu `$-polytabloids with coefficients $`\pm 1`$. For a pair of different nonbulky patterns, the corresponding pair of sets of occurring standard $`\mu `$-polytabloids is disjoint.
The summand $`\pm (\stackrel{ˇ}{a}_\lambda )^\gamma `$ is a standard polytabloid after having ordered its columns downwards increasingly provided $`\mathrm{\Xi }_\gamma `$ is not bulky. Different ordered double paths yield different fillings of the columns, hence the coefficients remain $`\pm 1`$ and, moreover, the asserted disjointness holds.
For $`gpqk`$, a tuple of integers $`(\vartheta _\mathrm{\Xi })_{\mathrm{\Xi }𝒳[p,q]}`$ is called a reduced tuple of coefficients if the following conditions hold, expressed using the involution introduced in (2.2).
* For $`\mathrm{\Xi }𝒳[p,q]`$, the invariance $`\vartheta _{(\mathrm{\Xi })\iota _{q+1}}=\vartheta _\mathrm{\Xi }`$ holds. Given, in addition, $`r[p,q]`$ with $`e_{\mathrm{\Xi },r}=2`$, the invariance $`\vartheta _{(\mathrm{\Xi })\iota _{r+1}}=\vartheta _\mathrm{\Xi }`$ holds.
* In case $`\mathrm{\Xi }𝒳[p,q]`$ is bulky, the coefficient $`\vartheta _\mathrm{\Xi }`$ vanishes, in case not, it does not vanish.
We note that the tuple $`({\displaystyle \underset{j[p,q]}{}}X_j^{(2e_{\mathrm{\Xi },j})})_{\mathrm{\Xi }𝒳[g,p]}`$ satisfies $`(\text{I}_{p,q})`$ but in general not $`(\text{II}_{p,q})`$. Let it be remarked that e.g. $`(\text{II}_{p,p+1})`$ implies that $`\vartheta _{\text{}}\vartheta _{\text{}}`$ if $`p[g+1,k1]`$ and $`\lambda _p^{}=\lambda _{p+1}^{}`$, whereas $`e_{\text{}}e_{\text{}}`$.
###### Lemma 2.20
Let $`p[g+1,k]`$, let $`(\vartheta _\mathrm{\Xi })_{\mathrm{\Xi }𝒳[g,p]}`$ (resp. $`(\vartheta _\mathrm{\Xi })_{\mathrm{\Xi }𝒳[g+2,p]}`$) be a tuple of integers satisfying $`(\text{I}_{g,p})`$ (resp. $`(\text{I}_{g+2,p})`$), let $`\mathrm{\Phi }𝒳[p+1,k+1]`$. Then
$$\begin{array}{ccc}\hfill \underset{\mathrm{\Xi }𝒳[g,p]}{}\vartheta _\mathrm{\Xi }f_{\mathrm{\Xi }\mathrm{\Phi }}^{\prime \prime }& =& \underset{\mathrm{\Xi }𝒳[g,p]}{}\left(_{j[g+1,p]}e_{\mathrm{\Xi },j}!\right)\vartheta _\mathrm{\Xi }\stackrel{}{f}_{\mathrm{\Xi }\mathrm{\Phi }}^{\prime \prime }\hfill \\ \hfill \underset{\mathrm{\Xi }𝒳[g+2,p]}{}\vartheta _\mathrm{\Xi }A_{\mathrm{\Xi }\mathrm{\Phi }}^{u,v}& =& \underset{\mathrm{\Xi }𝒳[g+2,p]}{}\left(_{j[g+2,p]}e_{\mathrm{\Xi },j}!\right)\vartheta _\mathrm{\Xi }\stackrel{}{A}_{\mathrm{\Xi }\mathrm{\Phi }}^{u,v},\hfill \end{array}$$
respectively, and likewise for $`A^{u,}`$, $`A^{,v}`$ instead of $`A^{u,v}`$.
This follows, separately for each partial weight $`e_\mathrm{\Xi }`$, from (2.2, ii).
###### Lemma 2.21
Let $`pL(0)`$, $`\mathrm{\Psi }𝒳[g,p1]`$, $`\mathrm{\Phi }𝒳[p+2,k+1]`$. We obtain
$$\begin{array}{cccc}\hfill \text{(i)}& \hfill f_{\mathrm{\Psi }\text{}\mathrm{\Phi }}^{\prime \prime }2f_{\mathrm{\Psi }\text{}\mathrm{\Phi }}^{\prime \prime }2f_{\mathrm{\Psi }\text{}\mathrm{\Phi }}^{\prime \prime }+2f_{\mathrm{\Psi }\text{}\mathrm{\Phi }}^{\prime \prime }& =& 0\hfill \\ \hfill \text{(ii)}& \hfill f_{\mathrm{\Psi }\text{}\mathrm{\Phi }}^{\prime \prime }f_{\mathrm{\Psi }\text{}\mathrm{\Phi }}^{\prime \prime }f_{\mathrm{\Psi }\text{}\mathrm{\Phi }}^{\prime \prime }& =& 0\hfill \\ \hfill \text{(iii)}& \hfill f_{\mathrm{\Psi }\text{}\mathrm{\Phi }}^{\prime \prime }f_{\mathrm{\Psi }\text{}\mathrm{\Phi }}^{\prime \prime }f_{(\mathrm{\Psi })\iota _p\text{}\mathrm{\Phi }}^{\prime \prime }& =& 0\hfill \\ \hfill \text{(iii}\text{}\text{)}& \hfill f_{\mathrm{\Psi }\text{}\mathrm{\Phi }}^{\prime \prime }f_{(\mathrm{\Psi })\iota _p\text{}\mathrm{\Phi }}^{\prime \prime }f_{\mathrm{\Psi }\text{}\mathrm{\Phi }}^{\prime \prime }& =& 0\hfill \\ \hfill \text{(iv)}& \hfill f_{\mathrm{\Psi }\text{}\mathrm{\Phi }}^{\prime \prime }f_{\mathrm{\Psi }\text{}\mathrm{\Phi }}^{\prime \prime }& =& 0\hfill \\ \hfill \text{(iv}\text{}\text{)}& \hfill f_{\mathrm{\Psi }\text{}\mathrm{\Phi }}^{\prime \prime }f_{\mathrm{\Psi }\text{}\mathrm{\Phi }}^{\prime \prime }& =& 0\hfill \\ \hfill \text{(v)}& \hfill 2f_{\mathrm{\Psi }\text{}\mathrm{\Phi }}^{\prime \prime }f_{\mathrm{\Psi }\text{}\mathrm{\Phi }}^{\prime \prime }& =& 0.\hfill \end{array}$$
Equation (i) still holds in case $`pL(1)`$. Equations (i-v) continue to hold for $`pL(0)[g+2,k1]`$, $`f^{\prime \prime }`$ replaced by $`A^{u,v}`$, $`A^{u,}`$ or $`A^{,v}`$ and $`\mathrm{\Psi }𝒳[g,p1]`$ replaced by $`\mathrm{\Psi }𝒳[g+2,p1]`$.
Equation (i) ensues from the Garnir relation that arises, for a given double path $`\gamma \dot{\mathrm{\Gamma }}(\mathrm{\Xi }\text{}\mathrm{\Phi })`$, from the join of the elements in column $`p`$ in the image of $`\gamma `$ with column $`p+1`$.
We deduce (ii) from the Garnir relation that arises, for a given double path $`\gamma \dot{\mathrm{\Gamma }}(\mathrm{\Xi }\text{}\mathrm{\Phi })`$, from the join of a single element in column $`p`$ in the image of $`\gamma `$ with column $`p+1`$, using (2.2 ii) to obtain $`f_{(\mathrm{\Psi })\iota _p\text{}\mathrm{\Phi }}^{\prime \prime }=f_{(\mathrm{\Psi }\text{}\mathrm{\Phi })\iota _{p+1}}^{\prime \prime }=f_{\mathrm{\Psi }\text{}\mathrm{\Phi }}^{\prime \prime }`$.
Assertion (iii) follows using the Garnir relation that arises, for a given double path $`\gamma \dot{\mathrm{\Gamma }}(\mathrm{\Xi }\text{}\mathrm{\Phi })`$, from the join of the element in column $`p`$ in the image of $`\gamma `$ with column $`p+1`$.
Similarly, (iv) can be seen via the Garnir relation that arises, for a given double path $`\gamma \dot{\mathrm{\Gamma }}(\mathrm{\Xi }\text{}\mathrm{\Phi })`$, to the join of the element in column $`p`$ in the image of $`\gamma `$ with column $`p+1`$.
As a consequence of (i) and (ii), we obtain (v).
###### Lemma 2.22
Suppose $`g+1L(0)`$, $`\mathrm{\Phi }𝒳[g+3,k+1]`$. We obtain
$$\begin{array}{cccc}\hfill \text{(i)}& \hfill A_\text{}\mathrm{\Phi }^{u,v}2A_\text{}\mathrm{\Phi }^{u,v}2A_\text{}\mathrm{\Phi }^{u,v}+2A_\text{}\mathrm{\Phi }^{u,v}& =& 0\hfill \\ \hfill \text{(ii)}& \hfill A_\text{}\mathrm{\Phi }^{u,v}A_\text{}\mathrm{\Phi }^{u,v}A_\text{}\mathrm{\Phi }^{u,v}& =& 0\hfill \\ \hfill \text{(iii)}& \hfill A_\text{}\mathrm{\Phi }^{u,}A_\text{}\mathrm{\Phi }^{u,}A_\text{}\mathrm{\Phi }^{,u}& =& 0\hfill \\ \hfill \text{(iii}\text{}\text{)}& \hfill A_\text{}\mathrm{\Phi }^{,v}A_\text{}\mathrm{\Phi }^{v,}A_\text{}\mathrm{\Phi }^{,v}& =& 0\hfill \\ \hfill \text{(iv)}& \hfill A_\text{}\mathrm{\Phi }^{u,}A_\text{}\mathrm{\Phi }^{u,}& =& 0\hfill \\ \hfill \text{(iv}\text{}\text{)}& \hfill A_\text{}\mathrm{\Phi }^{,v}A_\text{}\mathrm{\Phi }^{,v}& =& 0\hfill \\ \hfill \text{(v)}& \hfill 2A_\text{}\mathrm{\Phi }^{u,v}A_\text{}\mathrm{\Phi }^{u,v}& =& 0.\hfill \end{array}$$
Equation (i) still holds in case $`g+1L(1)`$.
###### Lemma 2.23
Let $`pL(0)`$, let $`(\vartheta _\mathrm{\Xi })_{\mathrm{\Xi }𝒳[g,p]}`$ be a reduced tuple of coefficients, let $`\mathrm{\Phi }𝒳[p+2,k+1]`$. There is a reduced tuple of coefficients $`(\widehat{\vartheta }_{\widehat{\mathrm{\Xi }}})_{\widehat{\mathrm{\Xi }}𝒳[g,p+1]}`$ such that
$$\underset{\mathrm{\Xi }𝒳[g,p]}{}\underset{\xi 𝒳[p+1,p+1]}{}\vartheta _\mathrm{\Xi }\mathrm{\#}\xi !X_{p+1}^{(2\mathrm{\#}\xi )}\stackrel{}{f}_{\mathrm{\Xi }\xi \mathrm{\Phi }}^{\prime \prime }=\underset{\widehat{\mathrm{\Xi }}𝒳[g,p+1]}{}\widehat{\vartheta }_{\widehat{\mathrm{\Xi }}}\stackrel{}{f}_{\widehat{\mathrm{\Xi }}\mathrm{\Phi }}^{\prime \prime }$$
Fixing a partial pattern $`\stackrel{ˇ}{\mathrm{\Xi }}𝒳[g,p1]`$, $`\widehat{\vartheta }`$ can be obtained from $`\vartheta `$ letting
$$\begin{array}{ccc}\hfill \widehat{\vartheta }_{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}& =& \hfill X_p(X_p+1)\vartheta _{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}\\ \hfill \widehat{\vartheta }_{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}& =& \hfill (X_p+1)\vartheta _{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}\\ \hfill \widehat{\vartheta }_{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}& =& \hfill (X_p+1)\vartheta _{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}\\ \hfill \widehat{\vartheta }_{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}& =& \hfill (X_p1)(X_p+1)\vartheta _{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}\\ \hfill \widehat{\vartheta }_{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}& =& \hfill (X_p1)(X_p+1)\vartheta _{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}\\ \hfill \widehat{\vartheta }_{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}& =& \hfill 2\vartheta _{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}\\ \hfill \widehat{\vartheta }_{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}=\widehat{\vartheta }_{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}& =& \hfill (X_p1)\vartheta _{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}\\ \hfill \widehat{\vartheta }_{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}& =& \hfill (X_p1)X_p\vartheta _{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}.\end{array}$$
In case $`pL(0)[g+2,k+1]`$, the assertion holds for $`\stackrel{}{A}^{u,v}`$, $`\stackrel{}{A}^{u,}`$ or $`\stackrel{}{A}^{,v}`$ instead of $`\stackrel{}{f}^{\prime \prime }`$, the reduced tuples of coefficients being indexed by $`\mathrm{\Xi }𝒳[g+2,p]`$ on the left hand side and by $`\widehat{\mathrm{\Xi }}𝒳[g+2,p+1]`$ on the right hand side.
Up to a rescaling, we rewrite summands of the left hand side using (2.21), to which the roman numbers refer. To this end, we substitute $`X_{p+1}=X_p1`$ and suppose $`\stackrel{ˇ}{\mathrm{\Xi }}𝒳[g,p1]`$ be given.
$$\begin{array}{ccc}\hfill \vartheta _{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}((X_p1)X_pf_{\stackrel{ˇ}{\mathrm{\Xi }}\text{}\mathrm{\Phi }}^{\prime \prime }& & \\ \hfill +(X_p1)(f_{\stackrel{ˇ}{\mathrm{\Xi }}\text{}\mathrm{\Phi }}^{\prime \prime }+f_{\stackrel{ˇ}{\mathrm{\Xi }}\text{}\mathrm{\Phi }}^{\prime \prime })+f_{\stackrel{ˇ}{\mathrm{\Xi }}\text{}\mathrm{\Phi }}^{\prime \prime })& \stackrel{\text{(ii, v)}}{=}& \vartheta _{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}\frac{1}{2}X_p(X_p+1)f_{\stackrel{ˇ}{\mathrm{\Xi }}\text{}\mathrm{\Phi }}^{\prime \prime }\hfill \\ \hfill \vartheta _{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}\left((X_p1)X_pf_{\stackrel{ˇ}{\mathrm{\Xi }}\text{}\mathrm{\Phi }}^{\prime \prime }+(X_p1)f_{\stackrel{ˇ}{\mathrm{\Xi }}\text{}\mathrm{\Phi }}^{\prime \prime }\right)& \stackrel{\text{(vi)}}{=}& \vartheta _{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}(X_p1)(X_p+1)f_{\stackrel{ˇ}{\mathrm{\Xi }}\text{}\mathrm{\Phi }}^{\prime \prime }\hfill \\ \hfill \vartheta _{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}(\frac{1}{2}(X_p1)(f_{\stackrel{ˇ}{\mathrm{\Xi }}\text{}\mathrm{\Phi }}^{\prime \prime }+f_{(\stackrel{ˇ}{\mathrm{\Xi }})\iota _p\text{}\mathrm{\Phi }}^{\prime \prime })& & \\ \hfill +f_{\stackrel{ˇ}{\mathrm{\Xi }}\text{}\mathrm{\Phi }}^{\prime \prime })& \stackrel{\text{(iii)}}{=}& \vartheta _{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}\frac{1}{2}(X_p+1)f_{\stackrel{ˇ}{\mathrm{\Xi }}\text{}\mathrm{\Phi }}^{\prime \prime }\hfill \end{array}$$
For the last equation we made use of $`(\text{I}_{g,p})`$ in order to have $`\vartheta _{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}=\vartheta _{(\stackrel{ˇ}{\mathrm{\Xi }})\iota _p\text{}}`$.
The rewriting of the summands having in column $`p`$ proceeds by symmetry. Summands of index $`\mathrm{\Xi }=\stackrel{ˇ}{\mathrm{\Xi }}\text{}`$ remain unchanged. We pass to $`\stackrel{}{f}^{\prime \prime }`$ by inserting a factor $`2`$ where necessary. So we obtain an array of equations whose left hand sides add up to yield the left hand side of the equation we claimed. After substitution of $`\vartheta `$ by $`\widehat{\vartheta }`$ according to the table above, we sum up the corresponding right hand sides to obtain the right hand side of the equation we claimed.
The arguments are valid verbatim for $`\stackrel{}{A}^{u,v}`$, $`\stackrel{}{A}^{u,}`$, $`\stackrel{}{A}^{,v}`$ instead of $`\stackrel{}{f}^{\prime \prime }`$, except that we replace $`\stackrel{ˇ}{\mathrm{\Xi }}𝒳[g,p1]`$ by $`\stackrel{ˇ}{\mathrm{\Xi }}𝒳[g+2,p1]`$ and need to suppose that $`pg+2`$.
###### Lemma 2.24
Suppose $`g+1L(0)`$, let $`\mathrm{\Phi }𝒳[g+3,k+1]`$. We obtain
$$\begin{array}{ccc}\hfill \underset{\xi 𝒳[g+2,g+2]}{}\mathrm{\#}\xi !X_{g+2}^{(2\mathrm{\#}\xi )}\stackrel{}{A}_{\xi \mathrm{\Phi }}^{u,v}& =& X_{g+1}(X_{g+1}+1)\stackrel{}{A}_\text{}\mathrm{\Phi }^{u,v}\hfill \\ \hfill \underset{\xi 𝒳[g+2,g+2]}{}\mathrm{\#}\xi !X_{g+2}^{(2\mathrm{\#}\xi )}(\stackrel{}{A}_{\xi \mathrm{\Phi }}^{u,}+\stackrel{}{A}_{\xi \mathrm{\Phi }}^{,u})& =& (X_{g+1}+1)(\stackrel{}{A}_\text{}\mathrm{\Phi }^{u,}+\stackrel{}{A}_\text{}\mathrm{\Phi }^{,u})\hfill \\ & +& (X_{g+1}1)(X_{g+1}+1)(\stackrel{}{A}_\text{}\mathrm{\Phi }^{u,}+\stackrel{}{A}_\text{}\mathrm{\Phi }^{,u})\hfill \end{array}$$
This results from (2.22) by the argument in (2.23).
###### Lemma 2.25
Let $`pL(1)`$, let $`(\vartheta _\mathrm{\Xi })_{\mathrm{\Xi }𝒳[g,p]}`$ be a reduced tuple of coefficients, let $`\mathrm{\Phi }𝒳[p+2,k+1]`$. There is a reduced tuple of coefficients $`(\widehat{\vartheta }_{\widehat{\mathrm{\Xi }}})_{\widehat{\mathrm{\Xi }}𝒳[g,p+1]}`$ such that
$$\underset{\mathrm{\Xi }𝒳[g,p]}{}\underset{\xi 𝒳[p+1,p+1]}{}\vartheta _\mathrm{\Xi }\mathrm{\#}\xi !X_{p+1}^{(2\mathrm{\#}\xi )}\stackrel{}{f}_{\mathrm{\Xi }\xi \mathrm{\Phi }}^{\prime \prime }=\underset{\widehat{\mathrm{\Xi }}𝒳[g,p+1]}{}\widehat{\vartheta }_{\widehat{\mathrm{\Xi }}}\stackrel{}{f}_{\widehat{\mathrm{\Xi }}\mathrm{\Phi }}^{\prime \prime }$$
Fixing a partial pattern $`\stackrel{ˇ}{\mathrm{\Xi }}𝒳[g,p1]`$, $`\vartheta `$ and $`\widehat{\vartheta }`$ are related by
$$\begin{array}{ccc}\hfill \widehat{\vartheta }_{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}& =& \hfill X_p(X_p3)\vartheta _{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}\\ \hfill \widehat{\vartheta }_{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}=\widehat{\vartheta }_{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}& =& \hfill X_p(X_p2)\vartheta _{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}\\ \hfill \widehat{\vartheta }_{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}& =& \hfill 2\vartheta _{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}\\ \hfill \widehat{\vartheta }_{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}=\widehat{\vartheta }_{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}& =& \hfill (X_p2)\vartheta _{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}\\ \hfill \widehat{\vartheta }_{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}& =& \hfill (X_p2)(X_p1)\vartheta _{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}\\ \hfill \widehat{\vartheta }_{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}& =& \hfill 2\vartheta _{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}\end{array}$$
$$\begin{array}{ccc}\hfill \widehat{\vartheta }_{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}=\widehat{\vartheta }_{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}& =& \hfill (X_p2)\vartheta _{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}\\ \hfill \widehat{\vartheta }_{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}& =& \hfill (X_p2)(X_p1)\vartheta _{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}\\ \hfill \widehat{\vartheta }_{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}& =& \hfill 2\vartheta _{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}\\ \hfill \widehat{\vartheta }_{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}=\widehat{\vartheta }_{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}& =& \hfill (X_p2)\vartheta _{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}\\ \hfill \widehat{\vartheta }_{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}& =& \hfill (X_p2)(X_p1)\vartheta _{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}\end{array}$$
In case $`pL(1)[g+2,k+1]`$, the assertion holds for $`\stackrel{}{A}^{u,v}`$, $`\stackrel{}{A}^{u,}`$ or $`\stackrel{}{A}^{,v}`$ instead of $`f^{\prime \prime }`$, the reduced tuples of coefficients being indexed by $`\mathrm{\Xi }𝒳[g+2,p]`$ on the left and by $`\widehat{\mathrm{\Xi }}𝒳[g+2,p+1]`$ on the right hand side.
We rewrite summands of the left hand side using (2.21 i). To this end, we substitute $`X_{p+1}=X_p2`$ and suppose $`\stackrel{ˇ}{\mathrm{\Xi }}𝒳[g,p1]`$ be given.
$$\begin{array}{ccc}\hfill \vartheta _{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}((X_p2)(X_p1)f_{\stackrel{ˇ}{\mathrm{\Xi }}\text{}\mathrm{\Phi }}^{\prime \prime }& & \\ \hfill +(X_p2)(f_{\stackrel{ˇ}{\mathrm{\Xi }}\text{}\mathrm{\Phi }}^{\prime \prime }+f_{\stackrel{ˇ}{\mathrm{\Xi }}\text{}\mathrm{\Phi }}^{\prime \prime })+f_{\stackrel{ˇ}{\mathrm{\Xi }}\text{}\mathrm{\Phi }}^{\prime \prime })& \stackrel{\text{(i)}}{=}& \vartheta _{\stackrel{ˇ}{\mathrm{\Xi }}\text{}}((X_p2)X_p(f_{\stackrel{ˇ}{\mathrm{\Xi }}\text{}\mathrm{\Phi }}^{\prime \prime }+f_{\stackrel{ˇ}{\mathrm{\Xi }}\text{}\mathrm{\Phi }}^{\prime \prime })\hfill \\ & & \frac{1}{2}(X_p3)X_pf_{\stackrel{ˇ}{\mathrm{\Xi }}\text{}\mathrm{\Phi }}^{\prime \prime })\hfill \end{array}$$
All the other summands remain unchanged.
###### Lemma 2.26
Suppose $`g+1L(1)`$, let $`\mathrm{\Phi }𝒳[g+3,k+1]`$. We obtain
$$\underset{\xi 𝒳[g+2,g+2]}{}\mathrm{\#}\xi !X_{g+2}^{(2\mathrm{\#}\xi )}\stackrel{}{A}_{\xi \mathrm{\Phi }}^{u,v}=X_{g+1}(X_{g+1}3)\stackrel{}{A}_\text{}\mathrm{\Phi }^{u,v}+X_{g+1}(X_{g+1}2)(\stackrel{}{A}_\text{}\mathrm{\Phi }^{u,v}+\stackrel{}{A}_\text{}\mathrm{\Phi }^{u,v}).$$
This results from (2.22), cf. (2.25).
Henceforth, we denote by $`\vartheta `$ the tuple of tuples $`\vartheta =((\vartheta _\mathrm{\Xi })_{\mathrm{\Xi }𝒳[g,r]})_{r[g,k]}`$ of integers such that the entry $`\vartheta _{\text{}}=1`$ furnishes the tuple in case $`r=g`$, such that each entry $`(\vartheta _\mathrm{\Xi })_{\mathrm{\Xi }𝒳[g,r]}`$ is a reduced tuple of coefficients, such that for $`r[g,k1]`$ the tuples $`(\vartheta _\mathrm{\Xi })_{\mathrm{\Xi }𝒳[g,r]}`$ and $`(\vartheta _\mathrm{\Xi })_{\mathrm{\Xi }𝒳[g,r+1]}`$ are related by the formulas in (2.23) (resp. 2.25) in case $`rL(0)`$ (resp. $`rL(1)`$), and such that in the remaining case $`r\overline{L}`$ the equation
$$\vartheta _{\mathrm{\Xi }\xi }=\mathrm{\#}\xi !X_{r+1}^{(2\mathrm{\#}\xi )}\vartheta _\mathrm{\Xi }$$
holds ‘unchanged’ for $`\mathrm{\Xi }𝒳[g,r]`$, $`\xi 𝒳[r+1,r+1]`$. We reindex the tuple for $`r=k`$ by $`\vartheta _\mathrm{\Xi }:=\vartheta _{\mathrm{\Xi }_{[g,k]}}`$ for $`\mathrm{\Xi }𝒳[g,k+1]`$. Furthermore, for $`r,s[g+1,k]`$, we let
$$R_{[r,s]}:=\underset{i[0,1]}{}\underset{jL(i)[r,s]}{}X_j^{(2i)},$$
in which a product is to be read to be equal to $`1`$ if no factor is present. For $`r\text{Z}`$, we write
$$\dot{r}:=\mathrm{gcd}(r,2).$$
###### Lemma 2.27
Let $`g+1rqk`$ such that $`r1\overline{L}`$, let $`\mathrm{\Psi }𝒳[g,r1]`$, $`\mathrm{\Xi }𝒳[r,q]`$ such that $`\mathrm{\Psi }\mathrm{\Xi }`$ is nonbulky. Then $`\vartheta _\mathrm{\Psi }R_{[r,q1]}`$ divides $`\vartheta _{\mathrm{\Psi }\mathrm{\Xi }}`$. If, moreover, $`\mathrm{\Xi }_{[r,r]}=\text{}`$, then $`2\vartheta _\mathrm{\Xi }R_{[r,q1]}`$ divides $`\vartheta _{\mathrm{\Psi }\mathrm{\Xi }}`$.
We perform an induction on $`q`$ to prove the assertion that $`\vartheta _\mathrm{\Psi }R_{[r,q1]}`$ divides $`\vartheta _{\mathrm{\Psi }\mathrm{\Xi }}`$ and that the quotient $`\vartheta _{\mathrm{\Psi }\mathrm{\Xi }}/(\vartheta _\mathrm{\Psi }R_{[r,q1]})`$ is divisible by $`X_q^{(2e_{\mathrm{\Xi },q})}`$. The factor $`2`$ in the second assertion then follows from the start of the induction.
At $`q=r`$, we have $`R_{[r,r1]}=1`$, and the quotient $`\vartheta _{\mathrm{\Psi }\mathrm{\Xi }}/(\vartheta _\mathrm{\Psi }R_{[r,q1]})`$ equals $`\mathrm{\#}\mathrm{\Xi }!X_r^{(2\mathrm{\#}\mathrm{\Xi })}`$ since $`r1\overline{L}`$.
We assume our assertion to hold for $`q1r`$. Let $`\stackrel{ˇ}{\mathrm{\Xi }}:=\mathrm{\Xi }_{[r,q1]}`$.
Case $`q1\overline{L}`$. We have $`\vartheta _{\mathrm{\Psi }\mathrm{\Xi }}/\vartheta _{\mathrm{\Psi }\stackrel{ˇ}{\mathrm{\Xi }}}=\mathrm{\#}\mathrm{\Xi }_{[q,q]}!X_q^{(2e_{\mathrm{\Xi },q})}`$ and $`R_{[r,q1]}=R_{[r,q2]}`$.
Case $`q1L(0)`$. We have $`X_q=X_{q1}1`$ and $`R_{[r,q1]}=X_{q1}(X_{q1}+1)R_{[r,q2]}`$. By assumption, $`\vartheta _{\mathrm{\Psi }\stackrel{ˇ}{\mathrm{\Xi }}}/(R_{[r,q2]}\vartheta _\mathrm{\Psi })`$ is integral and divisible by $`X_{q1}^{(2e_{\mathrm{\Xi },q1})}`$. We need to see that
$$(\vartheta _{\mathrm{\Psi }\mathrm{\Xi }}/\vartheta _{\mathrm{\Psi }\stackrel{ˇ}{\mathrm{\Xi }}})(\vartheta _{\mathrm{\Psi }\stackrel{ˇ}{\mathrm{\Xi }}}/(R_{[r,q2]}\vartheta _\mathrm{\Psi }))$$
is divisible by $`X_{q1}(X_{q1}+1)X_q^{(2e_{\mathrm{\Xi },q})}`$. The left hand side factor of that expression taken from the table in (2.23), the assertion follows separately for each nonbulky subset $`\mathrm{\Xi }_{[q1,q]}`$.
Case $`q1L(1)`$. We have $`X_q=X_{q1}2`$ and $`R_{[r,q1]}=X_{q1}R_{[r,q2]}`$. By assumption, $`\vartheta _{\mathrm{\Psi }\stackrel{ˇ}{\mathrm{\Xi }}}/(R_{[r,q2]}\vartheta _\mathrm{\Psi })`$ is integral and divisible by $`X_{q1}^{(2e_{\mathrm{\Xi },q1})}`$. We need to see that
$$(\vartheta _{\mathrm{\Psi }\mathrm{\Xi }}/\vartheta _{\mathrm{\Psi }\stackrel{ˇ}{\mathrm{\Xi }}})(\vartheta _{\mathrm{\Psi }\stackrel{ˇ}{\mathrm{\Xi }}}/(R_{[r,q2]}\vartheta _\mathrm{\Psi }))$$
is divisible by $`X_{q1}X_q^{(2e_{\mathrm{\Xi },q})}`$. The left hand side factor of that expression taken from the table in (2.25), the assertion follows separately for each nonbulky subset $`\mathrm{\Xi }_{[q1,q]}`$.
###### Lemma 2.28
Let $`g+1rqk`$ such that $`r1\overline{L}`$ and $`[r,q1]L(1)`$, let $`\mathrm{\Psi }𝒳_{\text{nb}}[g,r1]`$. As ideal of coefficients, we obtain
$$\left(\vartheta _{\mathrm{\Psi }\mathrm{\Xi }}|\mathrm{\Xi }𝒳[r,q]\right)=\left(\vartheta _\mathrm{\Psi }R_{[r,q1]}\dot{X}_q\right)\text{Z}.$$
We shall perform the following calculations.
* For $`\mathrm{\Xi }𝒳_{\text{nb}}[r,q]`$ of partial weight $`e_\mathrm{\Xi }=(1,1,\mathrm{},1)`$, we have $`\vartheta _{\mathrm{\Psi }\mathrm{\Xi }}=\vartheta _\mathrm{\Psi }R_{[r,q1]}X_q`$.
* For $`\mathrm{\Xi }𝒳_{\text{nb}}[r,q]`$ of partial weight $`e_\mathrm{\Xi }=(1,1,\mathrm{},1,2)`$, we have $`\vartheta _{\mathrm{\Psi }\mathrm{\Xi }}=\vartheta _\mathrm{\Psi }R_{[r,q1]}2`$.
* Conversely, $`\vartheta _\mathrm{\Psi }R_{[r,q1]}2`$ or $`\vartheta _\mathrm{\Psi }R_{[r,q1]}X_q`$ divides $`\vartheta _{\mathrm{\Psi }\mathrm{\Xi }}`$ for any $`\mathrm{\Xi }𝒳_{\text{nb}}[r,q]`$.
Ad (a). We perform an induction on $`q`$, starting with $`q=r`$, $`R_{[r,r1]}=1`$ and $`\vartheta _\mathrm{\Psi }\text{}=\vartheta _\mathrm{\Psi }\text{}=\vartheta _\mathrm{\Psi }X_q`$. Assume the assertion known for $`\stackrel{ˇ}{\mathrm{\Xi }}:=\mathrm{\Xi }_{[r,q1]}`$. We calculate
$$\begin{array}{ccc}\hfill \vartheta _{\mathrm{\Psi }\stackrel{ˇ}{\mathrm{\Xi }}\text{}}& \stackrel{\text{(}\text{2.25}\text{)}}{=}& \vartheta _{\mathrm{\Psi }\stackrel{ˇ}{\mathrm{\Xi }}}(X_{q1}2)\hfill \\ & =& \vartheta _\mathrm{\Psi }R_{[r,q2]}X_{q1}(X_{q1}2)\hfill \\ & =& \vartheta _\mathrm{\Psi }R_{[r,q1]}X_q,\hfill \end{array}$$
dito for $`\xi \backslash \widehat{\mathrm{\Xi }}=\text{}`$
Ad (b). The argument for (a) can be applied, except that we obtain a factor $`2`$ instead of a factor $`(X_{q1}2)=X_q`$ in the last step.
Ad (c). We compare to the proof of (2.27). We need to show that there exists a step in the induction in which the quotient $`\vartheta _{\mathrm{\Psi }\mathrm{\Xi }}/(R_{[r,q1]}\vartheta _\mathrm{\Psi })`$ is divisible by $`2`$ or by $`X_q`$. In case $`e_{\mathrm{\Xi },r}=2`$, a factor $`2`$ enters at $`r`$ already, so that we may assume $`e_{\mathrm{\Xi },r}1`$. In case $`e_{\mathrm{\Xi },q}1`$, a factor $`X_q`$ enters at $`q`$, so that we may assume $`e_{\mathrm{\Xi },q}=2`$. Since $`\mathrm{\Xi }`$ is not bulky, necessarily or occurs as a subpattern of $`\mathrm{\Xi }`$. But then a factor $`2`$ enters at the induction step that adjoins the right column of that subpattern (2.25).
###### Lemma 2.29
Let $`g+1rqk`$ such that $`r1\overline{L}`$, $`[r,q1]L(0)L(1)`$ and $`[r,q1]L(0)\mathrm{}`$, let $`\mathrm{\Psi }𝒳_{\text{nb}}[g,r1]`$. As ideal of coefficients, we obtain
$$\left(\vartheta _{\mathrm{\Psi }\mathrm{\Xi }}|\mathrm{\Xi }𝒳[r,q]\right)=\left(\vartheta _\mathrm{\Psi }R_{[r,q1]}\right)\text{Z}.$$
Let $`s:=\mathrm{max}(L(0)[r,q1])+1[r+1,q]`$. We shall show the following assertions.
* For $`t[s,q]`$ and a partial pattern $`\mathrm{\Xi }(q,t)𝒳_{\text{nb}}[r,q]`$ of partial weight
$$e_{\mathrm{\Xi }(q,t)}=\{\begin{array}{cc}2\hfill & \text{ for }j[s,t]\hfill \\ 1\hfill & \text{ for }j[r,q]\backslash [s,t],\hfill \end{array}$$
which is a slight abuse of notation since $`\mathrm{\Xi }(q,t)`$ does not only depend on $`q`$ and $`t`$, we have
$$\vartheta _{\mathrm{\Psi }\mathrm{\Xi }(q,t)}=\{\begin{array}{cc}\pm \vartheta _\mathrm{\Psi }R_{[r,q1]}\left(\underset{j[s,t1]}{}(X_j3)\right)(X_{q1}2)\hfill & \text{ for }t[s,q1]\hfill \\ \pm \vartheta _\mathrm{\Psi }R_{[r,q1]}\left(\underset{j[s,q1]}{}(X_j3)\right)\hfill & \text{ for }t=q,\hfill \end{array}$$
the sign $`\pm `$ indicating that the equation holds up to sign.
* In case $`s[r+1,q2]`$, for $`\mathrm{\Xi }𝒳_{\text{nb}}[r,q]`$ of partial weight
$$e_{\mathrm{\Xi },j}=\{\begin{array}{cc}2\hfill & \text{ for }j\{s,q\}\hfill \\ 1\hfill & \text{ for }j[r,q]\backslash \{s,q\},\hfill \end{array}$$
we have
$$\vartheta _{\mathrm{\Psi }\mathrm{\Xi }}=\vartheta _\mathrm{\Psi }R_{[r,q1]}2.$$
* Assertions (a) and (b) imply the lemma.
Ad (c). The inclusion $``$ follows from (2.27). In case $`s=q`$, the lemma follows from $`\vartheta _{\mathrm{\Psi }\mathrm{\Xi }(q,q)}=\vartheta _\mathrm{\Psi }R_{[r,q1]}`$. In case $`s=q1`$, the lemma follows from
$$\begin{array}{ccc}\hfill \vartheta _{\mathrm{\Psi }\mathrm{\Xi }(q,q1)}& =& \pm \vartheta _\mathrm{\Psi }R_{[r,q1]}(X_{q1}2)\hfill \\ \hfill \vartheta _{\mathrm{\Psi }\mathrm{\Xi }(q,q)}& =& \pm \vartheta _\mathrm{\Psi }R_{[r,q1]}(X_{q1}3).\hfill \end{array}$$
In case $`s[r+1,q2]`$, it suffices to show that
$$𝔞:=\left(2,\underset{j[s,q1]}{}(X_j3),\left(\underset{j[s,t1]}{}(X_j3)\right)(X_{q1}2)|t[s,q1]\right)=(1).$$
We perform an induction on $`t`$ from $`t=q1`$ downwards to $`t=s`$, contending that $`{\displaystyle \underset{j[s,t1]}{}}(X_j3)𝔞`$. In fact,
$$\left(\underset{j[s,t]}{}(X_j3)\right)\left(\underset{j[s,t1]}{}(X_j3)\right)(X_{q1}2)=\left(\underset{j[s,t1]}{}(X_j3)\right)(X_tX_{q1}1)$$
and $`X_tX_{q1}=2(tq+1)`$.
Ad (a). Let $`\overline{q}[r,q]`$. For $`t\overline{q}+1`$, let $`\mathrm{\Xi }(\overline{q},t)𝒳_{\text{nb}}[r,\overline{q}]`$ denote a partial pattern of the same partial weight as $`\mathrm{\Xi }(\overline{q},\overline{q})`$. We perform an induction on $`\overline{q}`$ running from $`r`$ to $`q`$ to prove the assertion that
$$\vartheta _{\mathrm{\Psi }\mathrm{\Xi }(\overline{q},t)}=\{\begin{array}{cc}\vartheta _\mathrm{\Psi }R_{[r,\overline{q}1]}X_{\overline{q}}\hfill & \text{ for }\overline{q}[r,s1]\hfill \\ \pm \vartheta _\mathrm{\Psi }R_{[r,\overline{q}1]}\left(\underset{j[s,\overline{q}1]}{}(X_j3)\right)\hfill & \text{ for }\overline{q}[s,t]\hfill \\ \pm \vartheta _\mathrm{\Psi }R_{[r,\overline{q}1]}\left(\underset{j[s,t1]}{}(X_j3)\right)(X_{\overline{q}1}2)\hfill & \text{ for }\overline{q}[t+1,q].\hfill \end{array}$$
To begin with, we recall our situation, viz. $`r<r+1st`$ and $`r\overline{q}`$.
For $`\overline{q}=r`$, we have $`R_{[r,r1]}=1`$ and $`e_{\mathrm{\Xi }(\overline{q},t),r}=1`$, whence we obtain
$$\vartheta _{\mathrm{\Psi }\mathrm{\Xi }(\overline{q},t)}=\vartheta _\mathrm{\Psi }R_{[r,\overline{q}1]}X_{\overline{q}}$$
since $`r1\overline{L}`$.
Assume the assertion to be true for $`r\overline{q}1s2`$.
Case $`\overline{q}1L(0)`$. We calculate
$$\begin{array}{ccc}\hfill \vartheta _{\mathrm{\Psi }\mathrm{\Xi }(\overline{q},t)}& \stackrel{(\text{2.23})}{=}& \vartheta _{\mathrm{\Psi }\mathrm{\Xi }(\overline{q}1,t)}(X_{\overline{q}1}1)(X_{\overline{q}1}+1)\hfill \\ & =& \vartheta _\mathrm{\Psi }R_{[r,\overline{q}2]}X_{\overline{q}1}(X_{\overline{q}1}1)(X_{\overline{q}1}+1)\hfill \\ & =& \vartheta _\mathrm{\Psi }R_{[r,\overline{q}1]}X_{\overline{q}}.\hfill \end{array}$$
Case $`\overline{q}1L(1)`$. We calculate
$$\begin{array}{ccc}\hfill \vartheta _{\mathrm{\Psi }\mathrm{\Xi }(\overline{q},s)}& \stackrel{(\text{2.25})}{=}& \vartheta _{\mathrm{\Psi }\mathrm{\Xi }(\overline{q}1,t)}(X_{\overline{q}1}2)\hfill \\ & =& \vartheta _\mathrm{\Psi }R_{[r,\overline{q}2]}X_{\overline{q}1}(X_{\overline{q}1}2)\hfill \\ & =& \vartheta _\mathrm{\Psi }R_{[r,\overline{q}1]}X_{\overline{q}}.\hfill \end{array}$$
Assume the assertion to be true for $`\overline{q}1=s1`$. Note that $`\overline{q}1L(0)`$, so
$$\begin{array}{ccc}\hfill \vartheta _{\mathrm{\Psi }\mathrm{\Xi }(\overline{q},t)}& \stackrel{(\text{2.23})}{=}& \vartheta _{\mathrm{\Psi }\mathrm{\Xi }(\overline{q}1,t)}(X_{\overline{q}1}+1)\hfill \\ & =& \vartheta _\mathrm{\Psi }R_{[r,\overline{q}2]}X_{\overline{q}1}(X_{\overline{q}1}+1)\hfill \\ & =& \vartheta _\mathrm{\Psi }R_{[r,\overline{q}1]}.\hfill \end{array}$$
Assume the assertion to be true for $`s\overline{q}1t1`$. Note that $`\overline{q}1L(1)`$, so
$$\begin{array}{ccc}\hfill \vartheta _{\mathrm{\Psi }\mathrm{\Xi }(\overline{q},t)}& \stackrel{(\text{2.25})}{=}& \vartheta _{\mathrm{\Psi }\mathrm{\Xi }(\overline{q}1,t)}X_{\overline{q}1}(X_{\overline{q}1}3)\hfill \\ & =& \pm \vartheta _\mathrm{\Psi }R_{[r,\overline{q}2]}\left(\underset{j[s,\overline{q}2]}{}(X_j3)\right)X_{\overline{q}1}(X_{\overline{q}1}3)\hfill \\ & =& \pm \vartheta _\mathrm{\Psi }R_{[r,\overline{q}1]}\left(\underset{j[s,\overline{q}1]}{}(X_j3)\right).\hfill \end{array}$$
Assume the assertion to be true for $`\overline{q}1=t`$. Note that $`\overline{q}1L(1)`$, so
$$\begin{array}{ccc}\hfill \vartheta _{\mathrm{\Psi }\mathrm{\Xi }(\overline{q},t)}& \stackrel{(\text{2.25})}{=}& \vartheta _{\mathrm{\Psi }\mathrm{\Xi }(\overline{q}1,t)}X_{\overline{q}1}(X_{\overline{q}1}2)\hfill \\ & =& \pm \vartheta _\mathrm{\Psi }R_{[r,\overline{q}2]}\left(\underset{j[s,\overline{q}2]}{}(X_j3)\right)X_{\overline{q}1}(X_{\overline{q}1}2)\hfill \\ & =& \pm \vartheta _\mathrm{\Psi }R_{[r,\overline{q}1]}\left(\underset{j[s,t1]}{}(X_j3)\right)(X_{\overline{q}1}2).\hfill \end{array}$$
Assume the assertion to be true for $`t+1\overline{q}1`$. Note that $`\overline{q}1L(1)`$, so
$$\begin{array}{ccc}\hfill \vartheta _{\mathrm{\Psi }\mathrm{\Xi }(\overline{q},t)}& \stackrel{(\text{2.25})}{=}& \vartheta _{\mathrm{\Psi }\mathrm{\Xi }(\overline{q}1,t)}(X_{\overline{q}1}2)\hfill \\ & =& \pm \vartheta _\mathrm{\Psi }R_{[r,\overline{q}2]}\left(\underset{j[s,t1]}{}(X_j3)\right)(X_{\overline{q}2}2)(X_{\overline{q}1}2)\hfill \\ & =& \pm \vartheta _\mathrm{\Psi }R_{[r,\overline{q}1]}\left(\underset{j[s,t1]}{}(X_j3)\right)(X_{\overline{q}1}2).\hfill \end{array}$$
Ad (b). We modify the last step in the proof of the formula for $`\vartheta _{\mathrm{\Psi }\mathrm{\Xi }(q,s)}`$ in (a) according to (2.25), i.e. we replace the factor $`(X_{q1}2)`$ by the factor $`2`$.
A subinterval $`[r,s]`$ of $`[g+1,k]`$ is called connected if $`[r,s1]L(0)L(1)`$. A maximal connected subinterval with respect to the ordering given by inclusion is called a component. We write $`[g+1,k]`$ as a disjoint decomposition into components
$$[g+1,k]=\underset{\kappa [1,K]}{}[p(\kappa ),q(\kappa )],$$
where $`p(\kappa )q(\kappa )`$ for $`\kappa [1,K]`$, $`p(1)=g+1`$, $`q(\kappa )+1=p(\kappa +1)`$ for $`\kappa [1,K1]`$ and $`q(K)=k`$. So $`\{q(\kappa )|\kappa [1,K]\}=\overline{L}\backslash \{g\}`$.
###### Lemma 2.30
The ideal of coefficients may be factored as
$$\left(\vartheta _\mathrm{\Xi }|\mathrm{\Xi }𝒳[g,k+1]\right)=\underset{\kappa [1,K]}{}𝔱_\kappa ,$$
where for $`\kappa [1,K]`$ we denote
$$𝔱_\kappa :=\left(\vartheta _{\mathrm{\Psi }\mathrm{\Xi }}/\vartheta _\mathrm{\Psi }|\mathrm{\Xi }𝒳[p(\kappa ),q(\kappa )]\right)\text{Z}$$
for some choice of a partial pattern $`\mathrm{\Psi }𝒳_{\text{nb}}[g,p(\kappa )1]`$, a choice which does not affect these quotients.
This factorization follows by the independence of choice just mentioned.
###### Proposition 2.31
Let the factor of redundancy be
$$R:=R_{[g+1,k]}\left(\underset{\stackrel{\kappa [1,K],}{[p(\kappa ),q(\kappa )1]L(1)}}{}\dot{X}_{q(\kappa )}\right).$$
The ideal of coefficients is given by
$$\left(\vartheta _\mathrm{\Xi }|\mathrm{\Xi }𝒳[g,k+1]\right)=\left(R\right)\text{Z}.$$
In particular, the $`\text{Z}𝒮_n`$-linear map
$$M^{\lambda ^{},}\text{}S^\mu $$
is divisible exactly by $`R`$, i.e. any matrix representing $`f^{}`$ linearly over Z has $`R`$ as greatest common divisor of its entries.
The assertion on the ideal of coefficients ensues from (2.28, 2.29, 2.30). We conclude by (2.23, 2.25) that the linear combination
$$f^{\prime \prime }=\underset{\mathrm{\Xi }𝒳[g,k+1]}{}\left(\underset{j[g+1,k]}{}(e_{\mathrm{\Xi },j}!X_j^{(2e_{\mathrm{\Xi },j})})\right)\stackrel{}{f}_\mathrm{\Xi }^{\prime \prime }:F^\lambda \text{}S^\mu ,$$
cf. (2.20), can be rewritten as
$$f^{\prime \prime }=\underset{\mathrm{\Xi }𝒳[g,k+1]}{}\vartheta _\mathrm{\Xi }\stackrel{}{f}_\mathrm{\Xi }^{\prime \prime },$$
whose coefficients are divisible by $`R`$ and vanish for $`\mathrm{\Xi }`$ bulky. Applying $`f^{\prime \prime }`$ to the tableau $`[\stackrel{ˇ}{a}_\lambda ]`$, we see by (2.19) that $`f^{\prime \prime }`$ is divisible exactly by $`R`$. The assertion remains true for $`f^{}`$, being the factorization of $`f^{\prime \prime }`$ over the epimorphism $`F^\lambda \text{}M^{\lambda ^{},}`$.
The modulus remains to be taken under consideration.
###### Lemma 2.32
We may rewrite
$$\begin{array}{ccc}\hfill 2\underset{\mathrm{\Xi }𝒳[g+2,k+1]}{}\left(\underset{j[g+2,k]}{}(e_{\mathrm{\Xi },j}!X_j^{(2e_{\mathrm{\Xi },j})})\right)\stackrel{}{A}_\mathrm{\Xi }^{u,v}& =& \underset{\mathrm{\Xi }𝒳[g+2,k+1]}{}\vartheta _\text{}\mathrm{\Xi }\stackrel{}{A}_\mathrm{\Xi }^{u,v}\hfill \\ \hfill X_{g+1}\underset{\mathrm{\Xi }𝒳[g+2,k+1]}{}\left(\underset{j[g+2,k]}{}(e_{\mathrm{\Xi },j}!X_j^{(2e_{\mathrm{\Xi },j})})\right)(\stackrel{}{A}_\mathrm{\Xi }^{u,}+\stackrel{}{A}_\mathrm{\Xi }^{,u})& & \\ \hfill =\underset{\mathrm{\Xi }𝒳[g+2,k+1]}{}\left(\vartheta _\text{}\mathrm{\Xi }\stackrel{}{A}_\mathrm{\Xi }^{u,}+\vartheta _\text{}\mathrm{\Xi }\stackrel{}{A}_\mathrm{\Xi }^{,u}\right)& & \end{array}$$
This ensues from (2.24, 2.23, 2.26, 2.25).
###### Lemma 2.33
Let $`g+1qk`$, let $`\xi =\text{}`$ or $`\xi =\text{}`$ in $`𝒳[g+1,g+1]`$. Suppose $`[g+1,q1]L(1)`$. As ideal of coefficients, we obtain
$$\left(\vartheta _\mathrm{\Xi }|\mathrm{\Xi }𝒳[g,q],\mathrm{\Xi }_{[g+1,g+1]}=\xi \right)=\left(R_{[g+1,q1]}\dot{X}_q\right).$$
Consequently, by comparison with (2.31),
$$\left(\vartheta _\mathrm{\Xi }|\mathrm{\Xi }𝒳[g,p+1],\mathrm{\Xi }_{[g+1,g+1]}=\xi \right)=\left(R\right).$$
Consider the partial weights indexing generators of the ideal of coefficients having no restriction imposed on column $`g+1`$. According to assertions (a) and (b) of the proof of (2.28), they can be chosen to have column $`g+1`$ equal to $`\xi `$.
###### Lemma 2.34
Let $`g+1qk`$, let $`\xi =\text{}`$ or $`\xi =\text{}`$ in $`𝒳[g+1,g+1]`$. Suppose $`[g+1,q1]L(0)L(1)`$ such that $`[g+1,q1]L(0)\mathrm{}`$. As ideal of coefficients we obtain
$$\left(\vartheta _\mathrm{\Xi }|\mathrm{\Xi }𝒳[g,q],\mathrm{\Xi }_{[g+1,g+1]}=\xi \right)=\left(R_{[g+1,q1]}\right).$$
Whence, by comparison with (2.31), we have
$$\left(\vartheta _\mathrm{\Xi }|\mathrm{\Xi }𝒳[g,p+1],\mathrm{\Xi }_{[g+1,g+1]}=\xi \right)=\left(R\right).$$
Consider the partial weights indexing generators of the ideal of coefficients having no restriction imposed on column $`g+1`$. According to assertions (a) and (b) of the proof of (2.29), they can be chosen to have column $`g+1`$ equal to $`\xi `$.
###### Lemma 2.35
Let $`g+1qk`$. Suppose $`[g+1,q1]L(1)`$. As ideal of coefficients, we obtain
$$\left(\vartheta _\mathrm{\Xi }|\mathrm{\Xi }𝒳[g,q],\mathrm{\Xi }_{[g+1,g+1]}=\text{}\right)=\left(2R_{[g+1,q1]}\right).$$
Therefore, by comparison with (2.31),
$$\left(\vartheta _\mathrm{\Xi }|\mathrm{\Xi }𝒳[g,k+1],\mathrm{\Xi }_{[g+1,g+1]}=\text{}\right)=\left(2R/\dot{X}_{q(1)}\right).$$
Note that $`\dot{X}_{q(1)}=\dot{X}_{g+1}`$.
We claim the following assertions.
* For $`t[g+1,q]`$, we let $`\mathrm{\Xi }(g,t)𝒳_{\text{nb}}[g,q]`$ denote a partial pattern of partial weight
$$e_{\mathrm{\Xi }(q,t),j}=\{\begin{array}{cc}2\hfill & \text{ for }j[g,t]\hfill \\ 1\hfill & \text{ for }j[t+1,q]\hfill \end{array}$$
and obtain
$$\vartheta _{\mathrm{\Xi }(q,t)}=\{\begin{array}{cc}\pm 2R_{[g+1,q1]}\left(\underset{j[g+1,t1]}{}(X_j3)\right)(X_{q1}2)\hfill & \text{ for }t[g+1,q1]\hfill \\ \pm 2R_{[g+1,q1]}\left(\underset{j[g+1,q1]}{}(X_j3)\right)\hfill & \text{ for }t=q.\hfill \end{array}$$
* In case $`qg+3`$, for $`\mathrm{\Xi }𝒳_{\text{nb}}[g,q]`$ of partial weight $`e_\mathrm{\Xi }=(2,2,1,\mathrm{},1,2)`$, we have
$$\vartheta _\mathrm{\Xi }=4R_{[g+1,q1]}.$$
* Assertions (a) and (b) imply the lemma.
(a, b, c) follow by the arguments of (2.29).
###### Lemma 2.36
Let $`g+1qk`$. Suppose $`[g+1,q1]L(0)L(1)`$, $`[g+1,q1]L(0)\mathrm{}`$. As ideal of coefficients, we obtain
$$\left(\vartheta _\mathrm{\Xi }|\mathrm{\Xi }𝒳[g,q],\mathrm{\Xi }_{[g+1,g+1]}=\text{}\right)=\left(2R_{[g+1,q1]}\right).$$
So, by comparison with (2.31),
$$\left(\vartheta _\mathrm{\Xi }|\mathrm{\Xi }𝒳[g,k+1],\mathrm{\Xi }_{[g+1,g+1]}=\text{}\right)=\left(2R\right).$$
Denote
$$\begin{array}{cccc}\hfill p& :=& \mathrm{min}\left(([g+1,q1]L(1))\{q\}\right)\hfill & [g+1,q]\hfill \\ \hfill s& :=& \mathrm{max}\left([g+1,q1]L(0)\right)+1\hfill & [g+2,q]\hfill \end{array}$$
so that we are in the situation $`g<g+1psq`$, more precisely, $`p=s`$ or $`ps2`$.
We shall show the following assertions.
* For $`t[s,q]`$ and a partial pattern $`\mathrm{\Xi }(q,t)𝒳_{\text{nb}}[g,q]`$ of partial weight
$$e_{\mathrm{\Xi }(q,t),j}=\{\begin{array}{cc}2\hfill & \text{ for }j[g,p]\hfill \\ 1\hfill & \text{ for }j[p+1,s1]\hfill \\ 2\hfill & \text{ for }j[s,t]\hfill \\ 1\hfill & \text{ for }j[t+1,q],\hfill \end{array}$$
we obtain
$$\vartheta _{\mathrm{\Xi }(q,t)}=\{\begin{array}{cc}\pm 2R_{[g+1,q1]}\left(\underset{j[s,t1]}{}(X_j3)\right)(X_{q1}2)\hfill & \text{ for }t[s,q1]\hfill \\ \pm 2R_{[g+1,q1]}\left(\underset{j[s,q1]}{}(X_j3)\right)\hfill & \text{ for }t=q.\hfill \end{array}$$
* In case $`s[g+2,q2]`$, for $`\mathrm{\Xi }𝒳_{\text{nb}}[g,q]`$ of partial weight
$$e_{\mathrm{\Xi },j}=\{\begin{array}{cc}2\hfill & \text{ for }j[g,p]\hfill \\ 1\hfill & \text{ for }j[p+1,s1]\hfill \\ 2\hfill & \text{ for }j=s\hfill \\ 1\hfill & \text{ for }j[s+1,q1]\hfill \\ 2\hfill & \text{ for }j=q,\hfill \end{array}$$
we obtain
$$\vartheta _\mathrm{\Xi }=4R_{[g+1,q1]}.$$
* Assertions (a) and (b) imply the lemma.
Ad (c). As in (2.29).
Ad (a). We perform an induction on $`\overline{q}[g+1,q]`$ to prove
$$\vartheta _{\mathrm{\Xi }(\overline{q},t)}=\{\begin{array}{cc}2R_{[g+1,\overline{q}1]}\hfill & \text{ for }\overline{q}[g+1,p]\hfill \\ 2R_{[g+1,\overline{q}1]}X_{\overline{q}}\hfill & \text{ for }\overline{q}[p+1,s1]\hfill \\ \pm 2R_{[g+1,\overline{q}1]}\left(\underset{j[s,\overline{q}1]}{}(X_j3)\right)\hfill & \text{ for }\overline{q}[s,t]\hfill \\ \pm 2R_{[g+1,\overline{q}1]}\left(\underset{j[s,t1]}{}(X_j3)\right)(X_{\overline{q}1}2)\hfill & \text{ for }\overline{q}[t+1,q],\hfill \end{array}$$
where $`e_{\mathrm{\Xi }(\overline{q},t)}=e_{\mathrm{\Xi }(\overline{q},\overline{q})}`$ for $`t\overline{q}+1`$. We recall that $`g<g+1pst`$.
We begin the induction by remarking that for $`\overline{q}=g+1`$ we have $`R_{[g+1,g]}=1`$, $`e_{\mathrm{\Xi }(\overline{q},t),g+1}=2`$, and thus $`\vartheta _{\mathrm{\Xi }(\overline{q},t)}=2R_{[g+1,\overline{q}1]}`$, as required.
Case I, $`ps2`$.
Assume the assertion known for $`\overline{q}1[g+1,p1]`$. Note that $`\overline{q}1L(0)`$, so
$$\begin{array}{ccc}\hfill \vartheta _{\mathrm{\Xi }(\overline{q},t)}& \stackrel{\text{(}\text{2.23}\text{)}}{=}& \vartheta _{\mathrm{\Xi }(\overline{q}1,t)}X_{\overline{q}1}(X_{\overline{q}1}+1)\hfill \\ & =& 2R_{[g+1,\overline{q}2]}X_{\overline{q}1}(X_{\overline{q}1}+1)\hfill \\ & =& 2R_{[g+1,\overline{q}1]}.\hfill \end{array}$$
Assume the assertion known for $`\overline{q}1=p`$. Note that $`\overline{q}1L(1)`$ since $`p<s`$, and that therefore $`[g+1,q1]L(1)\mathrm{}`$, so
$$\begin{array}{ccc}\hfill \vartheta _{\mathrm{\Xi }(\overline{q},t)}& \stackrel{\text{(}\text{2.25}\text{)}}{=}& \vartheta _{\mathrm{\Xi }(\overline{q}1,t)}X_{\overline{q}1}(X_{\overline{q}1}2)\hfill \\ & =& 2R_{[g+1,\overline{q}2]}X_{\overline{q}1}(X_{\overline{q}1}2)\hfill \\ & =& 2R_{[g+1,\overline{q}1]}X_{\overline{q}}.\hfill \end{array}$$
Assume the assertion known for $`\overline{q}1[p+1,s2]`$. In case $`\overline{q}1L(0)`$, we obtain
$$\begin{array}{ccc}\hfill \vartheta _{\mathrm{\Xi }(\overline{q},t)}& \stackrel{\text{(}\text{2.23}\text{)}}{=}& \vartheta _{\mathrm{\Xi }(\overline{q}1,t)}(X_{\overline{q}1}1)(X_{\overline{q}1}+1)\hfill \\ & =& 2R_{[g+1,\overline{q}2]}X_{\overline{q}1}(X_{\overline{q}1}1)(X_{\overline{q}1}+1)\hfill \\ & =& 2R_{[g+1,\overline{q}1]}X_{\overline{q}}.\hfill \end{array}$$
In case $`\overline{q}1L(1)`$, we obtain
$$\begin{array}{ccc}\hfill \vartheta _{\mathrm{\Xi }(\overline{q},t)}& \stackrel{\text{(}\text{2.25}\text{)}}{=}& \vartheta _{\mathrm{\Xi }(\overline{q}1,t)}(X_{\overline{q}1}2)\hfill \\ & =& 2R_{[g+1,\overline{q}2]}X_{\overline{q}1}(X_{\overline{q}1}2)\hfill \\ & =& 2R_{[g+1,\overline{q}1]}X_{\overline{q}}.\hfill \end{array}$$
Assume the assertion known for $`\overline{q}1=s1`$. Note that $`\overline{q}1L(0)`$, so
$$\begin{array}{ccc}\hfill \vartheta _{\mathrm{\Xi }(\overline{q},t)}& \stackrel{\text{(}\text{2.23}\text{)}}{=}& \vartheta _{\mathrm{\Xi }(\overline{q}1,t)}(X_{\overline{q}1}+1)\hfill \\ & =& 2R_{[g+1,\overline{q}2]}X_{\overline{q}1}(X_{\overline{q}1}+1)\hfill \\ & =& 2R_{[g+1,\overline{q}1]}.\hfill \end{array}$$
Assume the assertion known for $`\overline{q}1[s,t1]`$. Note that $`\overline{q}1L(1)`$, thus
$$\begin{array}{ccc}\hfill \vartheta _{\mathrm{\Xi }(\overline{q},t)}& \stackrel{\text{(}\text{2.25}\text{)}}{=}& \vartheta _{\mathrm{\Xi }(\overline{q}1,t)}(X_{\overline{q}1}(X_{\overline{q}1}3))\hfill \\ & =& \pm 2R_{[g+1,\overline{q}2]}\left(\underset{j[s,\overline{q}2]}{}(X_j3)\right)X_{\overline{q}1}(X_{\overline{q}1}3)\hfill \\ & =& \pm 2R_{[g+1,\overline{q}1]}\left(\underset{j[s,\overline{q}1]}{}(X_j3)\right).\hfill \end{array}$$
Assume the assertion known for $`\overline{q}1=t`$. Note that $`\overline{q}1L(1)`$, and consequently
$$\begin{array}{ccc}\hfill \vartheta _{\mathrm{\Xi }(\overline{q},t)}& \stackrel{\text{(}\text{2.25}\text{)}}{=}& \vartheta _{\mathrm{\Xi }(\overline{q}1,t)}X_{\overline{q}1}(X_{\overline{q}1}2)\hfill \\ & =& \pm 2R_{[g+1,\overline{q}2]}\left(\underset{j[s,t1]}{}(X_j3)\right)X_{\overline{q}1}(X_{\overline{q}1}2)\hfill \\ & =& \pm 2R_{[g+1,\overline{q}1]}\left(\underset{j[s,t1]}{}(X_j3)\right)(X_{\overline{q}1}2).\hfill \end{array}$$
Assume the assertion known for $`\overline{q}1[t+1,q1]`$. Note that $`\overline{q}1L(1)`$, whence
$$\begin{array}{ccc}\hfill \vartheta _{\mathrm{\Xi }(\overline{q},t)}& \stackrel{\text{(}\text{2.25}\text{)}}{=}& \vartheta _{\mathrm{\Xi }(\overline{q}1,t)}(X_{\overline{q}1}2)\hfill \\ & =& \pm 2R_{[g+1,\overline{q}2]}\left(\underset{j[s,t1]}{}(X_j3)\right)(X_{\overline{q}2}2)(X_{\overline{q}1}2)\hfill \\ & =& \pm 2R_{[g+1,\overline{q}1]}\left(\underset{j[s,t1]}{}(X_j3)\right)(X_{\overline{q}1}2).\hfill \end{array}$$
Case II, $`p=s`$. I.e. $`L(0)[g+1,q1]=[g+1,p1]`$, $`L(1)[g+1,q1]=[s,q1]`$.
The case $`\overline{q}[g+1,p]`$ follows from the argument for $`\overline{q}1[g+1,p1]`$ in Case I. The case $`\overline{q}[s+1,t]`$ follows from the argument for $`\overline{q}1[s,t1]`$ in Case I. The case $`\overline{q}[t+1,q]`$ follows from the arguments for $`\overline{q}1=t`$ and for $`\overline{q}1[t+1,q1]`$ in Case I.
Ad (b). We modify the last step in the proof of (a) according to (2.25), i.e. we replace the factor $`(X_{q1}2)`$ by the factor $`2`$.
###### Theorem 2.37
Suppose $`g<k`$. Let
$$m:=\{\begin{array}{ccc}X_g+2\hfill & \text{ if }\hfill & [p(1),q(1)1]L(1)\hfill \\ (X_g+2)/\mathrm{gcd}(2,X_g,X_{g+1})\hfill & \text{ if }\hfill & [p(1),q(1)1]L(1).\hfill \end{array}$$
Note that if $`g+1\overline{L}`$, the case $`\mathrm{}=[p(1),q(1)1]L(1)`$ applies. The integer
$$X_g+2=(\lambda _g^{}g)(\lambda _{k+1}^{}(k+1))+2$$
allows an interpretation as box shift length. The quotient $`f^{}/R`$ is integral and factors over
The morphism $`f`$ has order $`m`$ as an element of $`\text{Hom}_{\text{Z}𝒮_n}(S^\lambda ,S^\mu /m)`$. Recall that $`f^{}`$ maps
$$\begin{array}{ccc}\hfill M^{\lambda ^{},}& \text{}& S^\mu \hfill \\ \hfill \{a^{}\}^{}& \text{}& \underset{eE}{}(\underset{i[g+1,k]}{}X_i^{(2e_i)})\underset{\gamma \dot{\mathrm{\Gamma }}(e)}{}a^\gamma \epsilon _\gamma .\hfill \end{array}$$
The map $`f`$ operates accordingly, division taking place before going modulo $`m`$. Alternatively, by (2.31) we may write
$$\begin{array}{ccc}\hfill S^\lambda & \text{}& S^\mu /m\hfill \\ \hfill a& \text{}& \underset{\mathrm{\Xi }𝒳[g,k+1]}{}(\vartheta _\mathrm{\Xi }/R)\underset{\gamma \dot{\mathrm{\Gamma }}(\mathrm{\Xi })}{}a^\gamma \epsilon _\gamma ,\hfill \end{array}$$
the coefficients $`\vartheta _\mathrm{\Xi }/R`$ being integral.
Suppose $`g=k`$. The formula given for $`f^{}`$ in case $`g<k`$ holds in case $`g=k`$ as well, $`E`$ consisting of a single element, $`R=1`$ and $`m=(X_g+2)/\mathrm{gcd}(2,X_g)`$. We have the factorization diagram as above, the resulting morphism $`S^\lambda \text{}S^\mu /m`$ being of order $`m`$. Cf. (2.17).
Suppose $`g<k`$. By (2.31), the ideal generated by Z-linear matrix coefficients of $`f^{}/R`$ is $`(1)\text{Z}`$. In case $`[p(1),q(1)1]L(1)`$, the integer $`mR`$ divides each of the required expressions in (2.18) by (2.32) and by the ideals of coefficients calculated in (2.34, 2.36). In case $`[p(1),q(1)1]L(1)`$, the integer $`mR`$ divides each of the required expressions in (2.18) by (2.32) and by the ideals of coefficients calculated in (2.33, 2.35).
Suppose $`g=k`$. Concerning the assertion, we refer to (2.17) (alternatively, to 4.9). This case behaves as if the first component lay in $`L(1)`$ since $`X_{g+1}=0`$.
###### Remark 2.38
In the situation of (2.37), Carter and Payne \[CP 80\] have shown that
$$\text{Hom}_{K𝒮_n}(K_\text{Z}S^\lambda ,K_\text{Z}S^\mu )0,$$
$`K`$ being an infinite field of characteristic $`p`$ such that $`p`$ divides $`X_g+2`$ in case $`p3`$, respectively, such that $`4`$ divides $`X_g+2`$ in case $`p=2`$. This part of their result is recovered by (2.37).
###### Question 2.39
I do not know a description of $`\text{Hom}_{\text{Z}𝒮_n}(S^\lambda ,S^\mu /n!)`$ as an abelian group. See Section (2.4) for some examples.
### 2.4 Vertical examples
In the examples that follow, calculated by means of linear algebra (<sup>1</sup><sup>1</sup>1I thank the computer staff at the Fakultät für Mathematik in Bielefeld (‘root’) for kind cooperation.), we omit to denote the brackets that indicate polytabloids.
###### Example 2.40
Let $`\lambda =(3,3)`$, $`\mu =(2,2,1,1)`$. Then
$$\text{Hom}_{\text{Z}𝒮_6}(S^\lambda ,S^\mu /6!)\text{}\text{Hom}_{\text{Z}𝒮_6}(S^\lambda ,S^\mu /4)\text{Z}/4,$$
and a generator of $`\text{Hom}_{\text{Z}𝒮_6}(S^\lambda /4,S^\mu /4)`$ is given by
$$\begin{array}{ccc}\hfill \begin{array}{ccc}1\hfill & 3\hfill & 5\hfill \\ 2\hfill & 4\hfill & 6\hfill \end{array}& \text{}& 2\begin{array}{cc}1\hfill & 5\hfill \\ 2\hfill & 6\hfill \\ 3\hfill & \\ 4\hfill & \end{array}\left(\begin{array}{cc}1\hfill & 5\hfill \\ 2\hfill & 4\hfill \\ 3\hfill & \\ 6\hfill & \end{array}+\begin{array}{cc}1\hfill & 3\hfill \\ 2\hfill & 5\hfill \\ 4\hfill & \\ 6\hfill & \end{array}\right)\left(\begin{array}{cc}1\hfill & 6\hfill \\ 2\hfill & 4\hfill \\ 5\hfill & \\ 3\hfill & \end{array}+\begin{array}{cc}1\hfill & 3\hfill \\ 2\hfill & 6\hfill \\ 5\hfill & \\ 4\hfill & \end{array}\right)2\begin{array}{cc}1\hfill & 3\hfill \\ 2\hfill & 4\hfill \\ 5\hfill & \\ 6\hfill & \end{array}.\hfill \end{array}$$
###### Example 2.41
Let $`\lambda =(3,2,2)`$, $`\mu =(3,1,1,1,1)`$. Then
$$\text{Hom}_{\text{Z}𝒮_7}(S^\lambda ,S^\mu /7!)\text{}\text{Hom}_{\text{Z}𝒮_7}(S^\lambda ,S^\mu /6)\text{Z}/2\times \text{Z}/3,$$
generators of $`\text{Hom}_{\text{Z}𝒮_7}(S^\lambda /6,S^\mu /6)`$ being given by
$$\begin{array}{ccc}\hfill \begin{array}{ccc}1\hfill & 4\hfill & 7\hfill \\ 2\hfill & 5\hfill & \\ 3\hfill & 6\hfill & \end{array}& \text{}& 3\left(\underset{\left[b\right]\text{ standard }\mu \text{-tableau}}{}b\right)\hfill \end{array}$$
(cf. 3.19 below) and by
$$\begin{array}{ccc}\hfill \begin{array}{ccc}1\hfill & 4\hfill & 7\hfill \\ 2\hfill & 5\hfill & \\ 3\hfill & 6\hfill & \end{array}& \text{}& 2\left(\begin{array}{ccc}1\hfill & 6\hfill & 7\hfill \\ 2\hfill & & \\ 3\hfill & & \\ 4\hfill & & \\ 5\hfill & & \end{array}\begin{array}{ccc}1\hfill & 5\hfill & 7\hfill \\ 2\hfill & & \\ 3\hfill & & \\ 4\hfill & & \\ 6\hfill & & \end{array}+\begin{array}{ccc}1\hfill & 4\hfill & 7\hfill \\ 2\hfill & & \\ 3\hfill & & \\ 5\hfill & & \\ 6\hfill & & \end{array}\right).\hfill \end{array}$$
###### Example 2.42
Let $`\lambda =(3,3,1,1)`$, $`\mu =(2,2,1,1,1,1)`$. Then
$$\text{Hom}_{\text{Z}𝒮_8}(S^\lambda ,S^\mu /8!)\text{}\text{Hom}_{\text{Z}𝒮_8}(S^\lambda ,S^\mu /6)\text{Z}/6,$$
and a generator of $`\text{Hom}_{\text{Z}𝒮_8}(S^\lambda /6,S^\mu /6)`$ is given by
$$\begin{array}{ccc}\hfill \begin{array}{ccc}1\hfill & 5\hfill & 7\hfill \\ 2\hfill & 6\hfill & 8\hfill \\ 3\hfill & & \\ 4\hfill & & \end{array}& \text{}& 2\begin{array}{cc}1\hfill & 7\hfill \\ 2\hfill & 8\hfill \\ 3\hfill & \\ 4\hfill & \\ 5\hfill & \\ 6\hfill & \end{array}\left(\begin{array}{cc}1\hfill & 7\hfill \\ 2\hfill & 6\hfill \\ 3\hfill & \\ 4\hfill & \\ 5\hfill & \\ 8\hfill & \end{array}+\begin{array}{cc}1\hfill & 5\hfill \\ 2\hfill & 7\hfill \\ 3\hfill & \\ 4\hfill & \\ 6\hfill & \\ 8\hfill & \end{array}\right)\left(\begin{array}{cc}1\hfill & 8\hfill \\ 2\hfill & 6\hfill \\ 3\hfill & \\ 4\hfill & \\ 7\hfill & \\ 5\hfill & \end{array}+\begin{array}{cc}1\hfill & 5\hfill \\ 2\hfill & 8\hfill \\ 3\hfill & \\ 4\hfill & \\ 7\hfill & \\ 6\hfill & \end{array}\right)2\begin{array}{cc}1\hfill & 5\hfill \\ 2\hfill & 6\hfill \\ 3\hfill & \\ 4\hfill & \\ 7\hfill & \\ 8\hfill & \end{array}.\hfill \end{array}$$
###### Example 2.43
Let $`\lambda =(4,4)`$, $`\mu =(3,3,1,1)`$. Then
$$\text{Hom}_{\text{Z}𝒮_8}(S^\lambda ,S^\mu /8!)\text{}\text{Hom}_{\text{Z}𝒮_8}(S^\lambda ,S^\mu /5)\text{Z}/5,$$
and a generator of $`\text{Hom}_{\text{Z}𝒮_8}(S^\lambda /5,S^\mu /5)`$ is given by
$$\begin{array}{c}\begin{array}{cccc}1\hfill & 3\hfill & 5\hfill & 7\hfill \\ 2\hfill & 4\hfill & 6\hfill & 8\hfill \end{array}\text{}\hfill \\ 2\begin{array}{ccc}1\hfill & 5\hfill & 7\hfill \\ 2\hfill & 6\hfill & 8\hfill \\ 3\hfill & & \\ 4\hfill & & \end{array}\left(\begin{array}{ccc}1\hfill & 6\hfill & 7\hfill \\ 2\hfill & 4\hfill & 8\hfill \\ 5\hfill & & \\ 3\hfill & & \end{array}+\begin{array}{ccc}1\hfill & 3\hfill & 7\hfill \\ 2\hfill & 6\hfill & 8\hfill \\ 5\hfill & & \\ 4\hfill & & \end{array}\right)\left(\begin{array}{ccc}1\hfill & 5\hfill & 7\hfill \\ 2\hfill & 4\hfill & 8\hfill \\ 3\hfill & & \\ 6\hfill & & \end{array}+\begin{array}{ccc}1\hfill & 3\hfill & 7\hfill \\ 2\hfill & 5\hfill & 8\hfill \\ 4\hfill & & \\ 6\hfill & & \end{array}\right)\hfill \\ \left(\begin{array}{ccc}1\hfill & 5\hfill & 8\hfill \\ 2\hfill & 4\hfill & 6\hfill \\ 7\hfill & & \\ 3\hfill & & \end{array}+\begin{array}{ccc}1\hfill & 6\hfill & 5\hfill \\ 2\hfill & 4\hfill & 8\hfill \\ 7\hfill & & \\ 3\hfill & & \end{array}+\begin{array}{ccc}1\hfill & 3\hfill & 8\hfill \\ 2\hfill & 5\hfill & 6\hfill \\ 7\hfill & & \\ 4\hfill & & \end{array}+\begin{array}{ccc}1\hfill & 3\hfill & 5\hfill \\ 2\hfill & 6\hfill & 8\hfill \\ 7\hfill & & \\ 4\hfill & & \end{array}\right)\hfill \\ \left(\begin{array}{ccc}1\hfill & 5\hfill & 7\hfill \\ 2\hfill & 4\hfill & 6\hfill \\ 3\hfill & & \\ 8\hfill & & \end{array}+\begin{array}{ccc}1\hfill & 6\hfill & 5\hfill \\ 2\hfill & 4\hfill & 7\hfill \\ 3\hfill & & \\ 8\hfill & & \end{array}+\begin{array}{ccc}1\hfill & 3\hfill & 7\hfill \\ 2\hfill & 5\hfill & 6\hfill \\ 4\hfill & & \\ 8\hfill & & \end{array}+\begin{array}{ccc}1\hfill & 3\hfill & 5\hfill \\ 2\hfill & 6\hfill & 7\hfill \\ 4\hfill & & \\ 8\hfill & & \end{array}\right)\hfill \\ 2\begin{array}{ccc}1\hfill & 3\hfill & 7\hfill \\ 2\hfill & 4\hfill & 8\hfill \\ 5\hfill & & \\ 6\hfill & & \end{array}\left(\begin{array}{ccc}1\hfill & 3\hfill & 8\hfill \\ 2\hfill & 4\hfill & 6\hfill \\ 7\hfill & & \\ 5\hfill & & \end{array}+\begin{array}{ccc}1\hfill & 3\hfill & 5\hfill \\ 2\hfill & 4\hfill & 8\hfill \\ 7\hfill & & \\ 6\hfill & & \end{array}\right)\left(\begin{array}{ccc}1\hfill & 3\hfill & 7\hfill \\ 2\hfill & 4\hfill & 6\hfill \\ 5\hfill & & \\ 8\hfill & & \end{array}+\begin{array}{ccc}1\hfill & 3\hfill & 5\hfill \\ 2\hfill & 4\hfill & 7\hfill \\ 6\hfill & & \\ 8\hfill & & \end{array}\right)2\begin{array}{ccc}1\hfill & 3\hfill & 5\hfill \\ 2\hfill & 4\hfill & 6\hfill \\ 7\hfill & & \\ 8\hfill & & \end{array}.\hfill \end{array}$$
Note that $`R=6`$, $`\overline{m}=1`$ and
$$\begin{array}{ccc}\hfill \vartheta _{\text{}}& =& 12\hfill \\ \hfill \vartheta _{\text{}}=\vartheta _{\text{}}& =& 6\hfill \\ \hfill \vartheta _{\text{}}=\vartheta _{\text{}}& =& 6\hfill \\ \hfill \vartheta _{\text{}}& =& 12\hfill \\ \hfill \vartheta _{\text{}}=\vartheta _{\text{}}& =& 6\hfill \\ \hfill \vartheta _{\text{}}& =& 12.\hfill \end{array}$$
###### Example 2.44
Let $`\lambda =(3,3,2)`$, $`\mu =(2,2,2,1,1)`$. Then
$$\text{Hom}_{\text{Z}𝒮_8}(S^\lambda ,S^\mu /8!)\text{}\text{Hom}_{\text{Z}𝒮_8}(S^\lambda ,S^\mu /5)\text{Z}/5,$$
$`R=2`$, and a generator of $`\text{Hom}_{\text{Z}𝒮_8}(S^\lambda /5,S^\mu /5)`$ is given by
$$\begin{array}{ccc}\hfill \begin{array}{ccc}1\hfill & 4\hfill & 7\hfill \\ 2\hfill & 5\hfill & 8\hfill \\ 3\hfill & 6\hfill & \end{array}& \text{}& \left(\begin{array}{cc}1\hfill & 7\hfill \\ 2\hfill & 8\hfill \\ 3\hfill & 6\hfill \\ 4\hfill & \\ 5\hfill & \end{array}+\begin{array}{cc}1\hfill & 7\hfill \\ 2\hfill & 5\hfill \\ 3\hfill & 8\hfill \\ 4\hfill & \\ 6\hfill & \end{array}+\begin{array}{cc}1\hfill & 4\hfill \\ 2\hfill & 7\hfill \\ 3\hfill & 8\hfill \\ 5\hfill & \\ 6\hfill & \end{array}\right)\left(\begin{array}{cc}1\hfill & 7\hfill \\ 2\hfill & 5\hfill \\ 3\hfill & 6\hfill \\ 4\hfill & \\ 8\hfill & \end{array}+\begin{array}{cc}1\hfill & 4\hfill \\ 2\hfill & 7\hfill \\ 3\hfill & 6\hfill \\ 5\hfill & \\ 8\hfill & \end{array}+\begin{array}{cc}1\hfill & 4\hfill \\ 2\hfill & 5\hfill \\ 3\hfill & 7\hfill \\ 6\hfill & \\ 8\hfill & \end{array}\right)\hfill \\ & & \left(\begin{array}{cc}1\hfill & 8\hfill \\ 2\hfill & 5\hfill \\ 3\hfill & 6\hfill \\ 7\hfill & \\ 4\hfill & \end{array}+\begin{array}{cc}1\hfill & 4\hfill \\ 2\hfill & 8\hfill \\ 3\hfill & 6\hfill \\ 7\hfill & \\ 5\hfill & \end{array}+\begin{array}{cc}1\hfill & 4\hfill \\ 2\hfill & 5\hfill \\ 3\hfill & 8\hfill \\ 7\hfill & \\ 6\hfill & \end{array}\right)3\begin{array}{cc}1\hfill & 4\hfill \\ 2\hfill & 5\hfill \\ 3\hfill & 6\hfill \\ 7\hfill & \\ 8\hfill & \end{array}.\hfill \end{array}$$
###### Example 2.45
Let $`\lambda =(3,3,1,1)`$, $`\mu =(2,2,2,2)`$. Then
$$\text{Hom}_{\text{Z}𝒮_8}(S^\lambda ,S^\mu /8!)\text{}\text{Hom}_{\text{Z}𝒮_8}(S^\lambda ,S^\mu /3)\text{Z}/3,$$
and a generator of $`\text{Hom}_{\text{Z}𝒮_8}(S^\lambda /3,S^\mu /3)`$ is given by
$$\begin{array}{ccc}\hfill \begin{array}{ccc}1\hfill & 5\hfill & 7\hfill \\ 2\hfill & 6\hfill & 8\hfill \\ 3\hfill & & \\ 4\hfill & & \end{array}& \text{}& \begin{array}{cc}1\hfill & 5\hfill \\ 2\hfill & 6\hfill \\ 3\hfill & 7\hfill \\ 4\hfill & 8\hfill \end{array}.\hfill \end{array}$$
###### Example 2.46
Let $`\lambda =(4,3)`$, $`\mu =(2,2,2,1)`$. Then
$$\text{Hom}_{\text{Z}𝒮_7}(S^\lambda ,S^\mu /7!)\text{}\text{Hom}_{\text{Z}𝒮_7}(S^\lambda ,S^\mu /4)\text{Z}/4,$$
and a generator of $`\text{Hom}_{\text{Z}𝒮_7}(S^\lambda /4,S^\mu /4)`$ is given by
$$\begin{array}{ccc}\hfill \begin{array}{cccc}1\hfill & 3\hfill & 5\hfill & 7\hfill \\ 2\hfill & 4\hfill & 6\hfill & \end{array}& \text{}& 2\begin{array}{cc}1\hfill & 5\hfill \\ 2\hfill & 6\hfill \\ 3\hfill & 7\hfill \\ 4\hfill & \end{array}+\begin{array}{cc}1\hfill & 5\hfill \\ 2\hfill & 4\hfill \\ 3\hfill & 7\hfill \\ 6\hfill & \end{array}+\begin{array}{cc}1\hfill & 3\hfill \\ 2\hfill & 5\hfill \\ 4\hfill & 7\hfill \\ 6\hfill & \end{array}+\begin{array}{cc}1\hfill & 6\hfill \\ 2\hfill & 4\hfill \\ 5\hfill & 7\hfill \\ 3\hfill & \end{array}+\begin{array}{cc}1\hfill & 3\hfill \\ 2\hfill & 6\hfill \\ 5\hfill & 7\hfill \\ 4\hfill & \end{array}+2\begin{array}{cc}1\hfill & 3\hfill \\ 2\hfill & 4\hfill \\ 5\hfill & 7\hfill \\ 6\hfill & \end{array}.\hfill \end{array}$$
###### Example 2.47
Let $`\lambda =(4,4)`$, $`\mu =(2,2,2,2)`$. Then
$$\text{Hom}_{\text{Z}𝒮_8}(S^\lambda ,S^\mu /8!)\text{}\text{Hom}_{\text{Z}𝒮_8}(S^\lambda ,S^\mu /4)\text{Z}/4,$$
and a generator of $`\text{Hom}_{\text{Z}𝒮_8}(S^\lambda /4,S^\mu /4)`$ is given by
$$\begin{array}{ccc}\hfill \begin{array}{cccc}1\hfill & 3\hfill & 5\hfill & 7\hfill \\ 2\hfill & 4\hfill & 6\hfill & 8\hfill \end{array}& \text{}& 2\begin{array}{cc}1\hfill & 5\hfill \\ 2\hfill & 6\hfill \\ 3\hfill & 7\hfill \\ 4\hfill & 8\hfill \end{array}+\begin{array}{cc}1\hfill & 5\hfill \\ 2\hfill & 4\hfill \\ 3\hfill & 7\hfill \\ 6\hfill & 8\hfill \end{array}+\begin{array}{cc}1\hfill & 3\hfill \\ 2\hfill & 5\hfill \\ 4\hfill & 7\hfill \\ 6\hfill & 8\hfill \end{array}+\begin{array}{cc}1\hfill & 6\hfill \\ 2\hfill & 4\hfill \\ 5\hfill & 7\hfill \\ 3\hfill & 8\hfill \end{array}+\begin{array}{cc}1\hfill & 3\hfill \\ 2\hfill & 6\hfill \\ 5\hfill & 7\hfill \\ 4\hfill & 8\hfill \end{array}+2\begin{array}{cc}1\hfill & 3\hfill \\ 2\hfill & 4\hfill \\ 5\hfill & 7\hfill \\ 6\hfill & 8\hfill \end{array}.\hfill \end{array}$$
###### Question 2.48
I do not know a generic morphism that specializes to the morphism in (2.46) or to the morphism in (2.47).
## 3 Semistandard morphisms
The following considerations are based on methods of James \[J 78, 13\] and follow a hint of Ringel. We transpose our morphism from (2.37), cf. (3.16, 3.27), where transposition is given by dualization, followed by alternation and isomorphic substitution. Moreover, we establish a connection from the morphisms from $`M^{\lambda ,}`$ to $`S^\mu `$ exhibited in Section 2.1 to those of the form $`\mathrm{\Theta }_\phi ^{}\nu _S^\mu `$ (3.25, 3.26), where the morphisms of type $`\mathrm{\Theta }_\phi `$ occur in \[J 78, 13.13\].
Informal remark. The original aim to give upper bounds on the Hom-groups is still out of reach (cf. 2.39). It could perhaps be approached by showing that the method that has been sufficient to construct our morphism (2.37) is also necessary in some sense. The connection we establish is to be seen as the first step of this approach, namely that the construction of our building blocks $`f_e^{\prime \prime }`$ in Section 2.1 has been necessary in some sense, except at the prime $`2`$ in certain cases (3.12), and that there is a also reason to exclude bulky patterns by means of a reduced tuple of coefficients (3.28). For a particular case in which this aim could be achieved at least in the regular case, see (4.9).
### 3.1 Correspondoids
We shall recall the Carter-Lusztig-James construction of the semistandard basis of $`\text{Hom}_{\text{Z}𝒮_n}(S^\lambda ,M^\mu )`$ \[J 78, 13.13\], $`\lambda `$ and $`\mu `$ partitions of $`n`$, in our slightly modified language convenient for the purpose of dualization (cf. 3.2). Moreover, we shall give a recipe for the transposition (3.16).
In this section we write maps on the right. Inverse images of subsets are written on the left. Given a $`\text{Z}𝒮_n`$-lattice $`X`$, we denote its dual by $`X^{}=\text{Hom}_\text{Z}(X,\text{Z})`$, and, accordingly, we denote the dual of a morphism $`X\text{}Y`$ of $`\text{Z}𝒮_n`$-lattices by $`X^{}\text{}Y^{}`$. Let $`\lambda `$ and $`\mu `$ be arbitrary partitions of $`n`$. A correspondence from $`\mu `$ to $`\lambda `$ is a bijection
$$[\mu ]\text{}[\lambda ].$$
The set of correspondences from $`\mu `$ to $`\lambda `$ is denoted by $`\mathrm{\Phi }^{\mu ,\lambda }`$. We denote by
$$\{\phi \}:=\phi ^1\pi _R^\mu :[\lambda ]\text{}\text{N}$$
what is called a $`\lambda `$-tableau of type $`\mu `$ in \[J 78, 13.1\] and what in our context might also be called the correspondoid attached to the correspondence $`\phi \mathrm{\Phi }^{\mu ,\lambda }`$. The set of correspondoids from $`\mu `$ to $`\lambda `$ is denoted by $`\{\mathrm{\Phi }^{\mu ,\lambda }\}`$. James uses the bijection from the set of correspondoids $`\{\mathrm{\Phi }^{\mu ,\lambda }\}`$ to the set of $`\mu `$-tabloids that, for a fixed $`\lambda `$-tableau $`[a]`$, is given by
$$\{\phi \}=\phi ^1\pi _R^\mu \text{}[a]^1\phi ^1\pi _R^\mu =\{\phi [a]\}$$
as identification. Let
$$\begin{array}{ccc}\hfill R_\lambda & :=& \{\rho \mathrm{\Phi }^{\lambda ,\lambda }|\rho \pi _R^\lambda =\pi _R^\lambda \}\hfill \\ \hfill C_\lambda & :=& \{\kappa \mathrm{\Phi }^{\lambda ,\lambda }|\kappa \pi _C^\lambda =\pi _C^\lambda \}\hfill \end{array}$$
be the row stabilizer resp. the column stabilizer of $`\lambda `$. Given a correspondence $`\phi \mathrm{\Phi }^{\mu ,\lambda }`$, we write
$$r_\phi :=\mathrm{\#}(R_\lambda \phi ^1R_\mu \phi )$$
Note that $`r_\phi =r_{\phi ^1}`$. Given a correspondence $`\phi \mathrm{\Phi }^{\mu ,\lambda }`$ and a $`\lambda `$-tableau $`[a]`$, conjugation with the tableau $`[a]`$ yields a bijection $`R_\lambda \phi ^1R_\mu \phi \text{}R_{[a]}R_{\phi [a]}`$.
###### Lemma 3.1
Let $`\phi \mathrm{\Phi }^{\mu ,\lambda }`$ be a correspondence. The map
$$\begin{array}{ccc}\hfill M^\lambda & \text{}& M^\mu \hfill \\ \hfill \{a\}& \text{}& r_\phi ^1\underset{\rho R_{[a]}}{}\{\phi [a]\}\rho =r_\phi ^1\underset{\rho R_\lambda }{}\{\phi \rho [a]\}\hfill \end{array}$$
is well defined and $`\text{Z}𝒮_n`$-linear. For $`\rho R_\mu `$ and $`\rho ^{}R_\lambda `$, we have $`\mathrm{\Theta }_{\rho \phi \rho ^{}}=\mathrm{\Theta }_\phi `$. In particular, $`\mathrm{\Theta }_\phi `$ and $`\mathrm{\Theta }_{\phi ^1}`$ both depend only on the correspondoid $`\{\phi \}`$.
First, we note that the element $`{\displaystyle \underset{\rho R_{[a]}}{}}\{\phi [a]\}\rho M^\mu `$ is divisible by $`r_\phi =\mathrm{\#}(R_{[a]}R_{\phi [a]})`$ since $`\{\phi [a]\}\rho =\{\phi [a]\}\rho _0\rho `$ for $`\rho _0R_{[a]}R_{\phi [a]}`$. Furthermore, since $`R_{[a]\sigma }=(R_{[a]})^\sigma `$ for $`\sigma 𝒮_n`$, the map
$$\begin{array}{ccc}\hfill F^\lambda & \text{}& M^\mu \hfill \\ \hfill [a]& \text{}& r_\phi ^1\underset{\rho R_{[a]}}{}\{\phi [a]\}\rho \hfill \end{array}$$
is $`\text{Z}𝒮_n`$-linear. Finally, since $`{\displaystyle \underset{\rho R_{[a]}}{}}\{\phi [a]\}\rho ={\displaystyle \underset{\rho R_{[a]\sigma }}{}}\{\phi [a]\sigma \}\rho `$ for $`\sigma R_{[a]}`$, this map factors over
$$\begin{array}{ccc}\hfill F^\lambda & \text{}& M^\lambda \hfill \\ \hfill [a]& \text{}& \{a\}.\hfill \end{array}$$
Suppose given $`\rho R_\mu `$ and $`\rho ^{}R_\lambda `$. We have
$$\begin{array}{ccc}\hfill \{a\}\mathrm{\Theta }_{\rho \phi \rho ^{}}& =& r_{\rho \phi \rho ^{}}^1\underset{\rho ^{\prime \prime }R_\lambda }{}\{\rho \phi \rho ^{}\rho ^{\prime \prime }[a]\}\hfill \\ & =& r_\phi ^1\underset{\rho ^{\prime \prime }R_\lambda }{}\{\phi \rho ^{\prime \prime }[a]\}\hfill \\ & =& \{a\}\mathrm{\Theta }_\phi .\hfill \end{array}$$
A $`𝒮_n`$-invariant bilinear form on $`M^\lambda `$ is given by
$$(\{a\},\{b\})=\{\begin{array}{cc}1\hfill & \text{ for }\{a\}=\{b\}\hfill \\ 0\hfill & \text{ for }\{a\}\{b\},\hfill \end{array}$$
$`\{a\},\{b\}`$ being $`\lambda `$-tabloids. The resulting dualization isomorphism is denoted by
$$\begin{array}{ccc}\hfill M^\lambda & \text{}& M^{\lambda ,}\hfill \\ \hfill \{a\}& \text{}& (\{a\},).\hfill \end{array}$$
We also write such a $`\lambda `$-cotabloid as $`\{a\}^{}:=(\{a\},)`$.
###### Lemma 3.2
Given a correspondence $`\phi \mathrm{\Phi }^{\mu ,\lambda }`$, the diagram
commutes.
The dual of $`\mathrm{\Theta }_\phi `$ sends the $`\mu `$-cotabloid $`\{b\}^{}`$ to the map
$$\begin{array}{ccc}\hfill M^\lambda & \text{}& \text{Z}\hfill \\ \hfill \{a\}& \text{}& (\{b\},\{a\}\mathrm{\Theta }_\phi )\hfill \\ & =& r_\phi ^1\underset{\rho R_{[a]}}{}(\{b\},\{\phi [a]\}\rho )\hfill \\ & =& r_\phi ^1\underset{\rho R_{[a]}}{}\underset{\sigma R_{[b]}}{}\left\{\begin{array}{cc}1\hfill & \text{for }[b]\sigma =\phi [a]\rho \hfill \\ 0\hfill & \text{else}\hfill \end{array}\right\}\hfill \\ & =& r_{\phi ^1}^1\underset{\sigma R_{[b]}}{}\underset{\rho R_{[a]}}{}\left\{\begin{array}{cc}1\hfill & \text{for }\phi ^1[b]\sigma =[a]\rho \hfill \\ 0\hfill & \text{else}\hfill \end{array}\right\}\hfill \\ & =& r_{\phi ^1}^1\underset{\sigma R_{[b]}}{}(\{\phi ^1[b]\}\sigma ,\{a\})\hfill \\ & =& (\{b\}\mathrm{\Theta }_{\phi ^1},\{a\}).\hfill \end{array}$$
We denote the inclusion that is given by the definition of the Specht lattice by $`S^\lambda \text{}M^\lambda `$.
###### Remark 3.3
Let $`\nu `$ be still another partition of $`n`$, let $`\phi \mathrm{\Phi }^{\mu ,\lambda }`$, let $`\psi \mathrm{\Phi }^{\nu ,\mu }`$ and let $`[a]`$ be a $`\lambda `$-tableau. We obtain
$$\{a\}\mathrm{\Theta }_\phi \mathrm{\Theta }_\psi =\underset{\rho R_\mu }{}\frac{r_{\psi \rho \phi }}{r_\phi r_\psi }\{a\}\mathrm{\Theta }_{\psi \rho \phi }.$$
In fact,
$$\begin{array}{ccc}\hfill \{a\}\mathrm{\Theta }_\phi \mathrm{\Theta }_\psi & =& r_\phi ^1\left(\underset{\rho ^{}R_{[a]}}{}\{\phi [a]\rho ^{}\}\right)\mathrm{\Theta }_\psi \hfill \\ & =& r_\phi ^1r_\psi ^1\underset{\rho ^{}R_{[a]}}{}\underset{\rho R_{\phi [a]\rho ^{}}}{}\{\psi \phi [a]\rho ^{}\rho \}\hfill \\ & =& r_\phi ^1r_\psi ^1\underset{\rho R_{\phi [a]}}{}\underset{\rho ^{}R_{[a]}}{}\{(\psi \phi [a]\rho [a]^1\phi ^1\phi )[a]\}\rho ^{}\hfill \\ & =& r_\phi ^1r_\psi ^1\underset{\rho R_{\phi [a]}}{}r_{\psi \phi [a]\rho [a]^1\phi ^1\phi }\{a\}\mathrm{\Theta }_{\psi \phi [a]\rho [a]^1\phi ^1\phi }\hfill \\ & =& r_\phi ^1r_\psi ^1\underset{\rho R_\mu }{}r_{\psi \rho \phi }\{a\}\mathrm{\Theta }_{\psi \rho \phi }.\hfill \end{array}$$
###### Remark 3.4
$`X`$ being a $`\text{Z}𝒮_n`$-lattice, we may identify
$$\begin{array}{ccc}\hfill X^,& \text{}& X^,\hfill \\ \hfill fa& \text{}& (xb\text{}xfab)\hfill \\ \hfill (x\text{}(x1)g)1& \text{}& g.\hfill \end{array}$$
###### Proposition 3.5 (\[J 78, 6.7\])
The map
$$\begin{array}{ccc}\hfill S^\lambda ^{}& \text{}& S^{\lambda ,,}\hfill \\ \hfill a_\lambda ^{}\sigma & \text{}& (\{(a_\lambda ^{})^{}\}^{}|_{S^\lambda }1)\sigma ,\hfill \end{array}$$
where $`\sigma 𝒮_n`$, is a well defined $`\text{Z}𝒮_n`$-linear isomorphism. The diagram
commutes.
We abbreviate $`[a]=[(a_\lambda ^{})^{}]`$ in the course of the proof and claim that the kernels of the maps
$$\begin{array}{ccc}\hfill M^\lambda & \text{}& S^{\lambda ^{},}\hfill \\ \hfill \{a\}\sigma & \text{}& a^{}\sigma \epsilon _\sigma \hfill \\ \hfill M^\lambda & \text{}& S^{\lambda ,}\hfill \\ \hfill \{a\}\sigma & \text{}& (\{a\}\sigma )\delta ^\lambda \iota ^{\lambda ,}=\{a\sigma \}^{}|_{S^\lambda }\hfill \end{array}$$
coincide. Suppose given an element $`{\displaystyle \underset{\sigma 𝒮_n}{}}x_\sigma \{a\}\sigma M^\lambda `$, $`x_\sigma \text{Z}`$. To be contained in the kernel of $`\alpha `$ means
$$\begin{array}{ccc}\hfill 0& =& \underset{\sigma 𝒮_n}{}x_\sigma a^{}\sigma \epsilon _\sigma \hfill \\ & =& \underset{\sigma 𝒮_n,\rho R_{[a]}}{}x_\sigma \{a^{}\}\rho \epsilon _\rho \sigma \epsilon _\sigma \hfill \\ & =& \underset{\sigma 𝒮_n,\rho R_{[a]}}{}x_{\rho ^1\sigma }\{a^{}\}\sigma \epsilon _\sigma .\hfill \end{array}$$
Comparing coefficients, this is equivalent to saying that for all $`\tau 𝒮_n`$ we obtain
$$\underset{\kappa C_{[a]},\rho R_{[a]}}{}x_{\rho ^1\kappa \tau }\epsilon _{\kappa \tau }=0.$$
To be contained in the kernel of $`\beta `$ amounts to
$$\begin{array}{ccc}\hfill 0& =& (\underset{\sigma 𝒮_n}{}x_\sigma \{a\}\sigma ,a\tau )\hfill \\ & =& (\underset{\sigma 𝒮_n}{}x_\sigma \{a\}\sigma \tau ^1\underset{\kappa C_{[a]}}{}\kappa \epsilon _\kappa ,\{a\})\hfill \\ & =& (\underset{\kappa C_{[a]},\rho 𝒮_n}{}x_{\rho \kappa ^1\tau }\{a\}\rho \epsilon _\kappa ,\{a\})\hfill \end{array}$$
for all $`\tau 𝒮_n`$, i.e. that for all $`\tau 𝒮_n`$
$$\underset{\kappa C_{[a]},\rho R_{[a]}}{}x_{\rho \kappa ^1\tau }\epsilon _\kappa =0.$$
Since $`\alpha `$ is surjective, we may conclude that the map $`S^{\lambda ^{},}\text{}S^{\lambda ,}`$ is well defined and injective. To prove surjectivity, it remains to be shown that $`\beta =\delta ^\lambda \iota ^{\lambda ,}`$ is surjective, i.e. that $`\iota ^\lambda `$ is a pure monomorphism. But the tuple of standard tableaux gives a basis of the Specht module over $`\text{F}_p`$ for any prime $`p`$ \[J 78, 8.4\].
###### Lemma 3.6
Denoting evaluation by $`X\text{}X^{}`$, $`X`$ a $`\text{Z}𝒮_n`$-lattice, we may dualize to obtain
$$\begin{array}{ccccccccc}\hfill (M^\lambda & \text{}& M^{\lambda ,}\hfill & \text{}& M^{\lambda ,})\hfill & =& (M^\lambda \hfill & \text{}& M^{\lambda ,})\hfill \\ \hfill (S^\lambda & \text{}& S^{\lambda ,}\hfill & \text{}& S^{\lambda ^{},,})\hfill & =& (S^\lambda \hfill & \text{}& S^{\lambda ^{},,}),\hfill \end{array}$$
provided $`[a_\lambda ^{}]=[(a_\lambda )^{}]`$. Ignoring this condition, the second equality holds up to sign.
Suppose given $`\lambda `$-tableaux $`[a]`$ and $`[b]`$. Let $`\sigma =[a][b]^1`$. We have
$$(\{a\},b)=\{\begin{array}{cc}0\hfill & \text{for }\sigma R_\lambda C_\lambda \hfill \\ \epsilon _\kappa \hfill & \text{for }\sigma =\rho \kappa ,\rho R_\lambda ,\kappa C_\lambda .\hfill \end{array}$$
Therefore,
$$\text{(\#)}(\{a\},b)=\epsilon _\sigma (\{b^{}\},a^{}).$$
We abbreviate $`[a]=[(a_\lambda ^{})^{}]=[a_\lambda ]`$. The map $`S^\lambda ^{}\text{}S^{\lambda ,,}`$ sends $`a^{}\tau `$ to $`(\{[a]\tau \},)\epsilon _\tau `$, where $`\tau 𝒮_n`$, so that the composition sends
$$\begin{array}{ccccc}\hfill S^\lambda & \text{}& S^{\lambda ,}\hfill & \text{}& S^{\lambda ^{},,}\hfill \\ \hfill a& \text{}& a\text{eva}\hfill & \text{}& \left(a^{}\tau 1\text{}\epsilon _\tau (\{[a]\tau \},a)\stackrel{\text{(\#)}}{=}(\{a^{}\},a^{}\tau )\right).\hfill \end{array}$$
A correspondoid $`\{\phi \}\{\mathrm{\Phi }^{\mu ,\lambda }\}`$ is called semistandard in case the following conditions hold.
* For $`i\times j,i\times j^{}[\lambda ]`$ with $`j<j^{}`$, we have $`(i\times j)\{\phi \}(i\times j^{})\{\phi \}`$.
* For $`i\times j,i^{}\times j[\lambda ]`$ with $`i<i^{}`$, we have $`(i\times j)\{\phi \}<(i^{}\times j)\{\phi \}`$.
Let $`\{\mathrm{\Phi }^{\mu ,\lambda }\}_{\text{sst}}\{\mathrm{\Phi }^{\mu ,\lambda }\}`$ denote the subset of semistandard correspondoids.
We attach to each correspondence $`\phi \mathrm{\Phi }^{\mu ,\lambda }`$ its column distribution
$$\begin{array}{ccccc}\hfill \text{N}& \times & \text{N}\hfill & \text{}& \text{N}\hfill \\ \hfill j& \times & k\hfill & \text{}& |\phi |_{k,j}:=\mathrm{\#}\left(\{\phi \}^1(j)(\text{N}\times \{k\})\right).\hfill \end{array}$$
The map from the correspondences to the column distribution factors over the correspondoids, $`(\phi \text{}|\phi |)=(\phi \text{}\{\phi \}\text{}|\phi |)`$. For $`\phi \mathrm{\Phi }^{\mu ,\lambda }`$, $`\rho R_\mu `$ and $`\kappa C_\lambda `$, we have $`|\rho \phi \kappa |=|\phi |`$ because $`\{\rho \phi \kappa \}=\{\phi \kappa \}`$ and because of the bijection, $`j\times k\text{N}\times \text{N}`$,
$$\{\phi \}^1(j)(\text{N}\times \{k\})\text{}\{\phi \kappa \}^1(j)(\text{N}\times \{k\})$$
which is given by restriction of $`\kappa `$. Note that the column distribution may be regarded as column equivalence class in the sense of \[J 78, 13.8\].
A total order on the set of column distributions can be defined as follows. We order $`\text{N}\times \text{N}`$ lexicographically via $`j\times k<j^{}\times k^{}`$ if $`j<j^{}`$ or ($`j=j^{}`$ and $`k<k^{}`$). Suppose given correspondences $`\phi ,\stackrel{~}{\phi }\mathrm{\Phi }^{\mu ,\lambda }`$, $`|\phi ||\stackrel{~}{\phi }|`$. Let $`j\times k`$ be minimal with $`|\phi |_{k,j}|\stackrel{~}{\phi }|_{k,j}`$. We say that $`|\phi |<|\stackrel{~}{\phi }|`$ in case $`|\phi |_{k,j}<|\stackrel{~}{\phi }|_{k,j}`$.
The following arguments, needed to obtain (3.12), are taken from \[J 78\], and we reproduce them here in a slightly adapted manner. I thank M. Härterich for an explanation.
###### Lemma 3.7
For $`\{\phi \}\{\mathrm{\Phi }^{\mu ,\lambda }\}_{\text{sst}}`$ and $`\rho R_\lambda `$, we have $`|\phi ||\phi \rho |`$.
Suppose $`|\phi ||\phi \rho |`$ and let $`j\times k`$ be minimal with $`|\phi |_{k,j}|\phi \rho |_{k,j}`$. By induction on $`j^{}\times k^{}`$, we see that
$$()\{\phi \rho \}^1(j^{})(\text{N}\times \{k^{}\})=\{\phi \}^1(j^{})(\text{N}\times \{k^{}\})$$
for $`1\times 1j^{}\times k^{}<j\times k`$. In fact, assume that the two inverse images of the value $`j^{}`$ in column $`k^{}`$ differ in spite of having the same cardinality. Then there exists $`i\times k^{}[\lambda ]`$, $`i\times k^{\prime \prime }:=(i\times k^{})\rho ^1`$, such that $`j^{}=(i\times k^{\prime \prime })\{\phi \}(i\times k^{})\{\phi \}=:j^{\prime \prime }`$. Assume that $`j^{}>j^{\prime \prime }`$. But since $`i\times k^{}\{\phi \}^1(j^{\prime \prime })\backslash \{\phi \rho \}^1(j^{\prime \prime })`$, this yields a contradiction to ($``$) by induction hypothesis, so that we obtain $`j^{}<j^{\prime \prime }`$.
Let $`i\times k^{(s)}:=(i\times k^{})\rho ^s`$ for $`s0`$. We claim that $`(i\times k^{(s)})\{\phi \}=j^{}`$ for all $`s1`$, thus deriving a contradiction that establishes $`()`$. We perform an induction on $`s`$. But $`(i\times k^{(s)})\{\phi \}=j^{}`$ implies by semistandardness of $`\{\phi \}`$ that $`k^{(s)}<k^{}`$. Thus ($``$) applies by induction hypothesis and shows that $`(i\times k^{(s+1)})\{\phi \}=(i\times k^{(s)})\{\phi \rho \}=j`$.
Assume that $`|\phi |_{k,j}<|\phi \rho |_{k,j}`$. Then there exists $`i\times k[\lambda ]`$, $`i\times k^{}:=(i\times k)\rho ^1`$, such that $`j=(i\times k^{})\{\phi \}(i\times k)\{\phi \}=:j^{}`$. Assume that $`j>j^{}`$. But since $`i\times k\{\phi \}^1(j^{})\backslash \{\phi \rho \}^1(j^{})`$, this yields a contradiction to ($``$) and shows $`j<j^{}`$.
Let $`i\times k^{(s)}:=(i\times k)\rho ^s`$ for $`s0`$. We claim that $`(i\times k^{(s)})\{\phi \}=j`$ for all $`s1`$, thus deriving a contradiction that shows $`|\phi ||\phi \rho |`$. We perform an induction on $`s`$. But $`(i\times k^{(s)})\{\phi \}=j`$ implies by semistandardness of $`\{\phi \}`$ that $`k^{(s)}<k`$. Thus ($``$) shows that $`(i\times k^{(s+1)})\{\phi \}=(i\times k^{(s)})\{\phi \rho \}=j`$.
###### Lemma 3.8
Restricted to $`\{\mathrm{\Phi }^{\mu ,\lambda }\}_{\text{sst}}`$, the map $`\{\phi \}\text{}|\phi |`$ becomes injective.
To determine a semistandard correspondoid $`\{\phi \}`$, it suffices to know the set of its values on column $`k`$ for each $`k1`$, which can be written as $`\{j\text{N}||\phi |_{k,j}1\}`$.
###### Lemma 3.9
Assume $`\lambda `$ to be a $`2`$-regular partition of $`n`$, i.e. assume that $`i<i^{}`$ implies $`\lambda _i>\lambda _i^{}`$ for $`i,i^{}[1,\lambda _1^{}]`$. Suppose given $`\kappa ,\kappa ^{}C_\lambda `$. Let
$$\begin{array}{ccc}\hfill [\lambda ]& \text{}& [\lambda ]\hfill \\ \hfill i\times k& \text{}& i\times (\lambda _ik).\hfill \end{array}$$
If $`\kappa \iota \kappa ^{}R_\lambda `$ then $`\kappa =1`$, hence also $`\kappa ^{}=1`$.
Assume that $`\kappa 1`$, i.e. assume given some position $`i\times k[\lambda ]`$ such that $`i\times k(i\times k)\kappa `$. First, we choose $`i`$ minimal such that there exists a $`k`$ such that $`i\times k(i\times k)\kappa `$. Having chosen our $`i`$, we choose $`k`$ maximal such that $`i\times k(i\times k)\kappa `$. We write $`i^{}\times k:=(i\times k)\kappa `$ and remark that $`i<i^{}`$. We obtain
$$\begin{array}{ccccccccccccccc}\hfill i& \times & k\hfill & \text{}& \hfill i^{}& \times & k\hfill & \text{}& \hfill i^{}& \times & (\lambda _i^{}k)\hfill & \text{}& \hfill i^{\prime \prime \prime }& \times & (\lambda _i^{}k)\hfill \\ \hfill i& \times & (k+\lambda _i\lambda _i^{})\hfill & \text{}& \hfill i& \times & (k+\lambda _i\lambda _i^{})\hfill & \text{}& \hfill i& \times & (\lambda _i^{}k)\hfill & \text{}& \hfill i^{\prime \prime }& \times & (\lambda _i^{}k),\hfill \end{array}$$
where $`i^{\prime \prime }i^{\prime \prime \prime }`$, contradicting $`\kappa \iota \kappa ^{}R_\lambda `$.
###### Lemma 3.10
Suppose given $`\sigma R_\lambda C_\lambda `$. There exists a transposition $`\kappa _\sigma C_\lambda `$ such that $`\sigma \kappa _\sigma \sigma ^1R_\lambda `$.
Assume that the conclusion does not hold, i.e. that $`\sigma ^1`$ sends different entries of each given column to different rows. We claim that this implies that $`\sigma R_\lambda C_\lambda `$. Given a permutation $`\tau `$ of $`[\lambda ]`$, we let $`k_\tau `$ be the maximal column position $`k0`$ such that $`\tau `$ becomes the identity when restricted to $`[\lambda ](\text{N}\times [1,k])`$. We perform a downwards induction on $`k_\sigma `$, starting with $`k_\sigma =\lambda _1`$.
Different entries of column $`k_\sigma +1`$ are mapped to different rows. Hence there is a column permutation $`\kappa C_\lambda `$ acting trivially except perhaps in column $`k_\sigma +1`$ such that $`(i\times (k_\sigma +1))\kappa \sigma ^1\pi _R^\lambda =i`$ for $`i[1,\lambda _{k_\sigma +1}^{}]`$. Thus there is a row permutation $`\rho R_\lambda `$ such that $`k_{\kappa \sigma ^1\rho }>k_{\sigma ^1}`$, i.e. $`k_{\rho ^1\sigma \kappa ^1}>k_\sigma `$. By induction, $`\rho ^1\sigma \kappa ^1R_\lambda C_\lambda `$ ensues.
###### Lemma 3.11
Suppose given a semistandard correspondoid $`\{\phi \}\{\mathrm{\Phi }^{\mu ,\lambda }\}_{\text{sst}}`$ and a $`\lambda `$-tableau $`[a]`$. Writing
$$a\mathrm{\Theta }_\phi =\underset{\{\psi \}\{\mathrm{\Phi }^{\mu ,\lambda }\}}{}z_\psi \{\psi [a]\}M^\mu ,$$
we obtain $`z_{\{\phi \}}=1`$ and $`z_{\{\psi \}}=0`$ for $`|\psi |>|\phi |`$.
Note that
$$\begin{array}{ccc}\hfill a\mathrm{\Theta }_\phi & =& \underset{\kappa C_{[a]}}{}\underset{\rho (R_\lambda \phi ^1R_\mu \phi )\backslash R_\lambda }{}\{\phi \rho [a]\}\kappa \epsilon _\kappa \hfill \\ & =& \underset{\kappa C_\lambda }{}\underset{\rho (R_\lambda \phi ^1R_\mu \phi )\backslash R_\lambda }{}\{\phi \rho \kappa [a]\}\epsilon _\kappa .\hfill \end{array}$$
Since $`\{\phi \}`$ is semistandard, we have $`|\phi |\stackrel{\text{(}\text{3.7}\text{)}}{}|\phi \rho |=|\phi \rho \kappa |`$ for $`\rho R_\lambda `$, $`\kappa C_\lambda `$.
The equality $`\{\phi \}=\{\phi \rho \kappa \}`$ is equivalent to the existence of a $`\rho ^{}R_\mu `$ such that $`\kappa =\rho (\phi ^1\rho ^{}\phi )`$. This condition supposed to hold and given $`i\times k[\lambda ]`$, we obtain
$$i\times k\text{}i\times k^{}\text{}i^{}\times k,$$
where $`(i\times k^{})\{\phi \}=(i^{}\times k)\{\phi \}`$. Note that the resulting map $`i\text{}i^{}`$ is a bijection for each $`k`$. We perform an induction on $`k`$ to show that $`ii^{}`$. Suppose $`ii^{}`$ and hence $`i=i^{}`$ for $`k_0<k`$ (supposition possibly empty). Therefore, $`(i\times k_0)\{\phi \}=(i^{}\times k_0)\{\phi \}=(i\times k_0^{})\{\phi \}`$. Since $`\{\phi \}`$ is semistandard, this inhibits $`(i\times k)\{\phi \}>(i\times k^{})\{\phi \}`$, for this would mean that $`\rho `$ could not act on the set
$$\left\{i\times k_1\right|k_1[1,\lambda _i],(i\times k_1)\{\phi \}=(i\times k^{})\{\phi \}\}\{i\}\times [1,k1],$$
the inclusion given by semistandardness of $`\{\phi \}`$. Thus $`(i\times k)\{\phi \}(i\times k^{})\{\phi \}=(i^{}\times k)\{\phi \}`$, hence, by semistandardness of $`\{\phi \}`$, $`ii^{}`$, forcing $`i=i^{}`$ for all $`i\times k[\lambda ]`$. We conclude that $`\kappa =1`$ and that thus $`\phi ^1\rho ^{}\phi `$ inverts $`\rho `$ so that $`\rho `$ represents the trivial coset.
Let $`m0`$. $`\lambda `$ is called $`2`$-singular if there exists an $`i1`$ such that $`\lambda _i=\lambda _{i+1}>0`$. The case of $`m1`$, $`m(2)`$ and $`\lambda `$ being $`2`$-singular shall be called the singular case. The complementary case shall be called the regular case, comprising in particular the case $`m=0`$.
###### Theorem 3.12 (Carter, Lusztig \[CL 74, 3.5\], James \[J 78, 13.13\])
In the regular case, the tuple
$$\left(\mathrm{\Theta }_\phi |_{S^\lambda }_\text{Z}\text{Z}/m|\{\phi \}\{\mathrm{\Phi }^{\mu ,\lambda }\}_{\text{sst}}\right)$$
furnishes a $`\text{Z}/m`$-linear basis of $`\text{Hom}_{\text{Z}𝒮_n}(S^\lambda /m,M^\mu /m)`$.
In the singular case, the tuple
$$\left(2\mathrm{\Theta }_\phi |_{S^\lambda }_\text{Z}\text{Z}/m|\{\phi \}\{\mathrm{\Phi }^{\mu ,\lambda }\}_{\text{sst}}\right)$$
furnishes a $`\text{Z}/(m/2)`$-linear basis of $`2\text{Hom}_{\text{Z}𝒮_n}(S^\lambda /m,M^\mu /m)`$.
We claim that the respective tuple is linearly independent. Suppose given a $`\lambda `$-tableau $`a`$. We write down a matrix whose positions are indexed by $`\{\mathrm{\Phi }^{\mu ,\lambda }\}_{\text{sst}}\times \{\mathrm{\Phi }^{\mu ,\lambda }\}`$, and whose entry at position $`\{\phi \}\times \{\psi \}`$ is given by the multiplicity of the $`\mu `$-tabloid $`\{\psi [a]\}`$ in $`a\mathrm{\Theta }_\phi M^\mu `$. It suffices to show that the $`\{\mathrm{\Phi }^{\mu ,\lambda }\}_{\text{sst}}\times \{\mathrm{\Phi }^{\mu ,\lambda }\}_{\text{sst}}`$-part of our matrix is unipotent. But this follows from (3.11) in view of (3.8).
We claim that the respective tuple generates. Suppose given a morphism $`S^\lambda /m\text{}M^\mu /m`$ such that $`\mathrm{\Theta }0`$ in the regular case resp. such that $`2\mathrm{\Theta }0`$ in the singular case. We fix a $`\lambda `$-tableau $`[a]`$ and write
$$a\mathrm{\Theta }=\underset{\{\phi \}\{\mathrm{\Phi }^{\mu ,\lambda }\}}{}x_{\{\phi \}}\{\phi [a]\},$$
where all elements of $`S^\lambda `$ and of $`M^\mu `$ which occur in this expression and in the remainder of the proof are to be read as representing their residue classes modulo $`m`$. Suppose given $`\{\psi \}\{\mathrm{\Phi }^{\mu ,\lambda }\}`$ that allows for $`\{\psi \tau \}=\{\psi \}`$ for some $`i\times k,i^{}\times k[\lambda ]`$, $`ii^{}`$, $`\tau :=(i\times k,i\times k^{})`$. We wish to see that $`x_{\{\psi \}}=0`$ in the regular case and that $`2x_{\{\psi \}}=0`$ in the singular case. In both cases, we may compare the coefficients of
$$\begin{array}{ccc}\hfill \underset{\{\phi \}\{\mathrm{\Phi }^{\mu ,\lambda }\}}{}x_{\{\phi \tau \}}\{\phi [a]\}& =& \underset{\{\phi \}\{\mathrm{\Phi }^{\mu ,\lambda }\}}{}x_{\{\phi \}}\{\phi [a]\}(a_{k,i},a_{k,i^{}})\hfill \\ & =& a(a_{k,i},a_{k,i^{}})\mathrm{\Theta }\hfill \\ & =& \underset{\{\phi \}\{\mathrm{\Phi }^{\mu ,\lambda }\}}{}x_{\{\phi \}}\{\phi [a]\}\hfill \end{array}$$
to obtain
$$x_{\{\psi \}}=x_{\{\psi \tau \}}=x_{\{\psi \}}.$$
Now suppose the partition $`\lambda `$ to be $`2`$-regular. We need to see that $`x_{\{\psi \}}=0`$. Let $`\iota `$ be the permutation of $`[\lambda ]`$ defined by $`i\times k\text{}i\times (\lambda _ik)`$, $`i\times k[\lambda ]`$. We abbreviate $`\iota ^{[a]}:=[a]^1\iota [a]𝒮_n`$ and calculate
$$\begin{array}{ccc}\hfill \underset{\kappa ^{}C_{[a]}}{}a\iota ^{[a]}\kappa ^{}\epsilon _\kappa ^{}& =& \underset{\kappa ,\kappa ^{}C_{[a]}}{}\{a\}\kappa \iota ^{[a]}\kappa ^{}\epsilon _{\kappa \kappa ^{}}\hfill \\ & =& \underset{\kappa ,\kappa ^{}C_\lambda }{}\{\kappa \iota \kappa ^{}[a]\}\epsilon _{\kappa \kappa ^{}},\hfill \end{array}$$
where $`\kappa C_\lambda `$ has a sign $`\epsilon _\kappa `$ as a permutation of $`[\lambda ]`$, i.e. as an element of $`𝒮_{[\lambda ]}`$. Writing the resulting expression in the form $`{\displaystyle \underset{\sigma R_\lambda \backslash 𝒮_{[\lambda ]}}{}}y_\sigma \{\sigma [a]\}`$, $`y_\sigma \text{Z}`$, where we choose our coset representatives $`\sigma `$ to lie in $`C_\lambda `$ whenever (uniquely) possible, we obtain
$$y_\sigma =\underset{\kappa ,\kappa ^{}C_\lambda ,\kappa \iota \kappa ^{}R_\lambda \sigma }{}\epsilon _{\kappa \kappa ^{}}.$$
In case $`\sigma C_\lambda `$, the condition on the indexing elements yields $`\kappa =1`$, $`\kappa ^{}=\sigma `$ and thus $`y_\sigma =\epsilon _\sigma `$ (3.9). In case $`\sigma R_\lambda C_\lambda `$, we obtain, using (3.10), and choosing a fixed system of representatives of $`C_\lambda /\kappa _\sigma `$,
$$\begin{array}{ccc}\hfill y_\sigma & =& \underset{\kappa ,\kappa ^{}C_\lambda ,\kappa \iota \kappa ^{}R_\lambda \sigma }{}\epsilon _{\kappa \kappa ^{}}\hfill \\ & =& \underset{\kappa C_\lambda ,\kappa ^{}C_\lambda /\kappa _\sigma ,\kappa \iota \kappa ^{}R_\lambda \sigma }{}\epsilon _{\kappa \kappa ^{}}+\underset{\kappa C_\lambda ,\kappa ^{}C_\lambda /\kappa _\sigma ,\kappa \iota \kappa ^{}\kappa _\sigma R_\lambda \sigma }{}\epsilon _{\kappa \kappa ^{}\kappa _\sigma }\hfill \\ & =& 0.\hfill \end{array}$$
Hence
$$\underset{\kappa ^{}C_{[a]}}{}a\iota ^{[a]}\kappa ^{}\epsilon _\kappa ^{}=a,$$
and therefore
$$\begin{array}{ccc}\hfill a\mathrm{\Theta }& \stackrel{\text{1.}}{=}& \underset{\{\phi \}\{\mathrm{\Phi }^{\mu ,\lambda }\}}{}x_{\{\phi \}}\{\phi [a]\}\hfill \\ & \stackrel{\text{2.}}{=}& \underset{\kappa C_{[a]}}{}a\mathrm{\Theta }\iota ^{[a]}\kappa \epsilon _\kappa \hfill \\ & =& \underset{\{\phi \}\{\mathrm{\Phi }^{\mu ,\lambda }\}}{}x_{\{\phi \}}\left(\underset{\kappa C_\lambda }{}\{\phi \kappa [a]\}\epsilon _\kappa \right).\hfill \end{array}$$
For $`\{\phi \}\{\mathrm{\Phi }^{\mu ,\lambda }\}`$ and $`\kappa C_\lambda `$, the equality $`\{\psi \}=\{\phi \kappa \}`$ is equivalent to $`\phi \kappa \psi ^1R_\mu `$. Hence, as multiplicity of $`\{\psi [a]\}`$ in $`{\displaystyle \underset{\kappa C_\lambda }{}}\{\phi \kappa [a]\}`$ we get, choosing a fixed system of representatives for $`C_\lambda /\tau `$,
$$\begin{array}{ccc}\hfill \underset{\kappa C_\lambda ,\phi \kappa \psi ^1R_\mu }{}\epsilon _\kappa & =& \underset{\kappa C_\lambda /\tau ,\phi \kappa \psi ^1R_\mu }{}\epsilon _\kappa +\underset{\kappa C_\lambda /\tau ,\phi \kappa \tau \psi ^1R_\mu }{}\epsilon _{\kappa \tau }\hfill \\ & =& 0.\hfill \end{array}$$
Therefore, a comparison of 1. and 2. yields $`x_{\{\psi \}}=0`$.
Given $`\kappa C_\lambda `$, we obtain
$$\begin{array}{ccc}\hfill \underset{\{\phi \}\{\mathrm{\Phi }^{\mu ,\lambda }\}}{}x_{\{\phi \kappa ^1\}}\{\phi [a]\}& =& a\mathrm{\Theta }[a]^1\kappa [a]\hfill \\ & =& a\mathrm{\Theta }\epsilon _\kappa \hfill \\ & =& \underset{\{\phi \}\{\mathrm{\Phi }^{\mu ,\lambda }\}}{}\epsilon _\kappa x_{\{\phi \}}\{\phi [a]\},\hfill \end{array}$$
whence, up to sign, $`x_{\{\phi \}}`$ depends only on $`|\phi |`$. Let $`|\phi |`$ be the maximal column distribution such that $`x_{\{\phi \}}0`$ in the regular case, resp. such that $`2x_{\{\phi \}}0`$ in the singular case. By our assertion on $`x_{\{\psi \}}`$, we may assume that $`(i\times k)\{\phi \}<(i^{}\times k)\{\phi \}`$ for $`i\times k,i^{}\times k\lambda `$, $`i<i^{}`$. In order to show that $`\{\phi \}`$ is semistandard, we assume the contrary, i.e. the existence of positions $`i\times k,i\times (k+1)[\lambda ]`$ such that $`(i\times k)\{\phi \}>(i\times (k+1))\{\phi \}`$. Let $`\xi :=a_{k,[i,\lambda _i^{}]}`$, let $`\eta :=a_{k+1,[1,i]}`$. For $`\sigma `$ in $`𝒮_{\xi \eta }`$ but not in $`𝒮_\xi \times 𝒮_\eta `$, we have $`|\phi |<|\phi [a]\sigma [a]^1|`$ since the minimal value of $`\{\phi \}`$ changing columns under $`\sigma `$ cannot move to the right. Summing up over cosets, a Garnir relation gives
$$\begin{array}{ccc}\hfill 0& =& \underset{\sigma 𝒮_\xi \times 𝒮_\eta \backslash 𝒮_{\xi \eta }}{}a\mathrm{\Theta }\sigma \hfill \\ & =& \underset{\{\chi \}\{\mathrm{\Phi }^{\mu ,\lambda }\}}{}\left(\underset{\sigma 𝒮_\xi \times 𝒮_\eta \backslash 𝒮_{\xi \eta }}{}x_{\{\chi ([a]\sigma ^1[a]^1)\}}\right)\{\chi [a]\},\hfill \end{array}$$
whence $`{\displaystyle \underset{\sigma 𝒮_\xi \times 𝒮_\eta \backslash 𝒮_{\xi \eta }}{}}x_{\{\phi ([a]\sigma ^1[a]^1)\}}=0`$. For each $`\sigma `$ in $`𝒮_{\xi \eta }`$ but not in $`𝒮_\xi \times 𝒮_\eta `$, maximality of $`|\phi |`$ forces $`x_{\{\phi ([a]\sigma [a]^1)\}}=0`$ in the regular case, resp. $`2x_{\{\phi ([a]\sigma [a]^1)\}}=0`$ in the singular case, contradicting $`x_{\{\phi \}}0`$ in the regular case, resp. $`2x_{\{\phi \}}0`$ in the singular case. Hence $`\{\phi \}`$ is semistandard, and its column distribution $`|\phi |`$ shall be called the leading term of $`\mathrm{\Theta }`$. Given a correspondoid $`\{\chi \}\{\mathrm{\Phi }^{\mu ,\lambda }\}`$, we write $`\overline{\mathrm{\Theta }}_\chi :=\mathrm{\Theta }_\chi |_{S^\lambda }_\text{Z}\text{Z}/m`$.
We consider the regular case. By (3.11), the difference $`\mathrm{\Theta }x_{\{\phi \}}\overline{\mathrm{\Theta }}_\phi `$ has either a strictly smaller leading term than $`\mathrm{\Theta }`$ or it vanishes. Both alternatives allow to assume, by induction on the leading term of $`\mathrm{\Theta }`$, or, respectively, directly, that $`\mathrm{\Theta }x_{\{\phi \}}\overline{\mathrm{\Theta }}_\phi `$ is in the linear span of the tuple given above in the regular case.
We consider the singular case. By (3.11), the difference $`\mathrm{\Theta }x_{\{\phi \}}\overline{\mathrm{\Theta }}_\phi `$ has either a strictly smaller leading term than $`\mathrm{\Theta }`$ or it vanishes under multiplication with $`2`$. Both alternatives allow to assume, by induction on the leading term of $`\mathrm{\Theta }`$ or, respectively, directly, that $`2\mathrm{\Theta }x_{\{\phi \}}2\overline{\mathrm{\Theta }}_\phi `$ is in the linear span of the tuple given above in the singular case.
###### Corollary 3.13
In the regular case, the tuple
$$\left(\mathrm{\Theta }_{\phi ^1}^{}\nu _S^\lambda ^{}_\text{Z}\text{Z}/m|\{\phi \}\{\mathrm{\Phi }^{\mu ,\lambda }\}_{\text{sst}}\right)$$
furnishes a $`\text{Z}/m`$-linear basis of $`\text{Hom}_{\text{Z}𝒮_n}(M^{\mu ,}/m,S^\lambda ^{}/m)`$.
In the singular case, the tuple
$$\left(2\mathrm{\Theta }_{\phi ^1}^{}\nu _S^\lambda ^{}_\text{Z}\text{Z}/m|\{\phi \}\{\mathrm{\Phi }^{\mu ,\lambda }\}_{\text{sst}}\right)$$
furnishes a $`\text{Z}/(m/2)`$-linear basis of $`2\text{Hom}_{\text{Z}𝒮_n}(M^{\mu ,}/m,S^\lambda ^{}/m)`$.
We write $`()^{}:=\text{Hom}_\text{Z}(,\text{Z}/m)`$, neglect to denote $`()_\text{Z}\text{Z}/m`$ for maps and obtain $`\text{Z}/m`$-linear isomorphisms
$$\begin{array}{ccc}\hfill \text{Hom}_{\text{Z}𝒮_n}(S^\lambda /m,M^\mu /m)& \hfill \text{}& \text{Hom}_{\text{Z}𝒮_n}(M^\mu /m,S^{\lambda ,}/m)\hfill \\ \hfill \iota ^\lambda \mathrm{\Theta }_\phi & \hfill \text{}& ^\mu \mathrm{\Theta }_\phi ^{}\iota ^{\lambda ,}\stackrel{\text{(}\text{3.2}\text{)}}{=}\mathrm{\Theta }_{\phi ^1}^\lambda \iota ^{\lambda ,}\hfill \\ & \hfill \mathrm{}\text{}& \text{Hom}_{\text{Z}𝒮_n}(M^{\mu ,}/m,S^\lambda ^{}/m)\hfill \\ & \hfill \mathrm{}\text{}& \mathrm{\Theta }_{\phi ^1}^{}^{\lambda ,}\iota ^{\lambda ,,}(\delta ^\lambda )^1\stackrel{\text{(}\text{3.5}\text{)}}{=}\mathrm{\Theta }_{\phi ^1}^{}\nu _S^\lambda ^{}\hfill \end{array}$$
which send the bases of (3.12) to the tuples just described.
###### Corollary 3.14
In the regular case, the induced maps
$$\begin{array}{ccc}\hfill \text{Hom}_{\text{Z}𝒮_n}(M^\lambda /m,M^\mu /m)& \text{}& \text{Hom}_{\text{Z}𝒮_n}(S^\lambda /m,M^\mu /m)\hfill \\ \hfill \text{Hom}_{\text{Z}𝒮_n}(M^{\mu ,}/m,M^{\lambda ,}/m)& \text{}& \text{Hom}_{\text{Z}𝒮_n}(M^{\mu ,}/m,S^\lambda ^{}/m)\hfill \end{array}$$
are surjective. In the singular case, their cokernels are annihilated by multiplication by $`2`$. Or, slightly more precise, when passing to $`2\text{Hom}`$, the maps become surjective.
###### Question 3.15
I do not know a basis of $`\text{Hom}_{\text{F}_2𝒮_n}(S^\lambda /2,M^\mu /2)`$ in case $`\lambda `$ is $`2`$-singular.
###### Proposition 3.16
Suppose given a $`\text{Z}𝒮_n`$-linear map of the form
$$\begin{array}{ccc}\hfill S^\lambda /m& \text{}& S^\mu /m\hfill \\ \hfill a& \text{}& \underset{\{\phi \}\{\mathrm{\Phi }^{\mu ^{},\lambda ^{}}\}}{}x_{\{\phi \}}\{a^{}\}^{}\mathrm{\Theta }_\phi ^{}\nu _S^\mu ,\hfill \end{array}$$
where $`x_{\{\phi \}}\text{Z}/m`$. Neglecting to denote $`()_\text{Z}\text{Z}/m`$ for maps and using the isomorphism $`\delta `$ from (3.5), we obtain the transposition isomorphism
$$\begin{array}{ccc}\hfill \text{Hom}_{\text{Z}𝒮_n}(S^\lambda /m,S^\mu /m)& \text{}& \text{Hom}_{\text{Z}𝒮_n}(S^\mu ^{}/m,S^\lambda ^{}/m)\hfill \\ \hfill f& \text{}& f^\text{t}:=\delta ^\mu f^,(\delta ^\lambda )^1.\hfill \end{array}$$
This definition can be rephrased via the formula
$$b^{}f^\text{t}=\underset{\{\phi \}\{\mathrm{\Phi }^{\mu ^{},\lambda ^{}}\}}{}x_{\{\phi \}}b^{}\mathrm{\Theta }_{\phi ^1},$$
note that $`\mathrm{\Theta }_{\phi ^1}`$ is applicable to $`b^{}S^\mu ^{}/mM^\mu ^{}/m`$. We have $`(f^\text{t})^\text{t}=f`$.
We dualize via $`()^{}:=\text{Hom}_\text{Z}(,\text{Z}/m)`$ to obtain the following diagram.
The commutativity on top follows from (3.2). On the left and on the right, commutativity follows by dualizing the commutativity in (3.5) in the sense of (3.6). The commutativity in the back is the dualized assumption on $`f`$ and the commutativity at the bottom is the definition of $`f^\text{t}`$. The claimed formula is equivalent to the commutativity in the front face. Moreover, (3.6) allows to conclude that
$$\begin{array}{ccc}\hfill (f^\text{t})^\text{t}& =& \delta ^\lambda ^{}f^{\text{t},,}(\delta ^\mu ^{})^1\hfill \\ & =& \delta ^\lambda ^{}(\delta ^{\lambda ,,})^1f^{}\delta ^{\mu ,,}(\delta ^\mu ^{})^1\hfill \\ & =& \text{eva}f^{}\text{eva}^1\hfill \\ & =& f.\hfill \end{array}$$
###### Question 3.17 (main technical obstacle)
Given $`S^\lambda /m\text{}S^\mu /m`$ as in (3.16), I do not know a formula for $`S^\mu ^{}/m\text{}S^\lambda ^{}/m`$ in terms of a factorization over $`M^{\mu ,}/m\text{}S^\mu ^{}/m`$ of a linear combination of maps of the form $`\mathrm{\Theta }_\psi ^{}\nu _S^\lambda ^{}`$, $`\{\psi \}\{\mathrm{\Phi }^{\lambda ,\mu }\}`$. I.e., I do not know a formula ‘in terms of polytabloids instead of tabloids’.
### 3.2 Remarks on transposition in characteristic $`2`$
Given an integer $`a1`$, we denote by $`a_2:=2^{v_2(a)}`$ its $`2`$-part. From the more general assertion in (4.2) below, but nevertheless readable independently, we take that for $`b1`$ we have $`v_2(\left(\begin{array}{c}a\\ i\end{array}\right))>0`$ for all $`i[1,b]`$ if and only if $`b<a_2`$. We say that a partition $`\lambda `$ of $`n`$ is $`2`$-convergent if $`\lambda _{p+1}^{}<(\lambda _p^{}+1)_2`$ for all $`p[1,\lambda _11]`$. For instance, if $`n=8`$, the list of $`2`$-convergent partitions is $`(1^8)^{}`$, $`(3,1^5)^{}`$, $`(3^2,1^2)^{}`$, $`(5,1^3)^{}`$, $`(3^2,2)^{}`$, $`(7,1)^{}`$, $`(8)^{}`$. Let $`[x]=[\mathrm{1\hspace{0.33em}2}\mathrm{}n]`$ denote the standard $`(n)`$-tableau.
###### Lemma 3.18 (cf. \[J 78, 24.4\], \[K 99, 4.3.35\])
Let $`\lambda `$ be a partition of $`n`$. We have
$$dim\text{Hom}_{\text{F}_2𝒮_n}(S^\lambda /2,S^{(n)}/2)=\{\begin{array}{cc}1\hfill & \text{ if }\lambda \text{ is }2\text{-convergent}\hfill \\ 0\hfill & \text{ else. }\hfill \end{array}$$
In case $`\lambda `$ is $`2`$-convergent, the nonzero morphism is of the form
$$\begin{array}{ccc}\hfill S^\lambda /2& \text{}& S^{(n)}/2\hfill \\ \hfill a& \text{}& x\hfill \end{array}$$
By transposition (3.16), we conclude that
$$dim\text{Hom}_{\text{F}_2𝒮_n}(S^{(n)}/2,S^\lambda /2)=\{\begin{array}{cc}1\hfill & \text{ if }\lambda ^{}\text{ is }2\text{-convergent}\hfill \\ 0\hfill & \text{ else. }\hfill \end{array}$$
By $`𝒮_n`$-linearity, a morphism $`S^\lambda /2\text{}S^{(n)}/2`$, if existent, is necessarily of the form just given. It remains to be seen that this map is well defined if and only if $`\lambda `$ is $`2`$-convergent.
Welldefinedness may be rephrased as the existence of a factorization of the map
$$\begin{array}{ccc}\hfill F^\lambda & \text{}& S^{(n)}/2\hfill \\ \hfill [a]& \text{}& x\hfill \end{array}$$
over $`F^\lambda \text{}S^\lambda `$. A one-step Garnir relation
$$G_{[a],\xi ,\eta }^{\prime \prime }=\underset{\sigma 𝒮_\xi \times 𝒮_\eta \backslash 𝒮_{\xi \eta }}{}[a]\sigma \epsilon _\sigma ,$$
where $`\xi a_p`$ and $`\eta a_{p+1}`$ are subsets such that $`\mathrm{\#}\xi +\mathrm{\#}\eta =\lambda _p^{}+1`$ (cf. Section 1.1), is mapped to the element $`\left(\begin{array}{c}\lambda _p^{}+1\\ \mathrm{\#}\eta \end{array}\right)xS^{(n)}/2`$. Thus the map is well defined if and only if $`\left(\begin{array}{c}\lambda _p^{}+1\\ i\end{array}\right)_20`$ for all $`p[1,\lambda _11]`$ and for all $`i[1,\lambda _{p+1}^{}]`$.
###### Lemma 3.19
Let $`\lambda `$ be a partition of $`n`$. Let $`\chi ^\lambda `$ be the characteristic function of the subset $`C_{[a_\lambda ]}R_{[a_\lambda ]}𝒮_n`$, i.e. $`\chi _\sigma ^\lambda =1`$ for $`\sigma C_{[a_\lambda ]}R_{[a_\lambda ]}`$, and $`\chi _\sigma ^\lambda =0`$ for $`\sigma 𝒮_n\backslash C_{[a_\lambda ]}R_{[a_\lambda ]}`$. Let $`P^\lambda 𝒮_n`$ be the subset of permutations $`\sigma `$ for which $`[a_\lambda ]\sigma `$ is a standard tableau. We suppose given a nonzero $`\text{F}_2𝒮_n`$-linear map (unique if existent, cf. 3.18)
$$\begin{array}{ccc}\hfill S^{(n)}/2& \text{}& S^\lambda /2\hfill \\ \hfill x& \text{}& \underset{\tau P^\lambda }{}u_\tau ^\lambda a_\lambda \tau .\hfill \end{array}$$
For each $`\sigma P^\lambda `$, we obtain
$$1_2\underset{\tau P^\lambda }{}u_\tau ^\lambda \chi _{\tau \sigma ^1}^\lambda .$$
Moreover, $`(\chi _{\tau \sigma ^1}^\lambda )_{\tau \times \sigma P^\lambda \times P^\lambda }\text{GL}_{\text{rk}S^\lambda }(\text{F}_2)`$.
Note that the transpose of $`u^\lambda `$ maps each polytabloid to the nonzero element (3.18).
Disregarding alternation, we obtain on the one hand, $`\sigma 𝒮_n`$,
$$x\text{}\underset{\tau P^\lambda }{}u_\tau ^\lambda (\{a_\lambda \}\sigma ,a_\lambda \tau )$$
and on the other hand
$$x\text{}(\{x\},x)=1.$$
It remains to be shown that $`(\{a_\lambda \}\sigma ,a_\lambda \tau )_2\chi _{\tau \sigma ^1}^\lambda `$ for $`\tau ,\sigma P^\lambda `$. But $`(\{a_\lambda \}\sigma ,a_\lambda \tau )_21`$ is equivalent to the existence of $`\kappa C_{[a_\lambda ]}`$ and $`\rho R_{[a_\lambda ]}`$ such that $`[a_\lambda ]\rho =[a_\lambda ]\kappa \tau \sigma ^1`$. The invertibility of the matrix $`(\chi _{\tau \sigma ^1}^\lambda )_{\tau \times \sigma P^\lambda \times P^\lambda }`$ follows from (3.5). Cf. \[K 99, 6.2.8\].
For instance, in case $`\lambda =(3,2)`$, $`[a_\lambda ]=[\stackrel{ˇ}{a}_\lambda ]`$, and thus $`P^{(3,2)}=\{1,(45),(23),(23)(45),(2453)\}`$, this invertible matrix over $`\text{F}_2`$ turns out to be
$$(\chi _{\tau \sigma ^1}^{(3,2)})_{\tau \times \sigma P^{(3,2)}\times P^{(3,2)}}=\left[\begin{array}{ccccc}1\hfill & 0\hfill & 0\hfill & 0\hfill & 0\hfill \\ 0\hfill & 1\hfill & 0\hfill & 0\hfill & 0\hfill \\ 0\hfill & 0\hfill & 1\hfill & 0\hfill & 0\hfill \\ 0\hfill & 0\hfill & 0\hfill & 1\hfill & 0\hfill \\ 1\hfill & 0\hfill & 0\hfill & 0\hfill & 1\hfill \end{array}\right].$$
In case $`\lambda `$ is a hook partition, i.e. in case $`\lambda _i1`$ for $`i2`$, the matrix $`(\chi _{\tau \sigma ^1}^\lambda )_{\tau \times \sigma P^\lambda \times P^\lambda }`$ is the identity matrix. In particular, we recover \[K 99, 4.2.11\] by composition of a morphism as in (3.18) and the transpose, as given by (3.19), of a morphism as in (3.18). Cf. (3.31).
### 3.3 Paths as correspondoids
In order to transpose the vertical two-box-shift morphism in (2.37) via (3.16) to obtain a horizontal two-box-shift morphism, we translate the language of double paths to the language of correspondoids (3.25, 3.26) in order to be able to use our transposition formula (3.16) to obtain the provisional result (3.27, cf. 3.29). We shall use a slightly modified and generalized variant of the setup given in Section 2.1.
In this section maps are written on the right, with some exceptions made. Suppose given a partition $`\lambda `$ of $`n`$, an integer $`d1`$ and integers $`1gk\lambda _11`$ such that
$$\mu _i^{}:=\{\begin{array}{cc}\lambda _i^{}+d\hfill & \text{for }i=g\hfill \\ \lambda _i^{}d\hfill & \text{for }i=k+1\hfill \\ \lambda _i^{}\hfill & \text{else}\hfill \end{array}$$
defines a partition $`\mu `$. A weight $`e`$ is a map
$$\begin{array}{ccc}\hfill [1,\lambda _1]& \text{}& [0,d]\hfill \\ \hfill j& \text{}& e_j\hfill \end{array}$$
that maps $`g`$ and $`k+1`$ to $`e_g=e_{k+1}=d`$, and that maps $`j`$ to $`e_j=0`$ in case $`j[1,\lambda _1]\backslash [g,k+1]`$. A pattern $`\mathrm{\Xi }`$ of weight $`e`$ is a subset $`\mathrm{\Xi }[1,d]\times [g,k+1]`$ that has
$$\mathrm{\#}\left(\mathrm{\Xi }([1,d]\times \{j\})\right)=e_j$$
for $`j[g,k+1]`$. A $`d`$-fold path $`\gamma `$ of weight $`e`$ is an injection from a pattern $`\mathrm{\Xi }`$ of weight $`e`$ to $`[\lambda ][\mu ]`$ of the form
$$\begin{array}{ccc}\hfill \mathrm{\Xi }& \text{}& [\lambda ][\mu ]\hfill \\ \hfill i\times j& \text{}& \overline{\gamma }(j,i)\times j\hfill \end{array}$$
such that
$$\overline{\gamma }(g,i)=\lambda _g^{}+i\text{ for }i[1,d].$$
Sometimes, we denote its pattern by $`\mathrm{\Xi }_\gamma :=\mathrm{\Xi }`$. The set of $`d`$-fold paths of weight $`e`$ is denoted by $`\mathrm{\Gamma }(e)`$.
An ordered $`d`$-fold path of weight $`e`$ is a $`d`$-fold path $`\gamma `$ for which the application
$$\{i[1,d]|i\times j\mathrm{\Xi }_\gamma \}\text{}[1,\lambda _j^{}]$$
is increasing for each $`j[g+1,k+1]`$. The set of ordered $`d`$-fold paths of weight $`e`$ is denoted by $`\stackrel{}{\mathrm{\Gamma }}(e)`$.
Suppose given a weight $`e`$. For a $`d`$-fold path $`\gamma `$ of weight $`e`$, we let the successor permutation be defined by
$$\begin{array}{ccc}\hfill \mathrm{\Xi }_\gamma & \text{}& \mathrm{\Xi }_\gamma \hfill \\ \hfill i\times j& \text{}& \{\begin{array}{cc}i\times \mathrm{min}\{j^{}[j+1,k+1]|i\times j^{}\mathrm{\Xi }_\gamma \}\hfill & \text{ for }j[g,k]\hfill \\ i\times g\hfill & \text{ for }j=k+1.\hfill \end{array}\hfill \end{array}$$
Since $`\gamma `$ is an injection, we may define the permutation
$$[\lambda ^{}][\mu ^{}]\text{}[\lambda ^{}][\mu ^{}]$$
as the completion to a commutative diagram
in which the unlabeled arrows denote inclusions. Thus, roughly speaking, $`\widehat{\gamma }`$ is the identical prolongation of $`\text{suc}_\gamma `$ operating via $`\gamma `$. Given a $`d`$-fold path $`\gamma `$ of weight $`e`$, the permutation $`\sigma (\gamma )`$ of $`[\lambda ^{}][\mu ^{}]`$ is defined as being determined by the rule
$$\begin{array}{ccc}\hfill [\lambda ^{}][\mu ^{}]& \text{}& [\lambda ^{}][\mu ^{}]\hfill \\ \hfill j\times i& \text{}& \{\begin{array}{cccc}\hfill j& \times & i\hfill & \text{for }jk+1\hfill \\ \hfill (k+1)& \times & \overline{\gamma }(k+1,i\mu _{k+1}^{})\hfill & \text{for }j=k+1\text{ and }i[\mu _{k+1}^{}+1,\lambda _{k+1}^{}]\hfill \end{array}\hfill \end{array}$$
and by the requirement that its restriction to $`\{k+1\}\times [1,\mu _{k+1}^{}]`$ be of constant value $`k+1`$ in the first component and strictly increasing in the second component. Hence, the $`[\mu _{k+1}^{}+1,\lambda _{k+1}^{}]`$-part of the row $`k+1`$ of $`[\lambda ^{}][\mu ^{}]`$ is mapped under $`\sigma (\gamma )`$ to the image of $`\gamma \tau `$, whereas its $`[1,\mu _{k+1}^{}]`$-part is distributed, in a strictly increasing manner, over the complement of this image in that row. We define an element $`\stackrel{ˇ}{\gamma }\mathrm{\Phi }^{\mu ^{},\lambda ^{}}`$ as the completion to a commutative diagram
the unlabeled arrows denoting inclusions.
Given $`\kappa C_\lambda `$, we denote its transposition by $`\kappa ^{}R_\lambda ^{}`$, mapping $`j\times i[\lambda ^{}]`$ to $`(j\times i)\kappa ^{}:=(i\times j)\kappa \tau `$. The permutation of $`[\lambda ][\mu ]`$ that restricts to $`\kappa `$ on $`[\lambda ]`$ and to the identity on $`[\mu ]\backslash [\lambda ]`$ is denoted by $`\widehat{\kappa }`$. Similarly, the identical prolongation of $`\kappa ^{}`$ to a permutation of $`[\lambda ^{}][\mu ^{}]`$ is denoted by $`\widehat{\kappa }^{}`$. Likewise for $`\mu `$ instead of $`\lambda `$.
###### Lemma 3.20
Let $`\gamma `$ be a $`d`$-fold path of weight $`e`$, let $`\kappa C_\lambda `$. We note that the composition $`\gamma \widehat{\kappa }`$ is again a $`d`$-fold path of weight $`e`$, of the same pattern as $`\gamma `$, and obtain
$$\begin{array}{ccc}\hfill (\gamma \widehat{\kappa })\widehat{}& =& (\widehat{\kappa }^{})^1\widehat{\gamma }\widehat{\kappa }^{}\hfill \\ \hfill (\gamma \widehat{\kappa })\stackrel{ˇ}{}& =& (\kappa _0^{})^1\stackrel{ˇ}{\gamma }\kappa ^{}\hfill \end{array}$$
for some $`\kappa _0C_\mu `$.
The first equality holds since right conjugation of a cycle product in the symmetric group is performed by an application of the conjugating element to the cycle entries. To see the second equality, we consider the diagram
in which $`\kappa _0^1`$ is the completion by restriction and in which the unlabeled arrows denote inclusions. Note that $`\sigma (\gamma \widehat{\kappa })\widehat{\kappa }^1\sigma (\gamma )^1`$ restricts identically to the row k+1 of $`[\lambda ^{}][\mu ^{}]`$.
For each weight $`e`$, we fix a $`d`$-fold path $`\gamma _e`$ of weight $`e`$ and write
$$\epsilon (e):=\epsilon _{\sigma (\gamma _e)}\epsilon _{[a_\mu ^{}]^1\stackrel{ˇ}{\gamma }_e[a_\lambda ^{}]}.$$
###### Lemma 3.21
The sign $`\epsilon (e)`$ is independent of the choice of $`\gamma _e`$. More precisely, given a $`d`$-fold path $`\gamma `$ of weight $`e`$, we obtain
$$\epsilon _{\stackrel{ˇ}{\gamma }_e^1\stackrel{ˇ}{\gamma }}=\epsilon _{\sigma (\gamma _e)}\epsilon _{\sigma (\gamma )}.$$
Note that $`\epsilon _{\stackrel{ˇ}{\gamma }\stackrel{ˇ}{\gamma }_e^1}=\epsilon _{\sigma (\gamma )\widehat{\gamma }\widehat{\gamma }_e^1\sigma (\gamma _e)^1}`$ since $`\sigma (\gamma )\widehat{\gamma }\widehat{\gamma }_e^1\sigma (\gamma _e)^1`$ restricts to the identity on $`[\lambda ^{}]\backslash [\mu ^{}]`$. Moreover, we claim that $`\epsilon _{\widehat{\gamma }}`$ depends only on the weight of $`\gamma `$. First of all, it depends only on the pattern of $`\gamma `$, since this pattern determines the cycle type of $`\widehat{\gamma }`$. The sign of this cycle type in turn depends only on the cardinality $`_{j[g,k+1]}e_j`$ of that pattern, since the number of cycles of length $`2`$ that occur in $`\widehat{\gamma }`$ equals $`d`$. Altogether, after a reordering we obtain $`\epsilon _{\stackrel{ˇ}{\gamma }_e^1\stackrel{ˇ}{\gamma }}=\epsilon _{\sigma (\gamma )}\epsilon _{\sigma (\gamma _e)}`$.
###### Lemma 3.22
Given a $`d`$-fold path $`\gamma `$ of weight $`e`$ and a $`\lambda `$-tableau $`[a]`$, we obtain
$$(\{\stackrel{ˇ}{\gamma }[a^{}]\}1)\nu _S^\mu =(\stackrel{ˇ}{\gamma }[a^{}])^{}\epsilon _{\sigma (\gamma )}\epsilon (e)\epsilon _{[a]^1[a_\lambda ]}S^\mu .$$
We may conclude
$$\begin{array}{ccc}\hfill (\{\stackrel{ˇ}{\gamma }[a^{}]\}1)\nu _S^\mu & =& (\stackrel{ˇ}{\gamma }[a^{}])^{}\epsilon _{[a_\mu ^{}]^1\stackrel{ˇ}{\gamma }[a^{}]}\hfill \\ & =& (\stackrel{ˇ}{\gamma }[a^{}])^{}\epsilon _{[a_\mu ^{}]^1\stackrel{ˇ}{\gamma }_e[a_\lambda ^{}]}\epsilon _{[a_\lambda ^{}]^1[a^{}]}\epsilon _{[a^{}]^1\stackrel{ˇ}{\gamma }_e^1\stackrel{ˇ}{\gamma }[a^{}]}\hfill \\ & \stackrel{\text{(}\text{3.21}\text{)}}{=}& (\stackrel{ˇ}{\gamma }[a^{}])^{}\epsilon (e)\epsilon _{[a]^1[a_\lambda ]}\epsilon _{\sigma (\gamma )}.\hfill \end{array}$$
###### Lemma 3.23
Let $`\gamma `$ be a $`d`$-fold path of weight $`e`$. For $`\xi 𝒮_d`$ and $`j[g+1,k+1]`$ we let $`\xi `$ operate on the columns $`j`$ via
$$\begin{array}{ccc}\hfill \mathrm{\Xi }_\gamma & \text{}& [1,d]\times [g,k+1]\hfill \\ \hfill i\times j^{}& \text{}& \{\begin{array}{cccc}\hfill i& \times & j^{}\hfill & \text{ for }j^{}\{g\}[j+1,k+1]\hfill \\ \hfill (i)\xi & \times & j^{}\hfill & \text{ for }j^{}[g+1,j].\hfill \end{array}\hfill \end{array}$$
We suppose in addition that $`\xi _j`$ restricts identically to $`[1,d]\times \{j\}\backslash \mathrm{\Xi }_\gamma `$. Then the composition $`\left((\mathrm{\Xi }_\gamma )\xi _j\text{}\mathrm{\Xi }_\gamma \text{}[\lambda ][\mu ]\right)`$ is a $`d`$-fold path of weight $`e`$. Given a $`\lambda `$-tableau $`[a]`$, we obtain
$$(\{(\xi _j^1\gamma )\stackrel{ˇ}{}[a^{}]\}1)\nu _S^\mu =\{\begin{array}{cc}(\{\stackrel{ˇ}{\gamma }[a^{}]\}1)\nu _S^\mu \hfill & \text{ for }j[g+1,k]\hfill \\ (\{\stackrel{ˇ}{\gamma }[a^{}]\}1)\nu _S^\mu \epsilon _\xi \hfill & \text{ for }j=k+1.\hfill \end{array}$$
We may suppose $`\xi `$ to be a transposition, $`\xi =(s,t)`$, $`s\times j,t\times j\mathrm{\Xi }_j`$, $`st`$. In case $`j[g+1,k]`$ we obtain, composing with transpositions permuting $`[\mu ^{}]`$,
$$((s,t)_j\gamma )\stackrel{ˇ}{}=\stackrel{ˇ}{\gamma }((s\times g)\gamma ,(t\times g)\gamma )((s\times j)\gamma ,(t\times j)\gamma ),$$
yielding the required equality by two column permutations applied to the resulting polytabloid.
In case $`j=k+1`$, we obtain
$$((s,t)_{k+1}\gamma )\stackrel{ˇ}{}=\stackrel{ˇ}{\gamma }((s\times g)\gamma ,(t\times g)\gamma ),$$
yielding the required equality by a single column permutation applied to the resulting polytabloid.
Given a weight $`e`$, we write
$$\begin{array}{ccc}\hfill e!& :=& \underset{j[1,\lambda _1]\backslash \{g\}}{}e_j!\hfill \\ \hfill \lambda ^{}!& :=& \underset{j[1,\lambda _1]\backslash \{g\}}{}\lambda _j^{}!\hfill \\ \hfill (\lambda ^{}e)!& :=& \underset{j[1,\lambda _1]\backslash \{g\}}{}(\lambda _j^{}e_j)!\hfill \end{array}$$
###### Corollary 3.24
Let $`e`$ be a weight, let $`[a]`$ be a $`\lambda `$-tableau. We obtain
$$\underset{\gamma \mathrm{\Gamma }(e)}{}(\{\stackrel{ˇ}{\gamma }[a^{}]\}1)\nu _S^\mu \epsilon _{\sigma (\gamma )}=e!\underset{\gamma \stackrel{}{\mathrm{\Gamma }}(e)}{}(\{\stackrel{ˇ}{\gamma }[a^{}]\}1)\nu _S^\mu \epsilon _{\sigma (\gamma )}.$$
We rewrite this sum columnwise from right to left in that we use (3.23) to impose an ordering condition on $`\gamma `$ in column $`j`$, starting in column $`j=k+1`$ and ending in column $`j=g+1`$. So we obtain a factor $`e_j!`$ in column $`j`$.
Given a weight $`e`$, we let
$$\begin{array}{ccc}\hfill F^\lambda & \text{}& S^\mu \hfill \\ \hfill [a]& \text{}& [a]𝔣_e^{\prime \prime }:=\underset{\gamma \stackrel{}{\mathrm{\Gamma }}(e)}{}(\stackrel{ˇ}{\gamma }[a^{}])^{}\epsilon _{\sigma (\gamma )}.\hfill \end{array}$$
###### Remark 3.25 (connection to Section 2.1)
In case of $`d=2`$, we obtain that the set of ordered double paths, defined in Section 2.1, and the set of ordered $`2`$-fold paths, defined here, coincide. Moreover, $`[a]`$ being a $`\lambda `$-tableau, $`e`$ being a weight, we obtain
$$[a]f_e^{\prime \prime }=\frac{e!}{2}[a]𝔣_e^{\prime \prime },$$
the left hand side written in the notation of Section 2.1, the right hand side written in the notation introduced here.
The set $`\mathrm{\Gamma }(e)`$ is the disjoint union of $`\dot{\mathrm{\Gamma }}(e)`$ and $`(1,2)_{k+1}\dot{\mathrm{\Gamma }}(e)`$. We may rewrite the right hand side as
$$\begin{array}{ccc}\hfill \left(\underset{j[g+1,k]}{}e_j!\right)\underset{\gamma \stackrel{}{\mathrm{\Gamma }}(e)}{}(\stackrel{ˇ}{\gamma }[a^{}])^{}\epsilon _{\sigma (\gamma )}& \stackrel{\text{(}\text{3.24}\text{)}}{=}& (e_{k+1}!)^1\underset{\gamma \mathrm{\Gamma }(e)}{}(\stackrel{ˇ}{\gamma }[a^{}])^{}\epsilon _{\sigma (\gamma )}\hfill \\ & \stackrel{\text{(}\text{3.23}\text{)}}{=}& \underset{\gamma \dot{\mathrm{\Gamma }}(e)}{}(\stackrel{ˇ}{\gamma }[a^{}])^{}\epsilon _{\sigma (\gamma )}\hfill \\ & =& \underset{\gamma \dot{\mathrm{\Gamma }}(e)}{}(\stackrel{ˇ}{\gamma }[a^{}])^{}\epsilon _\gamma \hfill \\ & =& [a]f_e^{\prime \prime },\hfill \end{array}$$
where $`\epsilon _\gamma =(1)^{\overline{\gamma }(k+1,1)+\overline{\gamma }(k+1,2)}`$ as introduced in Section 2.1, and where for the last equality we translate $`[a^\gamma ]=(\stackrel{ˇ}{\gamma }[a^{}])^{}`$.
###### Proposition 3.26
Let $`e`$ be a weight. We may reformulate
$$[a]𝔣_e^{\prime \prime }=\frac{1}{e!(\lambda ^{}e)!\lambda _g^{}}\underset{\gamma \mathrm{\Gamma }(e)/C_\lambda }{}r_{\stackrel{ˇ}{\gamma }}\{a^{}\}^{}\mathrm{\Theta }_{\stackrel{ˇ}{\gamma }}^{}\nu _S^\mu \epsilon (e),$$
where $`\gamma \mathrm{\Gamma }(e)/C_\lambda `$ signifies that $`\gamma `$ runs over a set of orbit representatives of $`\mathrm{\Gamma }(e)`$ under the operation of $`C_\lambda `$ in the sense of (3.20). In particular, there is a factorization
$$(F^\lambda \text{}S^\mu )=(F^\lambda \text{}M^{\lambda ,}\text{}S^\mu ).$$
Given a $`d`$-fold path $`\gamma `$ of weight $`e`$, the stabilizer $`\{\kappa C_\lambda |\gamma \widehat{\kappa }=\gamma \}`$ of $`\gamma `$ in $`C_\lambda `$ has cardinality $`(\lambda ^{}e)!\lambda _g^{}`$. So we calculate
$$\begin{array}{ccc}\hfill [a]𝔣_e^{\prime \prime }& =& \underset{\gamma \stackrel{}{\mathrm{\Gamma }}(e)}{}(\stackrel{ˇ}{\gamma }[a^{}])^{}\epsilon _{\sigma (\gamma )}\hfill \\ & \stackrel{\text{(}\text{3.24}\text{)}}{=}& \frac{1}{e!}\underset{\gamma \mathrm{\Gamma }(e)}{}(\stackrel{ˇ}{\gamma }[a^{}])^{}\epsilon _{\sigma (\gamma )}\hfill \\ & =& \frac{1}{e!(\lambda ^{}e)!\lambda _g^{}}\underset{\gamma \mathrm{\Gamma }(e)/C_\lambda }{}r_{\stackrel{ˇ}{\gamma }}r_{\stackrel{ˇ}{\gamma }}^1\underset{\kappa C_\lambda }{}(\gamma \widehat{\kappa })\stackrel{ˇ}{}[a^{}])^{}\epsilon _{\sigma (\gamma )}\hfill \\ & \stackrel{\text{(}\text{3.22}\text{)}}{=}& \frac{1}{e!(\lambda ^{}e)!\lambda _g^{}}\underset{\gamma \mathrm{\Gamma }(e)/C_\lambda }{}r_{\stackrel{ˇ}{\gamma }}r_{\stackrel{ˇ}{\gamma }}^1\underset{\kappa C_\lambda }{}(\{(\gamma \widehat{\kappa })\stackrel{ˇ}{}[a^{}]\}1)\nu _S^\mu \epsilon (e)\epsilon _{[a]^1[a_\lambda ]}\hfill \end{array}$$
$$\begin{array}{ccc}& \stackrel{\text{(}\text{3.20}\text{)}}{=}& \frac{1}{e!(\lambda ^{}e)!\lambda _g^{}}\underset{\gamma \mathrm{\Gamma }(e)/C_\lambda }{}r_{\stackrel{ˇ}{\gamma }}\left(\left(r_{\stackrel{ˇ}{\gamma }}^1\underset{\kappa ^{}R_\lambda ^{}}{}\{\stackrel{ˇ}{\gamma }\kappa ^{}[a^{}]\}\right)1\right)\nu _S^\mu \epsilon (e)\epsilon _{[a]^1[a_\lambda ]}\hfill \\ & \stackrel{\text{(}\text{3.1}\text{)}}{=}& \frac{1}{e!(\lambda ^{}e)!\lambda _g^{}}\underset{\gamma \mathrm{\Gamma }(e)/C_\lambda }{}r_{\stackrel{ˇ}{\gamma }}\left(\{a^{}\}\mathrm{\Theta }_{\stackrel{ˇ}{\gamma }}1\right)\nu _S^\mu \epsilon (e)\epsilon _{[a]^1[a_\lambda ]}\hfill \\ & =& \frac{1}{e!(\lambda ^{}e)!\lambda _g^{}}\underset{\gamma \mathrm{\Gamma }(e)/C_\lambda }{}r_{\stackrel{ˇ}{\gamma }}\{a^{}\}^{}\mathrm{\Theta }_{\stackrel{ˇ}{\gamma }}^{}\nu _S^\mu \epsilon (e).\hfill \end{array}$$
###### Corollary 3.27 (to 2.37 via 3.26)
Let $`d=2`$. The transpose in the sense of (3.16) of the morphism given in (2.37) is obtained as a factorization
where
$$[b^{}]f^{\text{t},0}:=\underset{eE}{}\frac{1}{(\lambda ^{}e)!2R\lambda _g^{}}\left(\underset{i[g+1,k]}{}X_i^{(2e_i)}\right)\underset{\gamma \mathrm{\Gamma }(e)/C_\lambda }{}r_{\stackrel{ˇ}{\gamma }}b^{}\mathrm{\Theta }_{\stackrel{ˇ}{\gamma }^1}\epsilon (e),$$
and where $`f`$, $`E`$, $`R`$, $`m`$ and $`X_i`$ are defined as for (2.37).
Given $`t,T2`$, Z-lattices $`X`$ and $`Y`$ and a Z-linear map $`X/t\text{}Y/t`$, we denote
$$(X/(tT)\text{}Y/(tT)):=(X/(tT)\text{}X/t\text{}Y/t\text{}Y/(tT)).$$
Let $`M:=R\lambda ^{}!2\lambda _g^{}`$ play the role of a large enough integer and consider the morphism
$$\begin{array}{ccc}\hfill S^\lambda /(mM)& \text{}& S^\mu /(mM)\hfill \\ \hfill a& \text{}& \lambda ^{}!2\lambda _g^{}\underset{eE}{}\left(\underset{i[g+1,k]}{}X_i^{(2e_i)}\right)\{a^{}\}^{}f_e^{}\hfill \\ & \stackrel{\text{(}\text{3.25}\text{}\text{3.26}\text{)}}{=}& \lambda ^{}!\lambda _g^{}\underset{eE}{}\left(\underset{i[g+1,k]}{}X_i^{(2e_i)}\right)\{a^{}\}^{}𝔣_e^{}\hfill \\ & \stackrel{\text{(}\text{3.26}\text{)}}{=}& \underset{eE}{}\frac{\lambda ^{}!}{(\lambda ^{}e)!}\left(\underset{i[g+1,k]}{}X_i^{(2e_i)}\right)\underset{\gamma \mathrm{\Gamma }(e)/C_\lambda }{}r_{\stackrel{ˇ}{\gamma }}\{a^{}\}^{}\mathrm{\Theta }_{\stackrel{ˇ}{\gamma }}^{}\nu _S^\mu \epsilon (e).\hfill \end{array}$$
For $`\phi \mathrm{\Phi }^{\mu ^{},\lambda ^{}}`$, the coefficient $`x_{\{\phi \}}`$ in the sense of (3.16) can be written as
$$x_{\{\phi \}}=\underset{eE}{}\frac{\lambda ^{}!}{(\lambda ^{}e)!}(\underset{i[g+1,k]}{}X_i^{(2e_i)})\underset{\gamma \mathrm{\Gamma }(e)/C_\lambda }{}r_{\stackrel{ˇ}{\gamma }}\left\{\begin{array}{cc}1\hfill & \text{ for }\left\{\stackrel{ˇ}{\gamma }\right\}=\left\{\phi \right\}\hfill \\ 0\hfill & \text{ else }\hfill \end{array}\right\}\epsilon (e).$$
Thus, by (3.16), the transpose of $`S^\lambda /(mM)\text{}S^\mu /(mM)`$ maps
$$\begin{array}{ccc}\hfill S^\mu ^{}/(mM)& \hfill \text{}& S^\lambda ^{}/(mM)\hfill \\ \hfill b^{}& \hfill \text{}& \underset{eE}{}\frac{\lambda ^{}!}{(\lambda ^{}e)!}(\underset{i[g+1,k]}{}X_i^{(2e_i)})\underset{\gamma \mathrm{\Gamma }(e)/C_\lambda }{}r_{\stackrel{ˇ}{\gamma }}b^{}\mathrm{\Theta }_{\stackrel{ˇ}{\gamma }^1}\epsilon (e)\hfill \end{array}$$
Finally, there is a commutative diagram
with vertical injections given by multiplication by $`M`$, so that $`f^\text{t}M=(fM)^\text{t}`$.
###### Remark 3.28
Consider the set $`\stackrel{}{\mathrm{\Gamma }}_0(e)`$ of ordered $`d`$-fold paths $`\gamma `$ of weight $`e`$ that have $`\overline{\gamma }(j,i)[\lambda _j^{}e_j,\lambda _j^{}]`$ for $`i\times j\mathrm{\Xi }_{[g+1,k+1]}`$, i.e. that have their positions sitting at the bottoms of the columns. We say that a correspondoid $`\{\phi \}\{\mathrm{\Phi }^{\mu ^{},\lambda ^{}}\}`$ is semistandard up to row permutation if there exists a row permutation $`\rho R_\mu ^{}`$ such that $`\{\rho \phi \}\{\mathrm{\Phi }^{\mu ^{},\lambda ^{}}\}_{\text{sst}}`$ becomes semistandard. Note that $`\mathrm{\Theta }_\phi =\mathrm{\Theta }_{\rho \phi }`$ by (3.1). For $`\gamma \stackrel{}{\mathrm{\Gamma }}_0(e)`$, the correspondoid $`\{\stackrel{ˇ}{\gamma }^1\}`$ is semistandard up to row permutation as long as $`\mathrm{\Xi }_\gamma `$ is not bulky in the sense introduced at the beginning of Section 2.3. Cf. the bases in (3.13), cf. the reduced tuple of coefficients employed in (2.37).
Note that $`\stackrel{}{\mathrm{\Gamma }}_0(e)`$ is a set of representatives of $`\mathrm{\Gamma }(e)/C_\lambda `$.
###### Question 3.29
I do not know a formula for the transpose $`f^\text{t}`$ given in (3.27) in terms of $`\lambda ^{}`$-polytabloids. Cf. (3.17).
### 3.4 Horizontal examples
We shall transpose some examples of Section 2.4 with respect to the convention that $`[a_\lambda ]=[\stackrel{ˇ}{a}_\lambda ]`$ for all partitions $`\lambda `$, cf. Section 1.1. Again, we omit to denote the brackets that indicate polytabloids, moreover, we omit commas in cycles. If possible (cf. 3.13), we sort the resulting linear combination of polytabloids into images under $`\mathrm{\Theta }_\phi ^{}\nu _S^\lambda ^{}`$ for various (implicitely present) correspondoids $`\{\phi \}`$. The formula in (3.27) does not suffice to give the coefficients of this expression (cf. 3.17, 4.15).
###### Example 3.30
Let $`\lambda =(3,3)`$, $`\mu =(2,2,1,1)`$. The transpose of the morphism given in (2.40) is
$$\begin{array}{ccc}\hfill S^\mu ^{}/4& \text{}& S^\lambda ^{}/4\hfill \\ \hfill \begin{array}{cccc}1\hfill & 3\hfill & 5\hfill & 6\hfill \\ 2\hfill & 4\hfill & & \end{array}& \text{}& \begin{array}{cc}1\hfill & 4\hfill \\ 2\hfill & 5\hfill \\ 3\hfill & 6\hfill \end{array}\left(1\left(34\right)\right)+2\begin{array}{cc}1\hfill & 3\hfill \\ 2\hfill & 4\hfill \\ 5\hfill & 6\hfill \end{array}.\hfill \end{array}$$
###### Example 3.31
Let $`\lambda =(3,2,2)`$, $`\mu =(3,1,1,1,1)`$. The transposes of the morphisms given in (2.41) are
$$\begin{array}{ccc}\hfill S^\mu ^{}/3& \text{}& S^\lambda ^{}/3\hfill \\ \hfill \begin{array}{ccccc}1\hfill & 4\hfill & 5\hfill & 6\hfill & 7\hfill \\ 2\hfill & & & & \\ 3\hfill & & & & \end{array}& \text{}& \begin{array}{ccc}1\hfill & 4\hfill & 6\hfill \\ 2\hfill & 5\hfill & 7\hfill \\ 3\hfill & & \end{array}\begin{array}{ccc}1\hfill & 4\hfill & 5\hfill \\ 2\hfill & 6\hfill & 7\hfill \\ 3\hfill & & \end{array}.\hfill \end{array}$$
and
$$\begin{array}{ccc}\hfill S^\mu ^{}/2& \text{}& S^\lambda ^{}/2\hfill \\ \hfill \begin{array}{ccccc}1\hfill & 4\hfill & 5\hfill & 6\hfill & 7\hfill \\ 2\hfill & & & & \\ 3\hfill & & & & \end{array}& \text{}& \begin{array}{ccc}1\hfill & 2\hfill & 6\hfill \\ 3\hfill & 4\hfill & 7\hfill \\ 5\hfill & & \end{array}+\begin{array}{ccc}1\hfill & 4\hfill & 5\hfill \\ 2\hfill & 6\hfill & 7\hfill \\ 3\hfill & & \end{array}+\begin{array}{ccc}1\hfill & 3\hfill & 5\hfill \\ 2\hfill & 6\hfill & 7\hfill \\ 4\hfill & & \end{array}+\begin{array}{ccc}1\hfill & 2\hfill & 5\hfill \\ 3\hfill & 6\hfill & 7\hfill \\ 4\hfill & & \end{array}+\begin{array}{ccc}1\hfill & 3\hfill & 4\hfill \\ 2\hfill & 6\hfill & 7\hfill \\ 5\hfill & & \end{array}\hfill \\ & & +\begin{array}{ccc}1\hfill & 2\hfill & 4\hfill \\ 3\hfill & 6\hfill & 7\hfill \\ 5\hfill & & \end{array}+\begin{array}{ccc}1\hfill & 2\hfill & 3\hfill \\ 4\hfill & 6\hfill & 7\hfill \\ 5\hfill & & \end{array}+\begin{array}{ccc}1\hfill & 2\hfill & 5\hfill \\ 3\hfill & 4\hfill & 7\hfill \\ 6\hfill & & \end{array}+\begin{array}{ccc}1\hfill & 3\hfill & 4\hfill \\ 2\hfill & 5\hfill & 7\hfill \\ 6\hfill & & \end{array}+\begin{array}{ccc}1\hfill & 2\hfill & 4\hfill \\ 3\hfill & 5\hfill & 7\hfill \\ 6\hfill & & \end{array}\hfill \\ & & +\begin{array}{ccc}1\hfill & 2\hfill & 3\hfill \\ 4\hfill & 5\hfill & 7\hfill \\ 6\hfill & & \end{array}+\begin{array}{ccc}1\hfill & 2\hfill & 5\hfill \\ 3\hfill & 4\hfill & 6\hfill \\ 7\hfill & & \end{array}+\begin{array}{ccc}1\hfill & 3\hfill & 4\hfill \\ 2\hfill & 5\hfill & 6\hfill \\ 7\hfill & & \end{array}+\begin{array}{ccc}1\hfill & 2\hfill & 4\hfill \\ 3\hfill & 5\hfill & 6\hfill \\ 7\hfill & & \end{array}+\begin{array}{ccc}1\hfill & 2\hfill & 3\hfill \\ 4\hfill & 5\hfill & 6\hfill \\ 7\hfill & & \end{array}.\hfill \end{array}$$
Let us consider the morphism $`u^{\mu ,\text{t}}u^\lambda ^{}`$. Using $`[a_\lambda ^{}]=[\stackrel{ˇ}{a}_\lambda ^{}]`$, we obtain, in the notation of (3.19),
$$\begin{array}{c}P^\lambda ^{}=\{1,(56),(34),(34)(56),(354),(3564),(3654),(364),(37654),(3764),(234),(234)(56),\hfill \\ (2354),(23564),(23654),(2364),(237654),(23764),(24)(356),(24)(36),(24)(376)\}𝒮_7.\hfill \end{array}$$
Hence, with respect to the ordering of $`P^\lambda ^{}`$ just indicated, the matrix $`(\chi _{\tau \sigma ^1}^\lambda ^{})_{\tau \times \sigma P^\lambda ^{}\times P^\lambda ^{}}`$ takes the form
$$\left[\begin{array}{ccccccccccccccccccccc}\hfill 1& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 1& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 1& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill 1& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 1& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 1& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 1& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 1& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 1& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 1& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 1& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 1& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0\\ \hfill 1& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 1& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0\\ \hfill 1& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 1& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 1& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 1& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 1& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 1& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0\\ \hfill 1& \hfill 1& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 1& \hfill 0& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 1& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 1& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 1& \hfill 0& \hfill 1& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 1& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill 1& \hfill 0& \hfill 0& \hfill 1& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 1& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill 1& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 1& \hfill 0& \hfill 1& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill 1\end{array}\right].$$
By labelling our morphism as $`u^{\mu ,\text{t}}u^\lambda ^{}`$, we implicitely claimed the tuple $`(u_\tau ^\lambda ^{})_{\tau P^\lambda ^{}}`$ to be represented by the vector
$$[0,1,0,1,0,1,0,1,0,1,0,1,1,1,1,1,1,1,1,1,1].$$
But this row vector yields a row vector of constant entry $`1`$ by right multiplication with that matrix, which confirms this claim.
Furthermore, as remarked after (3.19), the matrix $`(\chi _{\tau \sigma ^1}^\mu )_{\tau \times \sigma P^\mu \times P^\mu }`$ equals the identity matrix.
###### Example 3.32
Let $`\lambda =(3,3,1,1)`$, $`\mu =(2,2,1,1,1,1)`$. The transpose of the morphism given in (2.42) is
$$\begin{array}{ccc}\hfill S^\mu ^{}/6& \text{}& S^\lambda ^{}/6\hfill \\ \hfill \begin{array}{cccccc}1\hfill & 3\hfill & 5\hfill & 6\hfill & 7\hfill & 8\hfill \\ 2\hfill & 4\hfill & & & & \end{array}& \text{}& \begin{array}{cccc}1\hfill & 4\hfill & 7\hfill & 8\hfill \\ 2\hfill & 5\hfill & & \\ 3\hfill & 6\hfill & & \end{array}\left(1\left(34\right)\right)2\begin{array}{cccc}1\hfill & 3\hfill & 7\hfill & 8\hfill \\ 2\hfill & 4\hfill & & \\ 5\hfill & 6\hfill & & \end{array}\hfill \\ & & \begin{array}{cccc}1\hfill & 4\hfill & 6\hfill & 8\hfill \\ 2\hfill & 5\hfill & & \\ 3\hfill & 7\hfill & & \end{array}\left(1\left(34\right)\right)2\begin{array}{cccc}1\hfill & 3\hfill & 6\hfill & 8\hfill \\ 2\hfill & 4\hfill & & \\ 5\hfill & 7\hfill & & \end{array}\hfill \\ & & \begin{array}{cccc}1\hfill & 4\hfill & 5\hfill & 8\hfill \\ 2\hfill & 6\hfill & & \\ 3\hfill & 7\hfill & & \end{array}\left(1\left(34\right)\right)2\begin{array}{cccc}1\hfill & 3\hfill & 5\hfill & 8\hfill \\ 2\hfill & 4\hfill & & \\ 6\hfill & 7\hfill & & \end{array}\hfill \\ & & \begin{array}{cccc}1\hfill & 4\hfill & 6\hfill & 7\hfill \\ 2\hfill & 5\hfill & & \\ 3\hfill & 8\hfill & & \end{array}\left(1\left(34\right)\right)2\begin{array}{cccc}1\hfill & 3\hfill & 6\hfill & 7\hfill \\ 2\hfill & 4\hfill & & \\ 5\hfill & 8\hfill & & \end{array}\hfill \\ & & \begin{array}{cccc}1\hfill & 4\hfill & 5\hfill & 7\hfill \\ 2\hfill & 6\hfill & & \\ 3\hfill & 8\hfill & & \end{array}\left(1\left(34\right)\right)2\begin{array}{cccc}1\hfill & 3\hfill & 5\hfill & 7\hfill \\ 2\hfill & 4\hfill & & \\ 6\hfill & 8\hfill & & \end{array}\hfill \\ & & \begin{array}{cccc}1\hfill & 4\hfill & 5\hfill & 6\hfill \\ 2\hfill & 7\hfill & & \\ 3\hfill & 8\hfill & & \end{array}\left(1\left(34\right)\right)2\begin{array}{cccc}1\hfill & 3\hfill & 5\hfill & 6\hfill \\ 2\hfill & 4\hfill & & \\ 7\hfill & 8\hfill & & \end{array}.\hfill \end{array}$$
###### Example 3.33
Let $`\lambda =(4,4)`$, $`\mu =(3,3,1,1)`$. The transpose of the morphism given in (2.41) is
$$\begin{array}{ccc}\hfill S^\mu ^{}/5& \text{}& S^\lambda ^{}/5\hfill \\ \hfill \begin{array}{cccc}1\hfill & 4\hfill & 7\hfill & 8\hfill \\ 2\hfill & 5\hfill & & \\ 3\hfill & 6\hfill & & \end{array}& \text{}& \begin{array}{cc}1\hfill & 5\hfill \\ 2\hfill & 6\hfill \\ 3\hfill & 7\hfill \\ 4\hfill & 8\hfill \end{array}\left(1\left(45\right)\left(46\right)\right)3\begin{array}{cc}1\hfill & 4\hfill \\ 2\hfill & 5\hfill \\ 3\hfill & 6\hfill \\ 7\hfill & 8\hfill \end{array}.\hfill \end{array}$$
###### Example 3.34
Let $`\lambda =(3,3,2)`$, $`\mu =(2,2,2,1,1)`$. The transpose of the morphism given in (2.41) is
$$\begin{array}{ccc}\hfill S^\mu ^{}/5& \text{}& S^\lambda ^{}/5\hfill \\ \hfill \begin{array}{ccccc}1\hfill & 3\hfill & 5\hfill & 7\hfill & 8\hfill \\ 2\hfill & 4\hfill & 6\hfill & & \end{array}& \text{}& \begin{array}{ccc}1\hfill & 4\hfill & 7\hfill \\ 2\hfill & 5\hfill & 8\hfill \\ 3\hfill & 6\hfill & \end{array}\left(1\left(34\right)\right)2\begin{array}{ccc}1\hfill & 4\hfill & 6\hfill \\ 2\hfill & 5\hfill & 8\hfill \\ 3\hfill & 7\hfill & \end{array}\left(1\left(34\right)\right)\left(1\left(56\right)\right)\hfill \\ & & \begin{array}{ccc}1\hfill & 3\hfill & 6\hfill \\ 2\hfill & 4\hfill & 8\hfill \\ 5\hfill & 7\hfill & \end{array}\left(1\left(56\right)\right)+\begin{array}{ccc}1\hfill & 4\hfill & 6\hfill \\ 2\hfill & 5\hfill & 7\hfill \\ 3\hfill & 8\hfill & \end{array}\left(1\left(34\right)\right)\left(1\left(56\right)\right)\hfill \\ & & \begin{array}{ccc}1\hfill & 3\hfill & 6\hfill \\ 2\hfill & 4\hfill & 7\hfill \\ 5\hfill & 8\hfill & \end{array}\left(1\left(56\right)\right)+2\begin{array}{ccc}1\hfill & 3\hfill & 5\hfill \\ 2\hfill & 4\hfill & 6\hfill \\ 7\hfill & 8\hfill & \end{array}.\hfill \end{array}$$
###### Example 3.35
Let $`\lambda =(3,3,1,1)`$, $`\mu =(2,2,2,2)`$. The transpose of the morphism given in (2.45) is
$$\begin{array}{ccc}\hfill S^\mu ^{}/3& \text{}& S^\lambda ^{}/3\hfill \\ \hfill \begin{array}{cccc}1\hfill & 3\hfill & 5\hfill & 7\hfill \\ 2\hfill & 4\hfill & 6\hfill & 8\hfill \end{array}& \text{}& \begin{array}{cccc}1\hfill & 4\hfill & 6\hfill & 8\hfill \\ 2\hfill & 5\hfill & & \\ 3\hfill & 7\hfill & & \end{array}\left(1\left(34\right)\right)\left(1\left(56\right)\right)\left(1\left(78\right)\right)\hfill \\ & & \begin{array}{cccc}1\hfill & 3\hfill & 6\hfill & 8\hfill \\ 2\hfill & 4\hfill & & \\ 5\hfill & 7\hfill & & \end{array}\left(1\left(56\right)\right)\left(1\left(78\right)\right).\hfill \end{array}$$
###### Remark 3.36
The morphisms in (2.46, 2.47) are transposed to their respective negative, $`g^\text{t}=g`$. For instance, the morphism $`S^{(4,3)}/4\text{}S^{(2,2,2,1)}/4`$ given in (2.46) can be obtained as
$$\nu _S^\lambda g_4(2\mathrm{\Theta }_{\phi _1}^{}\mathrm{\Theta }_{\phi _2}^{}+2\mathrm{\Theta }_{\phi _3}^{})\nu _S^\mu ,$$
$$\{\phi _1\}:=\left\{\begin{array}{cc}1\hfill & 3\hfill \\ 1\hfill & 3\hfill \\ 2\hfill & 4\hfill \\ 2\hfill & \end{array}\right\},\{\phi _2\}:=\left\{\begin{array}{cc}1\hfill & 2\hfill \\ 1\hfill & 3\hfill \\ 2\hfill & 4\hfill \\ 3\hfill & \end{array}\right\},\{\phi _3\}:=\left\{\begin{array}{cc}1\hfill & 2\hfill \\ 1\hfill & 2\hfill \\ 3\hfill & 4\hfill \\ 3\hfill & \end{array}\right\}.$$
which allows to apply (3.16) directly and to compare linear combinations of tabloids.
## 4 Two columns, several boxes
We shall generalize (2.17), treating the case of essentially two columns and two shifted boxes, to the case of essentially two columns and an arbitrary number of shifted boxes.
###### Lemma 4.1
Let $`x,y0`$. As polynomials in the free variable $`t`$, we obtain
$$\underset{j[0,x]}{}(1)^{xj}\left(\begin{array}{c}y+j\\ j\end{array}\right)\left(\begin{array}{c}t\\ xj\end{array}\right)=\left(\begin{array}{c}yt+x\\ x\end{array}\right)\text{Q}[t],$$
where
$$\left(\begin{array}{c}f\left(t\right)\\ z\end{array}\right):=z!^1\underset{i[0,z1]}{}(f(t)i)$$
for $`f(t)\text{Q}[t]`$, $`z0`$. So, in particular, $`\left(\begin{array}{c}f\left(t\right)\\ 0\end{array}\right)=1`$.
We name the claimed equality $`\text{Eq}(x,y)`$ and shall prove the following assertions.
* $`\text{Eq}(0,y)`$ holds for $`y0`$.
* $`\text{Eq}(x,0)`$ holds for $`x0`$.
* $`\text{Eq}(x,y1)`$ and $`\text{Eq}(x1,y)`$ together imply $`\text{Eq}(x,y)`$ for $`x,y1`$.
Ad (i). $`\text{Eq}(0,y)`$ writes $`1=1`$.
Ad (ii). $`\text{Eq}(x,0)`$ writes
$$\underset{j[0,x]}{}(1)^j\left(\begin{array}{c}t\\ j\end{array}\right)=\left(\begin{array}{c}xt\\ x\end{array}\right).$$
Proceeding by induction on $`x`$ and using (i) we are reduced to consider differences, i.e. to see that
$$(1)^x\left(\begin{array}{c}t\\ x\end{array}\right)=\left(\begin{array}{c}\left(tx+1\right)\\ x\end{array}\right).$$
Ad (iii). The right hand side of $`\text{Eq}(x,y)`$ equals
$$\begin{array}{ccc}\hfill \left(\begin{array}{c}yt+x\\ x\end{array}\right)& =& \left(\begin{array}{c}yt+\left(x1\right)\\ \left(x1\right)\end{array}\right)+\left(\begin{array}{c}\left(y1\right)t+x\\ x\end{array}\right)\hfill \\ & \stackrel{\text{by assumption}}{=}& \underset{j[0,x1]}{}(1)^{x1j}\left(\begin{array}{c}y+j\\ j\end{array}\right)\left(\begin{array}{c}t\\ x1j\end{array}\right)+\underset{j[0,x]}{}(1)^{xj}\left(\begin{array}{c}y1+j\\ j\end{array}\right)\left(\begin{array}{c}t\\ xj\end{array}\right)\hfill \\ & =& \underset{j[1,x]}{}(1)^{xj}\left(\begin{array}{c}y+j1\\ j1\end{array}\right)\left(\begin{array}{c}t\\ xj\end{array}\right)+\underset{j[0,x]}{}(1)^{xj}\left(\begin{array}{c}y+j1\\ j\end{array}\right)\left(\begin{array}{c}t\\ xj\end{array}\right)\hfill \\ & =& \underset{j[0,x]}{}(1)^{xj}\left(\begin{array}{c}y+j\\ j\end{array}\right)\left(\begin{array}{c}t\\ xj\end{array}\right),\hfill \end{array}$$
being the left hand side.
###### Lemma 4.2
Let $`1kx`$. We have
$$\mathrm{gcd}(\left(\begin{array}{c}x\\ 1\end{array}\right),\left(\begin{array}{c}x\\ 2\end{array}\right),\mathrm{},\left(\begin{array}{c}x\\ k\end{array}\right))=x\underset{p\text{ prime, }p|x}{}p^{\mathrm{min}(v_p(x),\text{ilog}_p(k))},$$
where
$$\text{ilog}_p(k):=\mathrm{max}\{i0|p^ik\}=\mathrm{max}v_p([1,k]).$$
Let $`p`$ be a prime, let $`j[0,x]`$. For an integer $`w0`$, written $`p`$-adically as $`w={\displaystyle \underset{i0}{}}w_ip^i`$, $`w_i[0,p1]`$, we denote its Quersumme by $`q_p(w):={\displaystyle \underset{i0}{}}w_i`$. We obtain
$$v_p(j!)=\underset{i0}{}w_i\frac{p^i1}{p1}=\frac{jq_p(j)}{p1},$$
and thus
$$v_p(\left(\begin{array}{c}x\\ j\end{array}\right))=(p1)^1(q_p(xj)(q_p(x)q_p(j))).$$
In particular, in case $`jp^{v_p(x)}`$ we get $`q_p(xj)=(q_p(x)1)+(p1)(v_p(x)v_p(j)1)+p1q_p(j)`$, whence
$$v_p(\left(\begin{array}{c}x\\ j\end{array}\right))=v_p(x)v_p(j).$$
Thus the minimal such valuation for $`j[1,k]`$ takes the value claimed above.
Let $`n1`$, $`d1`$. Assume given a partition $`\lambda `$ of $`n`$ and an integer $`g[1,\lambda _11]`$ such that
$$\mu _i^{}:=\{\begin{array}{cc}\lambda _i^{}+d\hfill & \text{for }i=g\hfill \\ \lambda _i^{}d\hfill & \text{for }i=g+1\hfill \\ \lambda _i^{}\hfill & \text{else}\hfill \end{array}$$
defines a partition $`\mu `$. Let $`Z^{}`$ be the set of injective maps
$$[1,d]\text{}[1,\lambda _{g+1}^{}].$$
Let $`ZZ^{}`$ be the subset of strictly monotone maps. Let the sign of $`\zeta Z`$ be given by $`\epsilon _\zeta :=(1)^{_{i[1,d]}\zeta (i)}`$. For $`\zeta ^{}Z^{}`$ there is a unique factorization $`\zeta ^{}(i)=\stackrel{}{\zeta }^{}(i\sigma )`$, $`i[1,d]`$, with $`\sigma 𝒮_d`$ and $`\stackrel{}{\zeta }^{}Z`$. Let the sign of $`\zeta ^{}Z^{}`$ be given by $`\epsilon _\zeta ^{}:=\epsilon _\sigma \epsilon _\stackrel{}{\zeta }^{}`$. Suppose given a $`\lambda `$-tableau $`[a]`$ and an injection $`\zeta ^{}Z^{}`$. We let
$$\begin{array}{ccc}\hfill [0,\lambda _{g+1}^{}]& \text{}& [0,\mu _{g+1}^{}]\hfill \\ \hfill i& \text{}& \mathrm{\#}\left([1,i]\backslash \zeta ^{}([1,d])\right)\hfill \\ \hfill \mathrm{min}(\phi ^1(\{j\}))& \text{}& j\hfill \end{array}$$
and define the $`\mu `$-tableau $`[a^\zeta ^{}]`$ by
$$\begin{array}{cccc}\hfill a_{j,i}^\zeta ^{}& :=& a_{j,i}\hfill & \text{for }(j[1,\lambda _1]\backslash \{g,g+1\}\text{ and }i[1,\mu _j^{}])\text{ or }(j=g\text{ and }i[1,\lambda _g^{}])\hfill \\ \hfill a_{g,\lambda _g^{}+i}^\zeta ^{}& :=& a_{g+1,\zeta ^{}(i)}\hfill & \text{for }i[1,d]\hfill \\ \hfill a_{g+1,i}^\zeta ^{}& :=& a_{g+1,\psi (i)}\hfill & \text{for }i[1,\mu _{g+1}^{}].\hfill \end{array}$$
Using this place operation, we define the $`\text{Z}𝒮_n`$-linear map
$$\begin{array}{ccc}\hfill F^\lambda & \text{}& S^\mu \hfill \\ \hfill [a]& \text{}& \underset{\zeta Z}{}a^\zeta \epsilon _\zeta =\frac{1}{d!}\underset{\zeta ^{}Z^{}}{}a^\zeta ^{}\epsilon _\zeta ^{}.\hfill \end{array}$$
For a $`\lambda `$-tableau $`[a]`$ and for an injection $`\zeta ^{}Z^{}`$, we denote the ‘$`[a]`$-realization’ of $`\zeta ^{}`$ by
$$([1,d]\text{}[1,n]):=\left([1,d]\text{}[1,\lambda _{g+1}^{}]\text{}([1,\lambda _{g+1}^{}]\times \{g+1\})\text{}[\lambda ]\text{}[1,n]\right).$$
###### Lemma 4.3
We have a factorization
$$(F^\lambda \text{}S^\mu )=(F^\lambda \text{}M^{\lambda ^{},}\text{}S^\mu ).$$
We need to show that given $`v,wa_{g+1}`$, $`vw`$, and a $`\lambda `$-tableau $`[a]`$, we have $`([a](v,w))f^{\prime \prime }=[a]f^{\prime \prime }`$. Suppose given $`\zeta Z`$. In case $`v,w\zeta _{[a]}([1,d])`$ or in case $`v,w\zeta _{[a]}([1,d])`$, we obtain $`a^\zeta \epsilon _\zeta (v,w)=a^\zeta \epsilon _\zeta `$. The maps $`\zeta Z`$ that satisfy $`v\zeta _{[a]}([1,d])`$, $`w\zeta _{[a]}([1,d])`$ furnish a subset $`Z_{10}Z`$, the maps $`\zeta Z`$ that satisfy $`v\zeta _{[a]}([1,d])`$, $`w\zeta _{[a]}([1,d])`$ furnish a subset $`Z_{01}Z`$. We have a bijection $`\iota `$ that sends $`\zeta Z_{10}`$ to the map $`\iota (\zeta )Z_{01}`$ determined by $`\iota (\zeta )_{[a]}\left([1,d]\right)=(\zeta _{[a]}\left([1,d]\right))(v,w)`$. We claim that
$$a^\zeta \epsilon _\zeta (v,w)=a^{\iota (\zeta )}\epsilon _{\iota (\zeta )},$$
thus proving the lemma. Let $`v=:a_{g+1,i}`$, let $`w=:a_{g+1,j}`$ and assume $`i<j`$. Comparing both sides, the entries in column $`g`$ of $`[\mu ]`$ yield a sign $`(1)^{\mathrm{\#}([i+1,j1]\zeta ([1,d]))}`$, the entries in column $`g+1`$ of $`[\mu ]`$ yield a sign $`(1)^{\mathrm{\#}([i+1,j1]\backslash \zeta ([1,d]))}`$, so, altogether, the entries yield a sign $`(1)^{ji1}`$. On the other hand, the quotient of the signs of the maps is $`\epsilon _{\iota (\zeta )}/\epsilon _\zeta =(1)^{ji}`$.
###### Remark 4.4
Let the correspondence $`\phi \mathrm{\Phi }^{\lambda ^{},\mu ^{}}`$ be defined by
$$\begin{array}{cccc}\hfill [\lambda ^{}]& \text{}& [\mu ^{}]\hfill & \\ \hfill j\times i& \text{}& j\times i\hfill & \text{ for }j\times i[\lambda ^{}][\mu ^{}]\hfill \\ \hfill (g+1)\times (\mu _{g+1}^{}+i)& \text{}& g\times (\lambda _{g+1}^{}+i)\hfill & \text{ for }i[1,d]\hfill \end{array}$$
The semistandard correspondoid $`\{\phi \}`$ is in fact the only element in $`\{\mathrm{\Phi }^{\lambda ^{},\mu ^{}}\}_{\text{sst}}`$ (cf. after 3.6). With respect to $`[a_\lambda ^{}]=[\stackrel{ˇ}{a}_\lambda ^{}]`$ and $`[a_\mu ^{}]=[\stackrel{ˇ}{a}_\mu ^{}]`$, we obtain $`f^{}=\mathrm{\Theta }_{\phi ^1}^{}\nu _S^\mu `$, which reproves (4.3) by means of (3.1).
From (3.14, 3.13) we take the following assertions. Suppose given $`\stackrel{~}{m}1`$.
In case $`\mu ^{}`$ is $`2`$-regular, the map
$$\text{Hom}_{\text{Z}𝒮_n}(M^{\lambda ^{},}/\stackrel{~}{m},M^{\mu ^{},}/\stackrel{~}{m})\text{}\text{Hom}_{\text{Z}𝒮_n}(M^{\lambda ^{},}/\stackrel{~}{m},S^\mu /\stackrel{~}{m})$$
is surjective. In particular, the group $`\text{Hom}_{\text{Z}𝒮_n}(S^\lambda /\stackrel{~}{m},S^\mu /\stackrel{~}{m})`$ is cyclic since we may embed
$$\text{Hom}_{\text{Z}𝒮_n}(S^\lambda /\stackrel{~}{m},S^\mu /\stackrel{~}{m})\text{}\text{Hom}_{\text{Z}𝒮_n}(M^{\lambda ^{},}/\stackrel{~}{m},S^\mu /\stackrel{~}{m}).$$
A generator is a factorization of some multiple of $`f^{}_\text{Z}\text{Z}/\stackrel{~}{m}`$ over $`\nu _S^\lambda `$.
In case $`\mu ^{}`$ is $`2`$-singular, the map
$$2\text{Hom}_{\text{Z}𝒮_n}(M^{\lambda ^{},}/\stackrel{~}{m},M^{\mu ^{},}/\stackrel{~}{m})\text{}\mathrm{\hspace{0.33em}\hspace{0.33em}2}\text{Hom}_{\text{Z}𝒮_n}(M^{\lambda ^{},}/\stackrel{~}{m},S^\mu /\stackrel{~}{m})$$
is surjective. In particular, the group $`2\text{Hom}_{\text{Z}𝒮_n}(S^\lambda /\stackrel{~}{m},S^\mu /\stackrel{~}{m})`$ is cyclic since we may embed
$$2\text{Hom}_{\text{Z}𝒮_n}(S^\lambda /\stackrel{~}{m},S^\mu /\stackrel{~}{m})\text{}\mathrm{\hspace{0.33em}\hspace{0.33em}2}\text{Hom}_{\text{Z}𝒮_n}(M^{\lambda ^{},}/\stackrel{~}{m},S^\mu /\stackrel{~}{m}).$$
A generator is a factorization of some multiple of $`2f^{}_\text{Z}\text{Z}/\stackrel{~}{m}`$ over $`\nu _S^\lambda `$.
Let $`\xi a_g`$, $`s:=\mathrm{\#}\xi `$, $`\overline{\xi }:=a_g\backslash \xi `$, $`t[1,\lambda _{g+1}^{}]`$, $`\eta :=a_{g+1,[1,t]}`$, $`\overline{\eta }:=a_{g+1,[t+1,\lambda _{g+1}^{}]}`$ be such that $`s+t=\lambda _g^{}+1`$. Let $`u:=\mathrm{\#}\overline{\eta }=\lambda _{g+1}^{}t`$.
For $`x[0,\mathrm{min}(d,t)]`$, we denote
$$\begin{array}{ccc}\hfill Z^{}[x]& :=& \{\zeta ^{}Z^{}|\zeta ^{}(i)=i\text{ for }i[1,x],\zeta ^{}(i)t+1\text{ for }i[x+1,d]\}\hfill \\ \hfill Z[x]& :=& ZZ^{}[x]\hfill \\ \hfill \mathrm{\Lambda }([a],x)& :=& \frac{1}{s!t!}\underset{\zeta Z[x]}{}a^\zeta \epsilon _\zeta (\xi \eta )\hfill \\ & =& \frac{1}{s!t!(dx)!}\underset{\zeta ^{}Z^{}[x]}{}a^\zeta ^{}\epsilon _\zeta ^{}(\xi \eta )\hfill \end{array}$$
so that we can recover
$$G_{[a],\xi ,\eta }^{}f^{}=\underset{x[0,\mathrm{min}(d,t)]}{}\left(\begin{array}{c}t\\ x\end{array}\right)\mathrm{\Lambda }([a],x).$$
For $`0xy\mathrm{min}(d,t)`$, there is a rectification map
$$\begin{array}{ccc}\hfill Z^{}[x]& \text{}& Z^{}[y]\hfill \\ \hfill \zeta ^{}& \text{}& \zeta ^{}[y]\hfill \end{array}$$
defined by $`\zeta ^{}[y](i):=i`$ for $`i[1,y]`$ and $`\zeta ^{}[y](i):=\zeta ^{}(i)`$ for $`i[y+1,d]`$. In case $`td`$, we denote by $`\zeta [d]`$ the unique element in $`Z[d]=Z^{}[d]`$, mapping each $`i[1,d]`$ to $`i[1,\lambda _{g+1}^{}]`$.
###### Lemma 4.5
Suppose $`td`$. Let $`\overline{\xi }=\overline{\xi }_1\overline{\xi }_2`$ be a disjoint decomposition such that $`\mathrm{\#}\overline{\xi }_1=d1`$ and such that $`\mathrm{\#}\overline{\xi }_2=td`$. Let $`\eta _2:=a_{g+1,[d+1,t]}`$. We obtain, in the notation of (1.6),
$$G_{[a],\xi ,\eta }^{}f^{}=\epsilon _{\zeta [d]}\underset{i[0,d1]}{}(1)^i\left(\begin{array}{c}su+d\\ di\end{array}\right)B_{[a^{\zeta [d]}](\overline{\xi }_2,\eta _2)_\alpha ,\overline{\xi }_1,\overline{\xi }_2}(i),$$
independent of the chosen bijection $`\overline{\xi }_2\text{}\eta _2`$.
Given $`d_1[0,d]`$ and a map $`\zeta Z[dd_1]`$, we denote
$$\begin{array}{ccc}\hfill \eta _1& :=& a_{g+1,[dd_1+1,d]}\hfill \\ \hfill \overline{\eta }_1(\zeta )& :=& \zeta _{[a]}([dd_1+1,d]).\hfill \end{array}$$
We obtain
$$\begin{array}{ccc}\hfill \mathrm{\Lambda }([a],dd_1)& =& \frac{1}{s!t!}\underset{\zeta Z[dd_1]}{}a^\zeta \epsilon _\zeta (\xi \eta )\hfill \\ & \stackrel{\text{(}\text{1.4}\text{)}}{=}& \frac{(s+dd_1)!}{s!t!(d1)!}\underset{\zeta Z[dd_1]}{}a^\zeta \epsilon _\zeta (\overline{\xi }_2,\eta _2)(\overline{\eta }_1(\zeta ),\eta _1)(\overline{\xi }\overline{\eta }_1(\zeta ))\hfill \\ & =& (1)^{d_1}\frac{(s+dd_1)!}{s!t!(d1)!}\underset{\stackrel{\overline{\eta }_1\overline{\eta },}{\mathrm{\#}\overline{\eta }_1=d_1}}{}a^{\zeta [d]}\epsilon _{\zeta [d]}(\overline{\xi }_2,\eta _2)(\overline{\xi }\overline{\eta }_1)\hfill \\ & \stackrel{\text{(cf. }\text{1.6}\text{)}}{=}& (1)^{d_1}\frac{(s+dd_1)!}{s!t!(d1)!}\epsilon _{\zeta [d]}C_{[a^{\zeta [d]}](\overline{\xi }_2,\eta _2)_\alpha ,\overline{\xi }_1,\overline{\xi }_2}(d_1)\hfill \\ & \stackrel{\text{(}\text{1.6}\text{)}}{=}& (1)^{d_1}\frac{(s+dd_1)!(d1)!(td+d_1)!}{s!t!(d1)!}\epsilon _{\zeta [d]}\hfill \\ & & \underset{i[0,\mathrm{min}(d_1,d1)]}{}\left(\begin{array}{c}ui\\ d_1i\end{array}\right)B_{[a^{\zeta [d]}](\overline{\xi }_2,\eta _2)_\alpha ,\overline{\xi }_1,\overline{\xi }_2}(i),\hfill \end{array}$$
whence our one-step Garnir relation is mapped to
$$\begin{array}{cc}& G_{[a],\xi ,\eta }^{}f^{}\hfill \\ \hfill =& \epsilon _{\zeta [d]}\underset{d_1[0,d]}{}\underset{i[0,\mathrm{min}(d_1,d1)]}{}\left(\begin{array}{c}t\\ dd_1\end{array}\right)(1)^{d_1}\frac{(s+dd_1)!(td+d_1)!}{s!t!}\left(\begin{array}{c}ui\\ d_1i\end{array}\right)\hfill \\ & B_{[a^{\zeta [d]}](\overline{\xi }_2,\eta _2)_\alpha ,\overline{\xi }_1,\overline{\xi }_2}(i)\hfill \\ \hfill =& \epsilon _{\zeta [d]}\underset{i[0,d1]}{}\left(\underset{j[0,di]}{}(1)^{i+dij}\left(\begin{array}{c}s+j\\ j\end{array}\right)\left(\begin{array}{c}ui\\ dij\end{array}\right)\right)B_{[a^{\zeta [d]}](\overline{\xi }_2,\eta _2)_\alpha ,\overline{\xi }_1,\overline{\xi }_2}(i)\hfill \\ \hfill \stackrel{\text{(}\text{4.1}\text{)}}{=}& \epsilon _{\zeta [d]}\underset{i[0,d1]}{}(1)^i\left(\begin{array}{c}su+d\\ di\end{array}\right)B_{[a^{\zeta [d]}](\overline{\xi }_2,\eta _2)_\alpha ,\overline{\xi }_1,\overline{\xi }_2}(i).\hfill \end{array}$$
###### Lemma 4.6
Suppose $`2d1\lambda _{g+1}^{}`$. Let $`m_0:=\lambda _g^{}\lambda _{g+1}^{}+d+1`$ be the box shift length, let
$$m:=m_0\underset{p\text{ prime, }p|m_0}{}p^{\mathrm{min}(v_p(m_0),\text{ilog}_p(d))},$$
In case $`\mu ^{}`$ is $`2`$-regular, the injection
$$\text{Hom}_{\text{Z}𝒮_n}(S^\lambda ,S^\mu /m)\text{}\text{Hom}_{\text{Z}𝒮_n}(S^\lambda ,S^\mu /n!)$$
is surjective. In case $`\mu ^{}`$ is $`2`$-singular, the cokernel of the injection
$$\text{Hom}_{\text{Z}𝒮_n}(S^\lambda ,S^\mu /m)\text{}\text{Hom}_{\text{Z}𝒮_n}(S^\lambda ,S^\mu /n!)$$
is annihilated by multiplication by $`2`$.
Consider (4.5) in the case $`[a]=[\stackrel{ˇ}{a}_\lambda ]`$, $`t=d`$, $`\xi =a_{g,[1,s]}`$, $`\overline{\xi }_1=\overline{\xi }`$ and $`\overline{\xi }_2=\mathrm{}`$. We obtain
$$\begin{array}{ccc}\hfill G_{[\stackrel{ˇ}{a}_\lambda ],\xi ,\eta }^{}f^{}& =& \epsilon _{\zeta [d]}\underset{i[0,d1]}{}(1)^i\left(\begin{array}{c}su+d\\ di\end{array}\right)B_{[(\stackrel{ˇ}{a}_\lambda )^{\zeta [d]}],\overline{\xi },\mathrm{}}(i)\hfill \\ & \stackrel{\text{(}\text{1.6}\text{)}}{=}& \epsilon _{\zeta [d]}\underset{i[0,d1]}{}\left(\begin{array}{c}m_0\\ di\end{array}\right)\underset{\stackrel{\overline{\xi }_0\overline{\xi },}{\mathrm{\#}\overline{\xi }_0=i}}{}\underset{\stackrel{\phi _0\overline{\eta },}{\mathrm{\#}\phi _0=i}}{}(\stackrel{ˇ}{a}_\lambda )^{\zeta [d]}(\overline{\xi }_0,\phi _0).\hfill \end{array}$$
Since $`s=\lambda _g^{}+1d\mu _{g+1}^{}+1`$, the elements $`(\stackrel{ˇ}{a}_\lambda )^{\zeta [d]}(\overline{\xi }_0,\phi _0)`$ occurring in this expression are standard up to column permutation, i.e. up to sign. Moreover, they are pairwise different because of different column fillings. We note that each $`i[0,d1]`$ indexes a nonzero summand of this expresseion, namely a nonzero linear combination of different standard polytabloids equipped with coefficients $`\pm \left(\begin{array}{c}m_0\\ di\end{array}\right)`$, since we assumed $`2d1\lambda _{g+1}^{}`$, i.e. since $`d1\mathrm{\#}\overline{\eta }`$ leaves space for a subset $`\phi _0\overline{\eta }`$ of cardinality $`i`$. Therefore, by (4.2), the element $`G_{[\stackrel{ˇ}{a}_\lambda ],\xi ,\eta }^{}f^{}S^\mu `$ is exactly divisible by $`m`$.
Suppose given a $`\text{Z}𝒮_n`$-linear map $`S^\lambda \text{}S^\mu /n!`$. Let $`S^\mu \text{}S^\mu /n!`$ denote the residue class morphism.
Suppose $`\mu ^{}`$ to be $`2`$-regular. By (4.4), there is an integer $`z`$ such that $`\nu _S^\lambda \stackrel{~}{f}=zf^{}\phi `$. Thus $`zG_{[\stackrel{ˇ}{a}_\lambda ],\xi ,\eta }^{}f^{}`$ is divisible by $`n!`$, and so $`zm`$ is divisible by $`n!`$. Hence, each element in the image of $`\stackrel{~}{f}`$ is contained in $`(n!/m)S^\lambda /n!S^\lambda `$, i.e. there exists a factorization
$$(S^\lambda \text{}S^\mu /n!)=(S^\lambda \text{}S^\mu /m\text{}S^\mu /n!).$$
Suppose $`\mu ^{}`$ to be $`2`$-singular. By (4.4), there is an integer $`z`$ such that $`2\nu _S^\lambda \stackrel{~}{f}=2zf^{}\phi `$. Thus $`2zG_{[\stackrel{ˇ}{a}_\lambda ],\xi ,\eta }^{}f^{}`$ is divisible by $`n!`$, and so $`2zm`$ is divisible by $`n!`$. Hence, each element in the image of $`2\stackrel{~}{f}`$ is contained in $`(n!/m)S^\lambda /n!S^\lambda `$, i.e. there exists a factorization
$$(S^\lambda \text{}S^\mu /n!)=(S^\lambda \text{}S^\mu /m\text{}S^\mu /n!).$$
###### Lemma 4.7
Suppose $`td1`$. We obtain, in the notation of (1.7),
$$G_{[a],\xi ,\eta }^{}f^{}=\epsilon _{\zeta [d]}\underset{i[0,t1]}{}(1)^i\left(\begin{array}{c}su+d\\ ti\end{array}\right)B_{[a^{\zeta [d]}],\overline{\xi }\overline{\eta }_2,\overline{\eta }_2,\overline{\eta }\backslash \overline{\eta }_2}^{}(i).$$
Given $`d_0[0,t]`$ and an injection $`\zeta ^{}Z^{}[d_0]`$, we denote
$$\begin{array}{ccc}\hfill \eta _1& :=& \zeta _{[a]}^{}([1,d_0])=a_{g+1,[1,d_0]}\hfill \\ \hfill \overline{\eta }_1(\zeta ^{})& :=& \zeta _{[a]}^{}([d_0+1,t])\hfill \\ \hfill \overline{\eta }_2(\zeta ^{})& :=& \zeta _{[a]}^{}([t+1,d])\hfill \\ \hfill \overline{\eta }_2& :=& a_{g+1,[t+1,d]}.\hfill \end{array}$$
Given $`\zeta ^{}Z^{}[t]`$, we get
$$()a^\zeta ^{}\epsilon _\zeta ^{}=a^{\zeta ^{}[d]}\sigma \epsilon _\sigma \epsilon _{\zeta ^{}[d]}$$
for any permutation $`\sigma 𝒮_{\overline{\eta }}`$ that maps $`(a_{g+1,i})\sigma =\zeta _{[a]}^{}(i)`$ for $`i[t+1,d]`$. Note that $`\mathrm{\#}\{\sigma 𝒮_{\overline{\eta }}|(a_{g+1,i})\sigma =\zeta _{[a]}^{}(i)\text{ for }i[t+1,d]\}=(ud+t)!`$.
For $`d_0[0,t]`$, we obtain
$$\begin{array}{cc}& \mathrm{\Lambda }([a],d_0)\hfill \\ =& \frac{1}{s!t!(dd_0)!}\underset{\zeta ^{}Z^{}[d_0]}{}a^\zeta ^{}\epsilon _\zeta ^{}(\xi \eta )\hfill \\ \stackrel{\text{(}\text{1.4}\text{)}}{=}& \frac{(s+d_0)!}{s!t!(dd_0)!(d1)!}\underset{\zeta ^{}Z^{}[d_0]}{}a^\zeta ^{}\epsilon _\zeta ^{}(\eta \backslash \eta _1,\overline{\eta }_1(\zeta ^{}))(\overline{\xi }\overline{\eta }_1(\zeta ^{})\overline{\eta }_2(\zeta ^{}))\hfill \\ =& (1)^{td_0}\frac{(s+d_0)!(td_0)!}{s!t!(dd_0)!(d1)!}\underset{\zeta ^{}Z^{}[t]}{}\underset{\stackrel{\overline{\eta }_1\overline{\eta }\backslash \overline{\eta }_2\left(\zeta ^{}\right),}{\mathrm{\#}\overline{\eta }_1=td_0}}{}a^\zeta ^{}\epsilon _\zeta ^{}(\overline{\xi }\overline{\eta }_1\overline{\eta }_2(\zeta ^{}))\hfill \end{array}$$
$$\begin{array}{cc}\stackrel{()}{=}& (1)^{td_0}\frac{(s+d_0)!(td_0)!}{s!t!(dd_0)!(d1)!(ud+t)!}\underset{\stackrel{\overline{\eta }_1\overline{\eta }\backslash \overline{\eta }_2,}{\mathrm{\#}\overline{\eta }_1=td_0}}{}a^{\zeta [d]}\epsilon _{\zeta [d]}(\overline{\xi }\overline{\eta }_1\overline{\eta }_2)\overline{\eta }\hfill \\ \stackrel{\text{(cf. }\text{1.7}\text{)}}{=}& (1)^{td_0}\epsilon _{\zeta [d]}\frac{(s+d_0)!(td_0)!}{s!t!(dd_0)!(d1)!(ud+t)!}C_{[a^{\zeta [d]}],\overline{\xi }\overline{\eta }_2,\overline{\eta }_2,\overline{\eta }\backslash \overline{\eta }_2}^{}(td_0)\hfill \\ \stackrel{\text{(}\text{1.7}\text{)}}{=}& (1)^{td_0}\epsilon _{\zeta [d]}\frac{(s+d_0)!(td_0)!}{s!t!(dd_0)!(d1)!(ud+t)!}\hfill \\ & \left(ud+t\right)!\left(d1\right)!\left(dd_0\right)!\underset{i[0,\mathrm{min}(td_0,t1)]}{}\left(\begin{array}{c}ud+ti\\ td_0i\end{array}\right)B_{[a^{\zeta [d]}],\overline{\xi }\overline{\eta }_2,\overline{\eta }_2,\overline{\eta }\backslash \overline{\eta }_2}^{}(i),\hfill \end{array}$$
whence
$$\begin{array}{cc}& G_{[a],\xi ,\eta }^{}f^{}\hfill \\ \hfill =& \epsilon _{\zeta [d]}\underset{d_0[0,t]}{}\underset{i[0,\mathrm{min}(td_0,t1)]}{}\left(\begin{array}{c}t\\ d_0\end{array}\right)(1)^{td_0}\frac{(s+d_0)!(td_0)!}{s!t!}\hfill \\ & \left(\begin{array}{c}ud+ti\\ td_0i\end{array}\right)B_{[a^{\zeta [d]}],\overline{\xi }\overline{\eta }_2,\overline{\eta }_2,\overline{\eta }\backslash \overline{\eta }_2}^{}(i)\hfill \\ \hfill =& \epsilon _{\zeta [d]}\underset{i[0,t1]}{}\left(\underset{d_0[0,ti]}{}(1)^{i+tid_0}\left(\begin{array}{c}s+d_0\\ d_0\end{array}\right)\left(\begin{array}{c}ud+ti\\ tid_0\end{array}\right)\right)B_{[a^{\zeta [d]}],\overline{\xi }\overline{\eta }_2,\overline{\eta }_2,\overline{\eta }\backslash \overline{\eta }_2}^{}(i)\hfill \\ \hfill \stackrel{\text{(}\text{4.1}\text{)}}{=}& \epsilon _{\zeta [d]}\underset{i[0,t1]}{}(1)^i\left(\begin{array}{c}su+d\\ ti\end{array}\right)B_{[a^{\zeta [d]}],\overline{\xi }\overline{\eta }_2,\overline{\eta }_2,\overline{\eta }\backslash \overline{\eta }_2}^{}(i).\hfill \end{array}$$
###### Remark 4.8
The morphism $`F^\lambda \text{}S^\mu `$ maps $`[\stackrel{ˇ}{a}_\lambda ]`$ to a linear combination of standard $`\mu `$-polytabloids with coefficients $`\pm 1`$.
We summarize.
###### Theorem 4.9
Let $`m_0:=\lambda _g^{}\lambda _{g+1}^{}+d+1`$ be the box shift length, let
$$m:=m_0\underset{p\text{ prime, }p|m_0}{}p^{\mathrm{min}(v_p(m_0),\text{ilog}_p(d))},$$
where $`\text{ilog}_p(k)=\mathrm{max}\{i0|p^ik\}`$. The $`\text{Z}𝒮_n`$-linear map $`M^{\lambda ^{},}\text{}S^\mu `$ factors over
The resulting morphism $`S^\lambda \text{}S^\mu /m`$ is of order $`m`$ in $`\text{Hom}_{\text{Z}𝒮_n}(S^\lambda ,S^\mu /m)`$.
In case $`\mu ^{}`$ is $`2`$-regular, the group $`\text{Hom}_{\text{Z}𝒮_n}(S^\lambda ,S^\mu /m)`$ is generated by $`f`$.
In case $`\mu ^{}`$ is $`2`$-singular, the group $`2\text{Hom}_{\text{Z}𝒮_n}(S^\lambda ,S^\mu /m)`$ is generated by $`2f`$.
In case $`2d1\lambda _{g+1}^{}`$ and $`\mu ^{}`$ is $`2`$-regular, there is an isomorphism
$$\text{Hom}_{\text{Z}𝒮_n}(S^\lambda ,S^\mu /m)\text{}\text{Hom}_{\text{Z}𝒮_n}(S^\lambda ,S^\mu /n!).$$
In case $`2d1\lambda _{g+1}^{}`$ and $`\mu ^{}`$ is $`2`$-singular, the cokernel of the inclusion
$$\text{Hom}_{\text{Z}𝒮_n}(S^\lambda ,S^\mu /m)\text{}\text{Hom}_{\text{Z}𝒮_n}(S^\lambda ,S^\mu /n!),$$
is annihilated by multiplication with $`2`$.
This follows by (4.5, 4.7, 4.2, 4.8, 4.4, 4.6).
###### Example 4.10
We have $`\text{Hom}_{\text{Z}𝒮_8}(S^{(2^3,1^2)},S^{(1^8)}/8!)=0`$.
###### Question 4.11
In case $`2d1>\lambda _{g+1}^{}`$, I do not know of a counterexample to the assertions in (4.9) that presuppose $`2d1\lambda _{g+1}^{}`$.
###### Remark 4.12
In the situation of (4.9), Carter and Payne \[CP 80\] have shown that
$$\text{Hom}_{K𝒮_n}(K_\text{Z}S^\lambda ,K_\text{Z}S^\mu )0,$$
$`K`$ being an infinite field of characteristic $`p`$ such that $`\text{ilog}_pd<v_p(m_0)`$. This part of their result is recovered by (4.9).
###### Corollary 4.13 (to (4.9) via (3.16, 4.4))
The transpose in the sense of (3.16) of the morphism given in (4.9), interpreted as $`S^\lambda /m\text{}S^\mu /m`$, is given by
$$\begin{array}{ccc}\hfill S^\mu ^{}/m& \text{}& S^\lambda ^{}/m\hfill \\ \hfill b^{}& \text{}& b^{}\mathrm{\Theta }_\phi ,\hfill \end{array}$$
where the correspondoid $`\{\phi \}\{\mathrm{\Phi }^{\lambda ^{},\mu ^{}}\}`$ is given by
$$\begin{array}{cccc}\hfill [\mu ^{}]& \text{}& \text{N}\hfill & \\ \hfill j\times i& \text{}& j\hfill & \text{ for }j\times i[\lambda ^{}][\mu ^{}]\hfill \\ \hfill g\times (\lambda _g^{}+i)& \text{}& g+1\hfill & \text{ for }i[1,d]\hfill \end{array}$$
All assertions on the Hom-groups in (4.9) have a counterpart, obtained by isomophic transport via the transposition isomorphism (3.16).
This follows by (4.9, 3.16) using (4.4). Concerning $`\mathrm{\Theta }_\phi `$, we refer to (3.1).
###### Question 4.14
I do not know a formula for the transpose $`f^\text{t}`$ in (4.13) in terms of $`\lambda ^{}`$-polytabloids. Cf. (3.17, 4.15).
###### Example 4.15 (a fixed point, cf. \[J 78, 24.4\], \[K 99, 4.4.1\])
Consider the case $`n4`$, $`d=2`$, $`\lambda =(2^2,1^{n4})`$ and $`\mu =(1^n)`$. The modulus becomes $`m=n1`$ if $`n1`$ is odd and $`m=(n1)/2`$ if $`n1`$ is even. We obtain
$$\begin{array}{ccc}\hfill [\mu ^{}]& \text{}& \text{N}\hfill \\ \hfill 1\times i& \text{}& \{\begin{array}{cc}1\hfill & \text{for }i[1,n2]\hfill \\ 2\hfill & \text{for }i[n1,n],\hfill \end{array}\hfill \end{array}$$
which allows to calculate the transpose as mapping
$$\begin{array}{ccc}\hfill S^\mu ^{}/m& \text{}& S^\lambda ^{}/m\hfill \\ \hfill x& \text{}& \underset{i,j[1,n],i<j}{}\{i,j\},\hfill \end{array}$$
where $`[x]`$ is the unique standard $`\mu ^{}`$-tableau and where we denote a $`\lambda ^{}`$-tabloid by its second row. Similarly, for a $`\lambda ^{}`$-polytabloid, we omit to denote the entries in the first row from column three onwards. I.e. we let a $`\lambda ^{}`$-polytabloid be determined by the first two entries of the first and of the second row - which is not to be confused with a $`(2,2)`$-polytabloid. Let
$$A:=\underset{j[4,n]}{}(j2)\begin{array}{cc}1& 3\\ 2& j\end{array}\underset{i,j[3,n],i<j}{}\begin{array}{cc}1& 2\\ i& j\end{array}S^\lambda ^{}.$$
For $`k[3,n]`$, $`l[4,n]`$, $`k<l`$, we obtain the following list of values of the scalar product in $`M^\lambda ^{}`$.
$$\begin{array}{cccccc}(A,\{k,l\})\hfill & =& +& 0\hfill & & 1\hfill \\ (A,\{1,l\})\hfill & =& & (l2)\hfill & +& (l3)\hfill \\ (A,\{2,l\})\hfill & =& +& (l2)\hfill & +& (nl)\hfill \\ (A,\{2,3\})\hfill & =& & (n2)(n1)/2+1\hfill & +& (n3)\hfill \\ (A,\{1,2\})\hfill & =& +& 0\hfill & & (n2)(n3)/2\hfill \\ (A,\{1,3\})\hfill & =& +& (n2)(n1)/21\hfill & +& 0.\hfill \end{array}$$
Hence, written in polytabloids, our map turns out to be
$$\begin{array}{ccc}\hfill S^\mu ^{}/m& \text{}& S^\lambda ^{}/m\hfill \\ \hfill x& \text{}& \underset{j[4,n]}{}(j2)\begin{array}{cc}1& 3\\ 2& j\end{array}\underset{i,j[3,n],i<j}{}\begin{array}{cc}1& 2\\ i& j\end{array}.\hfill \end{array}$$
## Appendix A Two lemmata
We present full versions of the two main lemmata of Section 2.2. For sake of completeness of the full versions, there is a certain overlap, concerning notation and conclusion, with their abridged versions (2.4, 2.10). We maintain the notation of Section 2.2.
###### Lemma A.1 (full version of (2.4))
Suppose $`g<p<k`$ and $`s,t2`$. The map $`f^{}`$ annihilates $`G_{[a],\xi ,\eta }^{}`$.
We fix a map
$$\begin{array}{ccc}\hfill [1,\lambda _1]\backslash \{p,p+1\}& \text{}& [0,2]\hfill \\ \hfill j& \text{}& \stackrel{~}{e}_j\hfill \end{array}$$
that maps $`g`$ and $`k+1`$ to $`e_g=e_{k+1}=2`$, and that maps $`j[1,\lambda _1]\backslash [g,k+1]`$ to $`e_j=0`$. For $`\alpha ,\beta [0,2]`$, we denote by $`\stackrel{~}{e}\alpha \beta `$ be the prolongation of $`\stackrel{~}{e}`$ to $`[1,\lambda _1]`$ defined by $`\stackrel{~}{e}\alpha \beta |_{[1,\lambda _1]\backslash \{p,p+1\}}:=\stackrel{~}{e}`$, $`\left(\stackrel{~}{e}\alpha \beta \right)_p:=\alpha `$ and $`\left(\stackrel{~}{e}\alpha \beta \right)_{p+1}:=\beta `$. We contend that
$$\underset{\alpha ,\beta [0,2]}{}X_p^{\left(2\alpha \right)}X_{p+1}^{\left(2\beta \right)}G_{\left[a\right],\xi ,\eta }^{}f_{\stackrel{~}{e}\alpha \beta }^{}=0,$$
from which the lemma ensues.
Given $`\alpha ,\beta [0,2]`$, given a subset $`\stackrel{~}{\mathrm{\Xi }}[1,2]\times \{p,p+1\}`$ with
$$\begin{array}{ccc}\hfill \mathrm{\#}\left(\stackrel{~}{\mathrm{\Xi }}\left([1,2]\times \left\{p\right\}\right)\right)& =& \alpha \hfill \\ \hfill \mathrm{\#}\left(\stackrel{~}{\mathrm{\Xi }}\left([1,2]\times \left\{p+1\right\}\right)\right)& =& \beta ,\hfill \end{array}$$
and given prescribed inverse images $`\xi ^{},\eta ^{}[1,2]`$ ‘within $`\stackrel{~}{\mathrm{\Xi }}`$’, i.e. such that
$$()\begin{array}{ccc}\hfill \xi ^{}\times \left\{p\right\}& & \stackrel{~}{\mathrm{\Xi }}\left([1,2]\times \left\{p\right\}\right)\hfill \\ \hfill \eta ^{}\times \left\{p+1\right\}& & \stackrel{~}{\mathrm{\Xi }}\left([1,2]\times \left\{p+1\right\}\right),\hfill \end{array}$$
we let
$$\begin{array}{ccccc}\hfill \dot{\mathrm{\Gamma }}(\stackrel{~}{e},\stackrel{~}{\mathrm{\Xi }},\xi ^{},\eta ^{})& :=& \hfill \{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}\alpha \beta \right)& |& \mathrm{\Xi }_\gamma \left([1,2]\times \{p,p+1\}\right)=\stackrel{~}{\mathrm{\Xi }},\hfill \\ & & & & \gamma ^1\left(\xi \right)=\xi ^{}\times \left\{p\right\},\gamma ^1\left(\eta \right)=\eta ^{}\times \{p+1\}\}\hfill \end{array}$$
and form the partial sum
$$\mathrm{\Lambda }(\left[a\right],\stackrel{~}{e},\stackrel{~}{\mathrm{\Xi }},\xi ^{},\eta ^{}):=\frac{1}{s!t!}\underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e},\stackrel{~}{\mathrm{\Xi }},\xi ^{},\eta ^{})}{}a^\gamma \epsilon _\gamma \left(\xi \eta \right)$$
so that we can recover
$$()G_{\left[a\right],\xi ,\eta }^{}f_{\stackrel{~}{e}\alpha \beta }^{}=\underset{\stackrel{~}{\mathrm{\Xi }}[1,2]\times \{p,p+1\}}{}\underset{\xi ^{},\eta ^{}[1,2]\text{ subject to }()}{}\mathrm{\Lambda }(\left[a\right],\stackrel{~}{e},\stackrel{~}{\mathrm{\Xi }},\xi ^{},\eta ^{}).$$
There exist elements $`x,y\xi `$, $`xy`$, $`x^{},y^{}\eta `$, $`x^{}y^{}`$, $`z\overline{\xi }`$, which we choose and fix. For $`\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)`$, we let $`x_\gamma :=a_{j,\overline{\gamma }(j,1)}`$, where $`j[p+2,k+1]`$ is minimal with $`1\times j\mathrm{\Xi }_\gamma `$, and $`y_\gamma :=a_{j,\overline{\gamma }(j,2)}`$, where $`j[p+2,k+1]`$ is minimal with $`2\times j\mathrm{\Xi }_\gamma `$. I.e. we pick the entries $`x_\gamma ,y_\gamma `$ that ‘cross the columns’ $`p`$ and $`p+1`$ under the operation of $`\gamma `$. We write
$$\begin{array}{ccc}\hfill U_\gamma & :=& \left(s+t2\right)!^1a^\gamma \epsilon _\gamma (\overline{\xi }\backslash z,\eta \backslash \{x^{},y^{}\})(x_\gamma ,x,x^{})(y_\gamma ,y,y^{})\left(\xi \eta \right)\hfill \\ \hfill V_{1,\gamma }& :=& \left(s+t1\right)!^1a^\gamma \epsilon _\gamma (\overline{\xi },\eta \backslash x^{})(x_\gamma ,x,x^{})\left(\xi \eta \right)\hfill \\ \hfill V_{2,\gamma }& :=& \left(s+t1\right)!^1a^\gamma \epsilon _\gamma (\overline{\xi },\eta \backslash y^{})(y_\gamma ,y,y^{})\left(\xi \eta \right),\hfill \end{array}$$
and let
$$\begin{array}{ccc}A\hfill & :=& \underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}U_\gamma \underset{w^{}\overline{\eta }}{}(w^{},z)\hfill \\ B\hfill & :=& \underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}U_\gamma \left(1\underset{w\overline{\xi }\backslash z}{}(w,z)\right)\hfill \\ C_1\hfill & :=& \underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}U_\gamma (z,x_\gamma )\hfill \\ C_2\hfill & :=& \underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}U_\gamma (z,y_\gamma )\hfill \\ D\hfill & :=& \underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}V_{1,\gamma }\underset{w\overline{\xi }}{}(w,y_\gamma )\hfill \\ & =& \underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}V_{2,\gamma }\underset{w\overline{\xi }}{}(w,x_\gamma )\hfill \\ H\hfill & :=& \underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}V_{1,\gamma }\underset{w^{}\overline{\eta }}{}(w^{},y_\gamma )\hfill \\ & =& \underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}V_{2,\gamma }\underset{w^{}\overline{\eta }}{}(w^{},x_\gamma )\hfill \\ F_1\hfill & :=& \underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}V_{1,\gamma }\hfill \\ F_2\hfill & :=& \underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}V_{2,\gamma }.\hfill \end{array}$$
The equalities herein follow by the argument of (2.2 i) and by the independence of of the respective expressions of the choice of $`x,y,x^{},y^{}`$.
Calculation of $`\mathrm{\Lambda }`$-values. We shall distinguish subcases and subsubcases according to the summation in $`()`$.
To distinguish subcases, e.g. $`\left[\begin{array}{cc}+& \\ +& +\end{array}\right]`$ designates the subset $`\stackrel{~}{\mathrm{\Xi }}=\{1\times p,\mathrm{\hspace{0.33em}2}\times p,\mathrm{\hspace{0.33em}2}\times \left(p+1\right)\}`$; in general, the first factor of an element of $`[1,2]\times \{p,p+1\}`$ counts rows, the second counts columns, and a plus sign $`+`$ denotes its appearance in $`\stackrel{~}{\mathrm{\Xi }}`$.
To distinguish subsubcases, we denote e.g. by $`\left[\begin{array}{cc}& \\ +& \end{array}\right]`$ the configuration $`\stackrel{~}{\mathrm{\Xi }}=\left[\begin{array}{cc}+& \\ +& +\end{array}\right]`$, $`\xi ^{}=\left\{1\right\}`$ and $`\eta ^{}=\left\{2\right\}`$; in general, the sign $``$ in the left column means that its row number is an element of $`\xi ^{}`$, the sign $``$ in the right column means that its row number is an element of $`\eta ^{}`$. Concerning the underlying set $`\stackrel{~}{\mathrm{\Xi }}`$, the symbols $``$ and $`+`$ are synonymous.
Case $`f_{\stackrel{~}{e}22}^{}`$.
Subcase $`\stackrel{~}{\mathrm{\Xi }}=\left[\begin{array}{cc}+& +\\ +& +\end{array}\right]`$.
Subsubcase $`(\stackrel{~}{\mathrm{\Xi }},\xi ^{},\eta ^{})=\left[\begin{array}{cc}& \\ & \end{array}\right]`$. We calculate
$$\begin{array}{cc}& \mathrm{\Lambda }(\left[a\right],\stackrel{~}{e},\text{},\{1,2\},\{1,2\})\hfill \\ =& \frac{1}{s!t!}\underset{\gamma \dot{\mathrm{\Gamma }}(e,\stackrel{~}{\mathrm{\Xi }},\{1,2\},\{1,2\})}{}a^\gamma \epsilon _\gamma \left(\xi \eta \right)\hfill \\ =& \frac{1}{s!t!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{x_1y_1\xi ,x_1^{}y_1^{}\eta }{}a^\gamma \epsilon _\gamma (x_\gamma ,x_1,x_1^{})(y_\gamma ,y_1,y_1^{})\left(\xi \eta \right)\hfill \\ =& \frac{s\left(s1\right)t\left(t1\right)}{s!t!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}a^\gamma \epsilon _\gamma (x_\gamma ,x,x^{})(y_\gamma ,y,y^{})\left(\xi \eta \right)\hfill \\ \stackrel{\text{(}\text{1.5}\text{)}}{=}& \frac{s\left(s1\right)t\left(t1\right)s!}{s!t!\left(s+t2\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}a^\gamma \epsilon _\gamma (\overline{\xi }\backslash z,\eta \backslash \{x^{},y^{}\})(x_\gamma ,x,x^{})(y_\gamma ,y,y^{})\left(\xi \eta \right)\overline{\xi }\hfill \\ =& s\left(s1\right)B.\hfill \end{array}$$
Subsubcase $`(\stackrel{~}{\mathrm{\Xi }},\xi ^{},\eta ^{})=\left[\begin{array}{cc}& \\ & +\end{array}\right]`$. We calculate
$$\begin{array}{cc}& \mathrm{\Lambda }(\left[a\right],\stackrel{~}{e},\text{},\{1,2\},\left\{1\right\})\hfill \\ =& \frac{s\left(s1\right)t}{s!t!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{y_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x,x^{})(y_\gamma ,y,y_1^{})\left(\xi \eta \right)\hfill \\ \stackrel{\text{(}\text{1.5}\text{)}}{=}& \frac{s\left(s1\right)t\left(s1\right)!}{s!t!\left(s+t2\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{y_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma \hfill \\ & (x_\gamma ,x,x^{})(y_\gamma ,y,y_1^{})(\overline{\xi }\backslash z,\eta \backslash \{x^{},y^{}\})(y_1^{},y^{})\left(\xi \eta \right)\left(\overline{\xi }y_1^{}\right)\hfill \\ =& \frac{s1}{\left(t1\right)!\left(s+t2\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{y_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma \hfill \\ & (x_\gamma ,x,x^{})(y_\gamma ,y,y_1^{})(\overline{\xi }\backslash z,\eta \backslash \{x^{},y^{}\})(y_1^{},y^{})(y_1^{},y_\gamma )\left(\xi \eta \right)\left(\overline{\xi }y_1^{}\right)\hfill \\ =& \frac{s1}{\left(t1\right)!\left(s+t2\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{y_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x,x^{})(y_\gamma ,y,y^{})(\overline{\xi }\backslash z,\eta \backslash \{x^{},y^{}\})\left(\xi \eta \right)\left(\overline{\xi }y_1^{}\right)\hfill \\ =& \left(s1\right)\left(uBA\right).\hfill \end{array}$$
Subsubcase $`(\stackrel{~}{\mathrm{\Xi }},\xi ^{},\eta ^{})=\left[\begin{array}{cc}& +\\ & \end{array}\right]`$. By symmetry, we obtain from subsubcase $`\left[\begin{array}{cc}& \\ & +\end{array}\right]`$
$$\mathrm{\Lambda }(\left[a\right],\stackrel{~}{e},\text{},\{1,2\},\left\{2\right\})=\left(s1\right)\left(uBA\right).$$
Subsubcase $`(\stackrel{~}{\mathrm{\Xi }},\xi ^{},\eta ^{})=\left[\begin{array}{cc}& +\\ & +\end{array}\right]`$. We calculate
$$\begin{array}{cc}& \mathrm{\Lambda }(\left[a\right],\stackrel{~}{e},\text{},\{1,2\},\left\{\right\})\hfill \\ =& \frac{s\left(s1\right)}{s!t!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{x_1^{}y_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x,x_1^{})(y_\gamma ,y,y_1^{})\left(\xi \eta \right)\hfill \\ \stackrel{\text{(}\text{1.5}\text{)}}{=}& \frac{s\left(s1\right)\left(s2\right)!}{s!t!\left(s+t2\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{x_1^{}y_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma \hfill \\ & (x_\gamma ,x,x_1^{})(y_\gamma ,y,y_1^{})(\overline{\xi }\backslash z,\eta \backslash \{x^{},y^{}\})(x_1^{},x^{})(y_1^{},y^{})\left(\xi \eta \right)\left(\overline{\xi }x_1^{}y_1^{}\right)\hfill \\ =& \frac{1}{t!\left(s+t2\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{x_1^{}y_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma \hfill \\ & (x_\gamma ,x,x_1^{})(y_\gamma ,y,y_1^{})(\overline{\xi }\backslash z,\eta \backslash \{x^{},y^{}\})(x_1^{},x^{})(y_1^{},y^{})(x_1^{},x_\gamma )(y_1^{},y_\gamma )\hfill \\ & \left(\xi \eta \right)\left(\overline{\xi }x_1^{}y_1^{}\right)\hfill \\ =& \frac{1}{t!\left(s+t2\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{x_1^{}y_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x,x^{})(y_\gamma ,y,y^{})(\overline{\xi }\backslash z,\eta \backslash \{x^{},y^{}\})\left(\xi \eta \right)\left(\overline{\xi }x_1^{}y_1^{}\right)\hfill \\ =& \left(u1\right)\left(uB2A\right).\hfill \end{array}$$
Subsubcase $`(\stackrel{~}{\mathrm{\Xi }},\xi ^{},\eta ^{})=\left[\begin{array}{cc}& \\ +& \end{array}\right]`$. We calculate
$$\begin{array}{cc}& \mathrm{\Lambda }(\left[a\right],\stackrel{~}{e},\text{},\left\{1\right\},\{1,2\})\hfill \\ =& \frac{st\left(t1\right)}{s!t!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{y_1\overline{\xi }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x,x^{})(y_\gamma ,y_1,y^{})\left(\xi \eta \right)\hfill \\ \stackrel{\text{(}\text{1.5}\text{)}}{=}& \frac{st\left(t1\right)\left(s+1\right)!\left(t2\right)!}{s!t!\left(s+t1\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{y_1\overline{\xi }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x,x^{})(y_\gamma ,y_1,y^{})(\overline{\xi }\backslash y_1,\eta \backslash \{x^{},y^{}\})\left(\xi \eta \right)\hfill \\ =& \frac{s\left(s+1\right)}{\left(s+t1\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{y_1\overline{\xi }}{}a^\gamma \epsilon _\gamma (\overline{\xi },\eta \backslash x^{})(x_\gamma ,x,x^{})(y_\gamma ,y_1)\left(\xi \eta \right)\hfill \\ =& \left(s+1\right)sD.\hfill \end{array}$$
Subsubcase $`(\stackrel{~}{\mathrm{\Xi }},\xi ^{},\eta ^{})=\left[\begin{array}{cc}& \\ +& +\end{array}\right]`$. We calculate
$$\begin{array}{cc}& \mathrm{\Lambda }(\left[a\right],\stackrel{~}{e},\text{},\left\{1\right\},\left\{1\right\})\hfill \\ =& \frac{st}{s!t!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{y_1\overline{\xi },y_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x,x^{})(y_\gamma ,y_1,y_1^{})\left(\xi \eta \right)\hfill \\ \stackrel{\text{(}\text{1.5}\text{)}}{=}& \frac{sts!\left(t1\right)!}{s!t!\left(s+t1\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{y_1\overline{\xi },y_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x,x^{})(y_\gamma ,y_1,y_1^{})(\overline{\xi }\backslash y_1,\eta \backslash \{x^{},y^{}\})(y_1^{},y^{})\left(\xi \eta \right)\hfill \\ =& \frac{s}{\left(s+t1\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{y_1\overline{\xi },y_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (\overline{\xi },\eta \backslash x^{})(x_\gamma ,x,x^{})(y_\gamma ,y_1,y^{},y_1^{})(y_1^{},y^{})(y_\gamma ,y_1^{})\left(\xi \eta \right)\hfill \\ =& \frac{su}{\left(s+t1\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{y_1\overline{\xi }}{}a^\gamma \epsilon _\gamma (\overline{\xi },\eta \backslash x^{})(x_\gamma ,x,x^{})(y_\gamma ,y_1)\left(\xi \eta \right)\hfill \\ =& suD.\hfill \end{array}$$
Subsubcase $`(\stackrel{~}{\mathrm{\Xi }},\xi ^{},\eta ^{})=\left[\begin{array}{cc}& +\\ +& \end{array}\right]`$. We calculate
$$\begin{array}{cc}& \mathrm{\Lambda }(\left[a\right],\stackrel{~}{e},\text{},\left\{1\right\},\left\{2\right\})\hfill \\ =& \frac{st}{s!t!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{y_1\overline{\xi },x_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x,x_1^{})(y_\gamma ,y_1,y^{})\left(\xi \eta \right)\hfill \\ \stackrel{\text{(}\text{1.5}\text{)}}{=}& \frac{sts!\left(t1\right)!}{s!t!\left(s+11\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{y_1\overline{\xi },x_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x,x_1^{})(y_\gamma ,y_1,y^{})(\overline{\xi }\backslash y_1,\eta \backslash \{x^{},y^{}\})(x_1^{},x^{})\left(\xi \eta \right)\hfill \\ =& \frac{s}{\left(s+11\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{y_1\overline{\xi },x_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (\overline{\xi },\eta \backslash x^{})(x_\gamma ,x,x_1^{})(y_\gamma ,y_1)(x_1^{},x^{})(x_\gamma ,x_1^{})\left(\xi \eta \right)\hfill \\ =& \frac{su}{\left(s+11\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{y_1\overline{\xi }}{}a^\gamma \epsilon _\gamma (\overline{\xi },\eta \backslash x^{})(x_\gamma ,x,x^{})(y_\gamma ,y_1)\left(\xi \eta \right)\hfill \\ =& suD.\hfill \end{array}$$
Subsubcase $`(\stackrel{~}{\mathrm{\Xi }},\xi ^{},\eta ^{})=\left[\begin{array}{cc}& +\\ +& +\end{array}\right]`$. We calculate
$$\begin{array}{cc}& \mathrm{\Lambda }(\left[a\right],\stackrel{~}{e},\text{},\left\{1\right\},\left\{\right\})\hfill \\ =& \frac{s}{s!t!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{y_1\overline{\xi },x_1^{}y_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x,x_1^{})(y_\gamma ,y_1,y_1^{})\left(\xi \eta \right)\hfill \\ \stackrel{\text{(}\text{1.5}\text{)}}{=}& \frac{s\left(s1\right)!t!}{s!t!\left(s+t1\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{y_1\overline{\xi },x_1^{}y_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma \hfill \\ & (x_\gamma ,x,x_1^{})(y_\gamma ,y_1,y_1^{})(\overline{\xi }\backslash y_1,\eta \backslash \{x^{},y^{}\})(x_1^{},x^{})(y_1^{},y^{})\left(\xi \eta \right)\hfill \\ =& \frac{1}{\left(s+t1\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{y_1\overline{\xi },x_1^{}y_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma \hfill \\ & (\overline{\xi },\eta \backslash x^{})(x_\gamma ,x,x_1^{})(y_\gamma ,y_1,y^{},y_1^{})(x_1^{},x^{})(y_1^{},y^{})(x_\gamma ,x_1^{})(y_\gamma ,y_1^{})\left(\xi \eta \right)\hfill \\ =& \frac{u\left(u1\right)}{\left(s+t1\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{y_1\overline{\xi }}{}a^\gamma \epsilon _\gamma (\overline{\xi },\eta \backslash x^{})(x_\gamma ,x,x^{})(y_\gamma ,y_1)\left(\xi \eta \right)\hfill \\ =& u\left(u1\right)D.\hfill \end{array}$$
Subsubcase $`(\stackrel{~}{\mathrm{\Xi }},\xi ^{},\eta ^{})=\left[\begin{array}{cc}+& \\ & \end{array}\right]`$. By symmetry, we obtain from subsubcase $`\left[\begin{array}{cc}& \\ +& \end{array}\right]`$
$$\mathrm{\Lambda }(\left[a\right],\stackrel{~}{e},\text{},\left\{2\right\},\{1,2\})=\left(s+1\right)sD.$$
Subsubcase $`(\stackrel{~}{\mathrm{\Xi }},\xi ^{},\eta ^{})=\left[\begin{array}{cc}+& \\ & +\end{array}\right]`$. By symmetry, we obtain from subsubcase $`\left[\begin{array}{cc}& +\\ +& \end{array}\right]`$
$$\mathrm{\Lambda }(\left[a\right],\stackrel{~}{e},\text{},\left\{2\right\},\left\{1\right\})=suD.$$
Subsubcase $`(\stackrel{~}{\mathrm{\Xi }},\xi ^{},\eta ^{})=\left[\begin{array}{cc}+& +\\ & \end{array}\right]`$. By symmetry, we obtain from subsubcase $`\left[\begin{array}{cc}& \\ +& +\end{array}\right]`$
$$\mathrm{\Lambda }(\left[a\right],\stackrel{~}{e},\text{},\left\{2\right\},\left\{2\right\})=suD.$$
Subsubcase $`(\stackrel{~}{\mathrm{\Xi }},\xi ^{},\eta ^{})=\left[\begin{array}{cc}+& +\\ & +\end{array}\right]`$. By symmetry, we obtain from subsubcase $`\left[\begin{array}{cc}& +\\ +& +\end{array}\right]`$
$$\mathrm{\Lambda }(\left[a\right],\stackrel{~}{e},\text{},\left\{2\right\},\left\{\right\})=u\left(u1\right)D.$$
Subsubcases $`\left[\begin{array}{cc}+& \\ +& \end{array}\right]`$, $`\left[\begin{array}{cc}+& \\ +& +\end{array}\right]`$, $`\left[\begin{array}{cc}+& +\\ +& \end{array}\right]`$, $`\left[\begin{array}{cc}+& +\\ +& +\end{array}\right]`$. The $`\mathrm{\Lambda }`$-value is zero by the Garnir relations.
Summing up $`\mathrm{\Lambda }`$-values over subcases and subsubcases, we obtain
$$G_{\left[a\right],\xi ,\eta }^{}f_{\stackrel{~}{e}22}^{}=2\left(su\right)A+\left(su\right)\left(su1\right)B+2\left(su\right)\left(su+1\right)D$$
Case $`f_{\stackrel{~}{e}12}^{}`$.
Subcase $`\stackrel{~}{\mathrm{\Xi }}=\left[\begin{array}{cc}+& +\\ & +\end{array}\right]`$.
Subsubcase $`(\stackrel{~}{\mathrm{\Xi }},\xi ^{},\eta ^{})=\left[\begin{array}{cc}& \\ & \end{array}\right]`$. We calculate
$$\begin{array}{cc}& \mathrm{\Lambda }(\left[a\right],\stackrel{~}{e},\text{},\left\{1\right\},\{1,2\})\hfill \\ =& \frac{st\left(t1\right)}{s!t!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}a^\gamma \epsilon _\gamma (x_\gamma ,x,x^{})(y_\gamma ,y^{})\left(\xi \eta \right)\hfill \\ \stackrel{\text{(}\text{1.5}\text{)}}{=}& \frac{st\left(t1\right)s!}{s!t!\left(s+t2\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}a^\gamma \epsilon _\gamma (x_\gamma ,x,x^{})(y_\gamma ,y^{})(\overline{\xi }\backslash z,\eta \backslash \{x^{},y^{}\})\left(\xi \eta \right)\overline{\xi }\hfill \\ =& \frac{s}{\left(t2\right)!\left(s+t2\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}a^\gamma \epsilon _\gamma (x_\gamma ,x,x^{})(y_\gamma ,y,y^{})(\overline{\xi }\backslash z,\eta \backslash \{x^{},y^{}\})\left(\xi \eta \right)\overline{\xi }\hfill \\ =& sB.\hfill \end{array}$$
Subsubcase $`(\stackrel{~}{\mathrm{\Xi }},\xi ^{},\eta ^{})=\left[\begin{array}{cc}& \\ & +\end{array}\right]`$. We calculate
$$\begin{array}{cc}& \mathrm{\Lambda }(\left[a\right],\stackrel{~}{e},\text{},\left\{1\right\},\left\{1\right\})\hfill \\ =& \frac{st}{s!t!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{y_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x,x^{})(y_\gamma ,y_1^{})\left(\xi \eta \right)\hfill \\ \stackrel{\text{(}\text{1.5}\text{)}}{=}& \frac{sts!\left(t1\right)!}{s!t!\left(s+t1\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{y_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x,x^{})(y_\gamma ,y_1^{})(\overline{\xi },\eta \backslash x^{})\left(\xi \eta \right)\hfill \\ =& sH.\hfill \end{array}$$
Subsubcase $`(\stackrel{~}{\mathrm{\Xi }},\xi ^{},\eta ^{})=\left[\begin{array}{cc}& +\\ & \end{array}\right]`$. We calculate
$$\begin{array}{cc}& \mathrm{\Lambda }(\left[a\right],\stackrel{~}{e},\text{},\left\{1\right\},\left\{2\right\})\hfill \\ =& \frac{st}{s!t!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{x_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x,x_1^{})(y_\gamma ,y^{})\left(\xi \eta \right)\hfill \\ \stackrel{\text{(}\text{1.5}\text{)}}{=}& \frac{st\left(s1\right)!}{s!t!\left(s+t2\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{x_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x,x_1^{})(y_\gamma ,y^{})(\overline{\xi }\backslash z,\eta \backslash \{x^{},y^{}\})(x_1^{},x^{})\left(\xi \eta \right)\left(\overline{\xi }x_1^{}\right)\hfill \\ =& \frac{1}{\left(t1\right)!\left(s+t2\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{x_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma \hfill \\ & (x_\gamma ,x,x_1^{})(y_\gamma ,y,y^{})(\overline{\xi }\backslash z,\eta \backslash \{x^{},y^{}\})(x_1^{},x^{})(x_1^{},x_\gamma )\left(\xi \eta \right)\left(\overline{\xi }x_1^{}\right)\hfill \\ =& \frac{1}{\left(t1\right)!\left(s+t2\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{x_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x,x^{})(y_\gamma ,y,y^{})(\overline{\xi }\backslash z,\eta \backslash \{x^{},y^{}\})\left(\xi \eta \right)\left(\overline{\xi }x_1^{}\right)\hfill \\ =& uBA.\hfill \end{array}$$
Subsubcase $`(\stackrel{~}{\mathrm{\Xi }},\xi ^{},\eta ^{})=\left[\begin{array}{cc}& +\\ & +\end{array}\right]`$. We calculate
$$\begin{array}{cc}& \mathrm{\Lambda }(\left[a\right],\stackrel{~}{e},\text{},\left\{1\right\},\left\{\right\})\hfill \\ =& \frac{s}{s!t!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{x_1^{}y_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x,x_1^{})(y_\gamma ,y_1^{})\left(\xi \eta \right)\hfill \\ \stackrel{\text{(}\text{1.5}\text{)}}{=}& \frac{s\left(s1\right)!t!}{s!t!\left(s+t1\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{x_1^{}y_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x,x_1^{})(y_\gamma ,y_1^{})(\overline{\xi },\eta \backslash x^{})(x_1^{},x^{})\left(\xi \eta \right)\hfill \\ =& \frac{1}{\left(s+t1\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{x_1^{}y_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x,x_1^{})(y_\gamma ,y_1^{})(\overline{\xi },\eta \backslash x^{})(x_1^{},x^{})(x_1^{},x_\gamma )\left(\xi \eta \right)\hfill \\ =& \frac{u1}{\left(s+t1\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{y_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x,x^{})(y_\gamma ,y_1^{})(\overline{\xi },\eta \backslash x^{})\left(\xi \eta \right)\hfill \\ =& \left(u1\right)H.\hfill \end{array}$$
Subsubcase $`(\stackrel{~}{\mathrm{\Xi }},\xi ^{},\eta ^{})=\left[\begin{array}{cc}+& \\ & \end{array}\right]`$. We calculate
$$\begin{array}{cc}& \mathrm{\Lambda }(\left[a\right],\stackrel{~}{e},\text{},\left\{\right\},\{1,2\})\hfill \\ =& \frac{t\left(t1\right)}{s!t!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{x_1\overline{\xi }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x_1,x^{})(y_\gamma ,y^{})\left(\xi \eta \right)\hfill \\ \stackrel{\text{(}\text{1.5}\text{)}}{=}& \frac{t\left(t1\right)\left(s+1\right)!\left(t2\right)!}{s!t!\left(s+t1\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{x_1\overline{\xi }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x_1,x^{})(y_\gamma ,y^{})(\overline{\xi }\backslash x_1,\eta \backslash \{x^{},y^{}\})\left(\xi \eta \right)\hfill \\ =& \frac{s+1}{\left(s+t1\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{x_1\overline{\xi }}{}a^\gamma \epsilon _\gamma (\overline{\xi },\eta \backslash y^{})(x_\gamma ,x_1)(y_\gamma ,y,y^{})\left(\xi \eta \right)\hfill \\ =& \left(s+1\right)D.\hfill \end{array}$$
Subsubcase $`(\stackrel{~}{\mathrm{\Xi }},\xi ^{},\eta ^{})=\left[\begin{array}{cc}+& +\\ & \end{array}\right]`$. We calculate
$$\begin{array}{cc}& \mathrm{\Lambda }(\left[a\right],\stackrel{~}{e},\text{},\left\{\right\},\left\{2\right\})\hfill \\ =& \frac{t}{s!t!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{x_1\overline{\xi },x_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x_1,x_1^{})(y_\gamma ,y^{})\left(\xi \eta \right)\hfill \\ \stackrel{\text{(}\text{1.5}\text{)}}{=}& \frac{ts!\left(t1\right)!}{s!t!\left(s+t1\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{x_1\overline{\xi },x_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x_1,x_1^{})(y_\gamma ,y^{})(\overline{\xi }\backslash x_1,\eta \backslash \{x^{},y^{}\})(x_1^{},x^{})\left(\xi \eta \right)\hfill \\ =& \frac{1}{\left(s+t1\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{x_1\overline{\xi },x_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (\overline{\xi },\eta \backslash y^{})(x_\gamma ,x_1,x^{},x_1^{})(y_\gamma ,y,y^{})(x_1^{},x^{})(x_1^{},x_\gamma )\left(\xi \eta \right)\hfill \\ =& \frac{u}{\left(s+t1\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{x_1\overline{\xi }}{}a^\gamma \epsilon _\gamma (\overline{\xi },\eta \backslash y^{})(x_\gamma ,x_1)(y_\gamma ,y,y^{})\left(\xi \eta \right)\hfill \\ =& uD.\hfill \end{array}$$
Subsubcases $`\left[\begin{array}{cc}+& \\ & +\end{array}\right]`$, $`\left[\begin{array}{cc}+& +\\ & +\end{array}\right]`$. The $`\mathrm{\Lambda }`$-value is zero by the Garnir relations.
Subcase $`\stackrel{~}{\mathrm{\Xi }}=\left[\begin{array}{cc}& +\\ +& +\end{array}\right]`$. Obtained, by symmetry, from subcase $`\stackrel{~}{\mathrm{\Xi }}=\left[\begin{array}{cc}+& +\\ & +\end{array}\right]`$.
Summing up $`\mathrm{\Lambda }`$-values over subcases and subsubcases, we obtain
$$G_{\left[a\right],\xi ,\eta }^{}f_{\stackrel{~}{e}12}^{}=2A2\left(su\right)B2\left(su+1\right)D+2\left(su+1\right)H.$$
Case $`f_{\stackrel{~}{e}02}^{}`$.
Subcase $`\stackrel{~}{\mathrm{\Xi }}=\left[\begin{array}{cc}& +\\ & +\end{array}\right]`$.
Subsubcase $`(\stackrel{~}{\mathrm{\Xi }},\xi ^{},\eta ^{})=\left[\begin{array}{cc}& \\ & \end{array}\right]`$. We calculate
$$\begin{array}{cc}& \mathrm{\Lambda }(\left[a\right],\stackrel{~}{e},\text{},\left\{\right\},\{1,2\})\hfill \\ =& \frac{t\left(t1\right)}{s!t!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}a^\gamma \epsilon _\gamma (x_\gamma ,x^{})(y_\gamma ,y^{})\left(\xi \eta \right)\hfill \\ \stackrel{\text{(}\text{1.5}\text{)}}{=}& \frac{t\left(t1\right)s!}{s!t!\left(s+t2\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}a^\gamma \epsilon _\gamma (x_\gamma ,x^{})(y_\gamma ,y^{})(\overline{\xi }\backslash z,\eta \backslash \{x^{},y^{}\})\left(\xi \eta \right)\overline{\xi }\hfill \\ =& \frac{1}{\left(t2\right)!\left(s+t2\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}a^\gamma \epsilon _\gamma (x_\gamma ,x,x^{})(y_\gamma ,y,y^{})(\overline{\xi }\backslash z,\eta \backslash \{x^{},y^{}\})\left(\xi \eta \right)\overline{\xi }\hfill \\ =& B.\hfill \end{array}$$
Subsubcase $`(\stackrel{~}{\mathrm{\Xi }},\xi ^{},\eta ^{})=\left[\begin{array}{cc}& \\ & +\end{array}\right]`$. We calculate
$$\begin{array}{cc}& \mathrm{\Lambda }(\left[a\right],\stackrel{~}{e},\text{},\left\{\right\},\left\{1\right\})\hfill \\ =& \frac{t}{s!t!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{y_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x^{})(y_\gamma ,y_1^{})\left(\xi \eta \right)\hfill \\ \stackrel{\text{(}\text{1.5}\text{)}}{=}& \frac{ts!\left(t1\right)!}{s!t!\left(s+t1\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{y_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x^{})(y_\gamma ,y_1^{})(\overline{\xi },\eta \backslash x^{})\left(\xi \eta \right)\hfill \\ =& \frac{1}{\left(s+t1\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{y_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x,x^{})(y_\gamma ,y_1^{})(\overline{\xi },\eta \backslash x^{})\left(\xi \eta \right)\hfill \\ =& H.\hfill \end{array}$$
Subsubcase $`(\stackrel{~}{\mathrm{\Xi }},\xi ^{},\eta ^{})=\left[\begin{array}{cc}& +\\ & \end{array}\right]`$. By symmetry, we obtain from subsubcase $`\left[\begin{array}{cc}& \\ & +\end{array}\right]`$
$$\mathrm{\Lambda }(\left[a\right],\stackrel{~}{e},\text{},\left\{\right\},\left\{2\right\})=H.$$
Subsubcase $`(\stackrel{~}{\mathrm{\Xi }},\xi ^{},\eta ^{})=\left[\begin{array}{cc}& +\\ & +\end{array}\right]`$. The $`\mathrm{\Lambda }`$-value is zero by the Garnir relations.
Summing up $`\mathrm{\Lambda }`$-values over subcases and subsubcases, we obtain
$$G_{\left[a\right],\xi ,\eta }^{}f_{\stackrel{~}{e}02}^{}=B2H.$$
Case $`f_{\stackrel{~}{e}21}^{}`$.
Subcase $`\stackrel{~}{\mathrm{\Xi }}=\left[\begin{array}{cc}+& +\\ +& \end{array}\right]`$.
Subsubcase $`(\stackrel{~}{\mathrm{\Xi }},\xi ^{},\eta ^{})=\left[\begin{array}{cc}& \\ & \end{array}\right]`$. We calculate
$$\begin{array}{cc}& \mathrm{\Lambda }(\left[a\right],\stackrel{~}{e},\text{},\{1,2\},\left\{1\right\})\hfill \\ =& \frac{s\left(s1\right)t}{s!t!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}a^\gamma \epsilon _\gamma (x_\gamma ,x,x^{})(y_\gamma ,y)\left(\xi \eta \right)\hfill \\ \stackrel{\text{(}\text{1.5}\text{)}}{=}& \frac{s\left(s1\right)t\left(s1\right)!}{s!t!\left(s+t2\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}a^\gamma \epsilon _\gamma (x_\gamma ,x,x^{})(y_\gamma ,y)(\overline{\xi }\backslash z,\eta \backslash \{x^{},y^{}\})(y_\gamma ,y^{})\left(\xi \eta \right)\left(\overline{\xi }y_\gamma \right)\hfill \\ =& \frac{\left(s1\right)}{\left(t1\right)!\left(s+t2\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}a^\gamma \epsilon _\gamma (x_\gamma ,x,x^{})(y_\gamma ,y,y^{})(\overline{\xi }\backslash z,\eta \backslash \{x^{},y^{}\})\left(\xi \eta \right)\left(\overline{\xi }y_\gamma \right)\hfill \\ =& \left(s1\right)\left(BC_2\right).\hfill \end{array}$$
Subsubcase $`(\stackrel{~}{\mathrm{\Xi }},\xi ^{},\eta ^{})=\left[\begin{array}{cc}& +\\ & \end{array}\right]`$. We calculate
$$\begin{array}{cc}& \mathrm{\Lambda }(\left[a\right],\stackrel{~}{e},\text{},\{1,2\},\left\{\right\})\hfill \\ =& \frac{s\left(s1\right)}{s!t!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{x_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x,x_1^{})(y_\gamma ,y)\left(\xi \eta \right)\text{ }\hfill \end{array}$$
$$\begin{array}{cc}\stackrel{\text{(}\text{1.5}\text{)}}{=}& \frac{s\left(s1\right)\left(s2\right)!}{s!t!\left(s+t2\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{x_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma \hfill \\ & (x_\gamma ,x,x_1^{})(y_\gamma ,y)(\overline{\xi }\backslash z,\eta \backslash \{x^{},y^{}\})(x_1^{},x^{})(y_\gamma ,y^{})\left(\xi \eta \right)\left(\overline{\xi }x_1^{}y_\gamma \right)\hfill \\ =& \frac{1}{t!\left(s+t2\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{x_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x,x_1^{})(y_\gamma ,y)\hfill \\ & (\overline{\xi }\backslash z,\eta \backslash \{x^{},y^{}\})(x_1^{},x^{})(y_\gamma ,y^{})(x_1^{},x_\gamma )\left(\xi \eta \right)\left(\overline{\xi }x_1^{}y_\gamma \right)\hfill \\ =& \frac{1}{t!\left(s+t2\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{x_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x,x^{})(y_\gamma ,y,y^{})(\overline{\xi }\backslash z,\eta \backslash \{x^{},y^{}\})\left(\xi \eta \right)\left(\overline{\xi }x_1^{}y_\gamma \right)\hfill \\ =& \left(uBuC_2A\right).\hfill \end{array}$$
Subsubcase $`(\stackrel{~}{\mathrm{\Xi }},\xi ^{},\eta ^{})=\left[\begin{array}{cc}& \\ +& \end{array}\right]`$. We calculate
$$\begin{array}{cc}& \mathrm{\Lambda }(\left[a\right],\stackrel{~}{e},\text{},\left\{1\right\},\left\{1\right\})\hfill \\ =& \frac{st}{s!t!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{y_1\overline{\xi }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x,x^{})(y_\gamma ,y_1)\left(\xi \eta \right)\hfill \\ \stackrel{\text{(}\text{1.5}\text{)}}{=}& \frac{sts!\left(t1\right)!}{s!t!\left(s+t1\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{y_1\overline{\xi }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x,x^{})(y_\gamma ,y_1)(\overline{\xi }\backslash y_1,\eta \backslash \{x^{},y^{}\})(y_\gamma ,y^{})\left(\xi \eta \right)\hfill \\ =& \frac{s}{\left(s+t1\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{y_1\overline{\xi }}{}a^\gamma \epsilon _\gamma (\overline{\xi },\eta \backslash x^{})(x_\gamma ,x,x^{})(y_\gamma ,y_1)\left(\xi \eta \right)\hfill \\ =& sD.\hfill \end{array}$$
Subsubcase $`(\stackrel{~}{\mathrm{\Xi }},\xi ^{},\eta ^{})=\left[\begin{array}{cc}& +\\ +& \end{array}\right]`$. We calculate
$$\begin{array}{cc}& \mathrm{\Lambda }(\left[a\right],\stackrel{~}{e},\text{},\left\{1\right\},\left\{\right\})\hfill \\ =& \frac{s}{s!t!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{y_1\overline{\xi },x_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x,x_1^{})(y_\gamma ,y_1)\left(\xi \eta \right)\hfill \\ \stackrel{\text{(}\text{1.5}\text{)}}{=}& \frac{s\left(s1\right)!t!}{s!t!\left(s+t1\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{y_1\overline{\xi },x_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma \hfill \\ & (x_\gamma ,x,x_1^{})(y_\gamma ,y_1)(\overline{\xi }\backslash y_1,\eta \backslash \{x^{},y^{}\})(y_\gamma ,y^{})(x_1^{},x^{})\left(\xi \eta \right)\hfill \\ =& \frac{1}{\left(s+t1\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{y_1\overline{\xi },x_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (\overline{\xi },\eta \backslash x^{})(x_\gamma ,x,x_1^{})(y_\gamma ,y_1)(x_1^{},x^{})(x_\gamma ,x_1^{})\left(\xi \eta \right)\hfill \\ =& \frac{u}{\left(s+t1\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{y_1\overline{\xi }}{}a^\gamma \epsilon _\gamma (\overline{\xi },\eta \backslash x^{})(x_\gamma ,x,x^{})(y_\gamma ,y_1)\left(\xi \eta \right)\hfill \\ =& uD.\hfill \end{array}$$
Subsubcase $`(\stackrel{~}{\mathrm{\Xi }},\xi ^{},\eta ^{})=\left[\begin{array}{cc}+& \\ & \end{array}\right]`$. We calculate
$$\begin{array}{cc}& \mathrm{\Lambda }(\left[a\right],\stackrel{~}{e},\text{},\left\{2\right\},\left\{1\right\})\hfill \\ =& \frac{st}{s!t!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{x_1\overline{\xi }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x_1,x^{})(y_\gamma ,y)\left(\xi \eta \right)\hfill \\ \stackrel{\text{(}\text{1.5}\text{)}}{=}& \frac{sts!\left(t1\right)!}{s!t!\left(s+t1\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{x_1\overline{\xi }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x_1,x^{})(y_\gamma ,y)(\overline{\xi }\backslash x_1,\eta \backslash \{x^{},y^{}\})(y_\gamma ,y^{})\left(\xi \eta \right)\hfill \\ =& \frac{s}{\left(s+t1\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{x_1\overline{\xi }}{}a^\gamma \epsilon _\gamma (\overline{\xi },\eta \backslash y^{})(x_\gamma ,x_1)(y_\gamma ,y,y^{})\left(\xi \eta \right)\hfill \\ =& sD.\hfill \end{array}$$
Subsubcase $`(\stackrel{~}{\mathrm{\Xi }},\xi ^{},\eta ^{})=\left[\begin{array}{cc}+& +\\ & \end{array}\right]`$. We calculate
$$\begin{array}{cc}& \mathrm{\Lambda }(\left[a\right],\stackrel{~}{e},\text{},\left\{2\right\},\left\{\right\})\hfill \\ =& \frac{s}{s!t!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{x_1\overline{\xi },x_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x_1,x_1^{})(y_\gamma ,y)\left(\xi \eta \right)\hfill \\ \stackrel{\text{(}\text{1.5}\text{)}}{=}& \frac{s\left(s1\right)!t!}{s!t!\left(s+t1\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{x_1\overline{\xi },x_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma \hfill \\ & (x_\gamma ,x_1,x_1^{})(y_\gamma ,y)(\overline{\xi }\backslash x_1,\eta \backslash \{x^{},y^{}\})(y_\gamma ,y^{})(x_1^{},x^{})\left(\xi \eta \right)\hfill \\ =& \frac{1}{\left(s+t1\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{x_1\overline{\xi },x_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (\overline{\xi },\eta \backslash y^{})(x_\gamma ,x_1,x^{},x_1^{})(y_\gamma ,y,y^{})(x_1^{},x^{})(x_\gamma ,x_1^{})\left(\xi \eta \right)\hfill \\ =& \frac{u}{\left(s+t1\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{x_1\overline{\xi }}{}a^\gamma \epsilon _\gamma (\overline{\xi },\eta \backslash y^{})(x_\gamma ,x_1)(y_\gamma ,y,y^{})\left(\xi \eta \right)\hfill \\ =& uD.\hfill \end{array}$$
Subsubcases $`\left[\begin{array}{cc}+& \\ +& \end{array}\right]`$, $`\left[\begin{array}{cc}+& +\\ +& \end{array}\right]`$. The $`\mathrm{\Lambda }`$-value is zero by the Garnir relations.
Subcase $`\stackrel{~}{\mathrm{\Xi }}=\left[\begin{array}{cc}+& \\ +& +\end{array}\right]`$. Obtained, by symmetry, from subcase $`\stackrel{~}{\mathrm{\Xi }}=\left[\begin{array}{cc}+& +\\ +& \end{array}\right]`$.
Summing up $`\mathrm{\Lambda }`$-values over subcases and subsubcases, we obtain
$$G_{\left[a\right],\xi ,\eta }^{}f_{\stackrel{~}{e}21}^{}=2A+2\left(su1\right)B\left(su1\right)\left(C_1+C_2\right)+4\left(su\right)D.$$
Case $`f_{\stackrel{~}{e}11}^{}`$.
Subcase $`\stackrel{~}{\mathrm{\Xi }}=\left[\begin{array}{cc}+& +\\ & \end{array}\right]`$.
Subsubcase $`(\stackrel{~}{\mathrm{\Xi }},\xi ^{},\eta ^{})=\left[\begin{array}{cc}& \\ & \end{array}\right]`$. We calculate
$$\begin{array}{cc}& \mathrm{\Lambda }(\left[a\right],\stackrel{~}{e},\text{},\left\{1\right\},\left\{1\right\})\hfill \\ =& \frac{st}{s!t!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}a^\gamma \epsilon _\gamma (x_\gamma ,x,x^{})\left(\xi \eta \right)\hfill \\ \stackrel{\text{(}\text{1.5}\text{)}}{=}& \frac{sts!\left(t1\right)!}{s!t!\left(s+t1\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}a^\gamma \epsilon _\gamma (x_\gamma ,x,x^{})(\overline{\xi },\eta \backslash x^{})\left(\xi \eta \right)\hfill \\ =& sF_1.\hfill \end{array}$$
Subsubcase $`(\stackrel{~}{\mathrm{\Xi }},\xi ^{},\eta ^{})=\left[\begin{array}{cc}& +\\ & \end{array}\right]`$. We calculate
$$\begin{array}{cc}& \mathrm{\Lambda }(\left[a\right],\stackrel{~}{e},\text{},\left\{1\right\},\left\{\right\})\hfill \\ =& \frac{s}{s!t!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{x_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x,x_1^{})\left(\xi \eta \right)\hfill \\ \stackrel{\text{(}\text{1.5}\text{)}}{=}& \frac{s\left(s1\right)!t!}{s!t!\left(s+t1\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{x_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x,x_1^{})(\overline{\xi },\eta \backslash x^{})(x_1^{},x^{})\left(\xi \eta \right)\hfill \\ =& \frac{1}{\left(s+t1\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{x_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x,x_1^{})(\overline{\xi },\eta \backslash x^{})(x_1^{},x^{})(x_1^{},x_\gamma )\left(\xi \eta \right)\hfill \\ =& \frac{u}{\left(s+t1\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}a^\gamma \epsilon _\gamma (\overline{\xi },\eta \backslash x^{})(x_\gamma ,x,x^{})\left(\xi \eta \right)\hfill \\ =& uF_1.\hfill \end{array}$$
Subsubcases $`\left[\begin{array}{cc}+& \\ & \end{array}\right]`$, $`\left[\begin{array}{cc}+& +\\ & \end{array}\right]`$. The $`\mathrm{\Lambda }`$-value is zero by the Garnir relations.
Subcase $`\stackrel{~}{\mathrm{\Xi }}=\left[\begin{array}{cc}& \\ +& +\end{array}\right]`$. Obtained, by symmetry, from subcase $`\stackrel{~}{\mathrm{\Xi }}=\left[\begin{array}{cc}+& +\\ & \end{array}\right]`$.
Subcase $`\stackrel{~}{\mathrm{\Xi }}=\left[\begin{array}{cc}+& \\ & +\end{array}\right]`$.
Subsubcase $`(\stackrel{~}{\mathrm{\Xi }},\xi ^{},\eta ^{})=\left[\begin{array}{cc}& \\ & \end{array}\right]`$. We calculate
$$\begin{array}{cc}& \mathrm{\Lambda }(\left[a\right],\stackrel{~}{e},\text{},\left\{1\right\},\left\{2\right\})\hfill \\ =& \frac{st}{s!t!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}a^\gamma \epsilon _\gamma (x_\gamma ,x)(y_\gamma ,y^{})\left(\xi \eta \right)\hfill \\ \stackrel{\text{(}\text{1.5}\text{)}}{=}& \frac{st\left(s1\right)!}{s!t!\left(s+t2\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}a^\gamma \epsilon _\gamma (x_\gamma ,x)(y_\gamma ,y^{})(\overline{\xi }\backslash z,\eta \backslash \{x^{},y^{}\})(x_\gamma ,x^{})\left(\xi \eta \right)\left(\overline{\xi }x_\gamma \right)\hfill \\ =& \frac{1}{\left(t1\right)!\left(s+t2\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}a^\gamma \epsilon _\gamma (x_\gamma ,x,x^{})(y_\gamma ,y,y^{})(\overline{\xi }\backslash z,\eta \backslash \{x^{},y^{}\})\left(\xi \eta \right)\left(\overline{\xi }x_\gamma \right)\hfill \\ =& \left(BC_1\right).\hfill \end{array}$$
Subsubcase $`(\stackrel{~}{\mathrm{\Xi }},\xi ^{},\eta ^{})=\left[\begin{array}{cc}+& \\ & \end{array}\right]`$. We calculate
$$\begin{array}{cc}& \mathrm{\Lambda }(\left[a\right],\stackrel{~}{e},\text{},\left\{\right\},\left\{2\right\})\hfill \\ =& \frac{t}{s!t!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{x_1\overline{\xi }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x_1)(y_\gamma ,y^{})\left(\xi \eta \right)\hfill \\ \stackrel{\text{(}\text{1.5}\text{)}}{=}& \frac{ts!\left(t1\right)!}{s!t!\left(s+t1\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{x_1\overline{\xi }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x_1)(y_\gamma ,y^{})(\overline{\xi }\backslash x_1,\eta \backslash \{x^{},y^{}\})(x_\gamma ,x^{})\left(\xi \eta \right)\hfill \\ =& \frac{1}{\left(s+t1\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{x_1\overline{\xi }}{}a^\gamma \epsilon _\gamma (\overline{\xi },\eta \backslash y^{})(x_\gamma ,x_1)(y_\gamma ,y,y^{})\left(\xi \eta \right)\hfill \\ =& D.\hfill \end{array}$$
Subsubcase $`(\stackrel{~}{\mathrm{\Xi }},\xi ^{},\eta ^{})=\left[\begin{array}{cc}& \\ & +\end{array}\right]`$. We calculate
$$\begin{array}{cc}& \mathrm{\Lambda }(\left[a\right],\stackrel{~}{e},\text{},\left\{1\right\},\left\{\right\})\hfill \\ =& \frac{s}{s!t!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{y_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x)(y_\gamma ,y_1^{})\left(\xi \eta \right)\hfill \\ \stackrel{\text{(}\text{1.5}\text{)}}{=}& \frac{s\left(s1\right)!t!}{s!t!\left(s+t1\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{y_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x)(y_\gamma ,y_1^{})(\overline{\xi },\eta \backslash x^{})(x_\gamma ,x^{})\left(\xi \eta \right)\hfill \\ =& \frac{1}{\left(s+t1\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{y_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x,x^{})(y_\gamma ,y_1^{})(\overline{\xi },\eta \backslash x^{})\left(\xi \eta \right)\hfill \\ =& H.\hfill \end{array}$$
Subsubcase $`(\stackrel{~}{\mathrm{\Xi }},\xi ^{},\eta ^{})=\left[\begin{array}{cc}+& \\ & +\end{array}\right]`$. The $`\mathrm{\Lambda }`$-value is zero by the Garnir relations.
Subcase $`\stackrel{~}{\mathrm{\Xi }}=\left[\begin{array}{cc}& +\\ +& \end{array}\right]`$. Obtained, by symmetry, from subcase $`\stackrel{~}{\mathrm{\Xi }}=\left[\begin{array}{cc}+& \\ & +\end{array}\right]`$.
Summing up $`\mathrm{\Lambda }`$-values over subcases and subsubcases, we obtain
$$G_{\left[a\right],\xi ,\eta }^{}f_{\stackrel{~}{e}11}^{}=2B+\left(C_1+C_2\right)2D+2H+\left(su\right)\left(F_1+F_2\right).$$
Case $`f_{\stackrel{~}{e}01}^{}`$.
Subcase $`\stackrel{~}{\mathrm{\Xi }}=\left[\begin{array}{cc}& +\\ & \end{array}\right]`$.
Subsubcase $`(\stackrel{~}{\mathrm{\Xi }},\xi ^{},\eta ^{})=\left[\begin{array}{cc}& \\ & \end{array}\right]`$. We calculate
$$\begin{array}{cc}& \mathrm{\Lambda }(\left[a\right],\stackrel{~}{e},\text{},\left\{\right\},\left\{1\right\})\hfill \\ =& \frac{t}{s!t!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}a^\gamma \epsilon _\gamma (x_\gamma ,x^{})\left(\xi \eta \right)\hfill \end{array}$$
$$\begin{array}{cc}\stackrel{\text{(}\text{1.5}\text{)}}{=}& \frac{ts!\left(t1\right)!}{s!t!\left(s+t1\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}a^\gamma \epsilon _\gamma (x_\gamma ,x^{})(\overline{\xi },\eta \backslash x^{})\left(\xi \eta \right)\hfill \\ =& \frac{1}{\left(s+t1\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}a^\gamma \epsilon _\gamma (x_\gamma ,x,x^{})(\overline{\xi },\eta \backslash x^{})\left(\xi \eta \right)\hfill \\ =& F_1.\hfill \end{array}$$
Subsubcase $`(\stackrel{~}{\mathrm{\Xi }},\xi ^{},\eta ^{})=\left[\begin{array}{cc}& +\\ & \end{array}\right]`$. The $`\mathrm{\Lambda }`$-value is zero by the Garnir relations.
Subcase $`\stackrel{~}{\mathrm{\Xi }}=\left[\begin{array}{cc}& \\ & +\end{array}\right]`$. Obtained, by symmetry, from subcase $`\stackrel{~}{\mathrm{\Xi }}=\left[\begin{array}{cc}& +\\ & \end{array}\right]`$.
Summing up $`\mathrm{\Lambda }`$-values over subcases and subsubcases, we obtain
$$G_{\left[a\right],\xi ,\eta }^{}f_{\stackrel{~}{e}01}^{}=\left(F_1+F_2\right).$$
Case $`f_{\stackrel{~}{e}20}^{}`$.
Subcase $`\stackrel{~}{\mathrm{\Xi }}=\left[\begin{array}{cc}+& \\ +& \end{array}\right]`$.
Subsubcase $`(\stackrel{~}{\mathrm{\Xi }},\xi ^{},\eta ^{})=\left[\begin{array}{cc}& \\ & \end{array}\right]`$. We calculate
$$\begin{array}{cc}& \mathrm{\Lambda }(\left[a\right],\stackrel{~}{e},\text{},\{1,2\},\left\{\right\})\hfill \\ =& \frac{s\left(s1\right)}{s!t!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}a^\gamma \epsilon _\gamma (x_\gamma ,x)(y_\gamma ,y)\left(\xi \eta \right)\hfill \\ \stackrel{\text{(}\text{1.5}\text{)}}{=}& \frac{s\left(s1\right)\left(s2\right)!}{s!t!\left(s+t2\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}a^\gamma \epsilon _\gamma \hfill \\ & (x_\gamma ,x)(y_\gamma ,y)(\overline{\xi }\backslash z,\eta \backslash \{x^{},y^{}\})(x_\gamma ,x^{})(y_\gamma ,y^{})\left(\xi \eta \right)\left(\overline{\xi }x_\gamma y_\gamma \right)\hfill \\ =& \frac{1}{t!\left(s+t2\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}a^\gamma \epsilon _\gamma (x_\gamma ,x,x^{})(y_\gamma ,y,y^{})(\overline{\xi }\backslash z,\eta \backslash \{x^{},y^{}\})\left(\xi \eta \right)\left(\overline{\xi }x_\gamma y_\gamma \right)\hfill \\ =& BC_1C_2.\hfill \end{array}$$
Subsubcase $`(\stackrel{~}{\mathrm{\Xi }},\xi ^{},\eta ^{})=\left[\begin{array}{cc}& \\ +& \end{array}\right]`$. We calculate
$$\begin{array}{cc}& \mathrm{\Lambda }(\left[a\right],\stackrel{~}{e},\text{},\left\{1\right\},\left\{\right\})\hfill \\ =& \frac{s}{s!t!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{y_1\overline{\xi }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x)(y_\gamma ,y_1)\left(\xi \eta \right)\hfill \\ \stackrel{\text{(}\text{1.5}\text{)}}{=}& \frac{s\left(s1\right)!t!}{s!t!\left(s+t1\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{y_1\overline{\xi }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x)(y_\gamma ,y_1)(\overline{\xi }\backslash y_1,\eta \backslash \{x^{},y^{}\})(x_\gamma ,x^{})(y_\gamma ,y^{})\left(\xi \eta \right)\hfill \\ =& \frac{1}{\left(s+t1\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}\underset{y_1\overline{\xi }}{}a^\gamma \epsilon _\gamma (\overline{\xi },\eta \backslash x^{})(x_\gamma ,x,x^{})(y_\gamma ,y_1)\left(\xi \eta \right)\hfill \\ =& D.\hfill \end{array}$$
Subsubcase $`(\stackrel{~}{\mathrm{\Xi }},\xi ^{},\eta ^{})=\left[\begin{array}{cc}+& \\ & \end{array}\right]`$. By symmetry, we obtain from subsubcase $`\left[\begin{array}{cc}& \\ +& \end{array}\right]`$
$$\mathrm{\Lambda }(\left[a\right],\stackrel{~}{e},\text{},\left\{2\right\},\left\{\right\})=D.$$
Subsubcase $`(\stackrel{~}{\mathrm{\Xi }},\xi ^{},\eta ^{})=\left[\begin{array}{cc}+& \\ +& \end{array}\right]`$. The $`\mathrm{\Lambda }`$-value is zero by the Garnir relations.
Summing up $`\mathrm{\Lambda }`$-values over subcases and subsubcases, we obtain
$$G_{\left[a\right],\xi ,\eta }^{}f_{\stackrel{~}{e}20}^{}=B\left(C_1+C_2\right)+2D.$$
Case $`f_{\stackrel{~}{e}10}^{}`$.
Subcase $`\stackrel{~}{\mathrm{\Xi }}=\left[\begin{array}{cc}+& \\ & \end{array}\right]`$.
Subsubcase $`(\stackrel{~}{\mathrm{\Xi }},\xi ^{},\eta ^{})=\left[\begin{array}{cc}& \\ & \end{array}\right]`$. We calculate
$$\begin{array}{cc}& \mathrm{\Lambda }(\left[a\right],\stackrel{~}{e},\text{},\left\{1\right\},\left\{\right\})\hfill \\ =& \frac{s}{s!t!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}a^\gamma \epsilon _\gamma (x_\gamma ,x)\left(\xi \eta \right)\hfill \\ \stackrel{\text{(}\text{1.5}\text{)}}{=}& \frac{s\left(s1\right)!t!}{s!t!\left(s+t1\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}a^\gamma \epsilon _\gamma (x_\gamma ,x)(\overline{\xi },\eta \backslash x^{})(x_\gamma ,x^{})\left(\xi \eta \right)\hfill \\ =& \frac{1}{\left(s+t1\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}00\right)}{}a^\gamma \epsilon _\gamma (x_\gamma ,x,x^{})(\overline{\xi },\eta \backslash x^{})\left(\xi \eta \right)\hfill \\ =& F_1.\hfill \end{array}$$
Subsubcase $`(\stackrel{~}{\mathrm{\Xi }},\xi ^{},\eta ^{})=\left[\begin{array}{cc}+& \\ & \end{array}\right]`$. The $`\mathrm{\Lambda }`$-value is zero by the Garnir relations.
Subcase $`\stackrel{~}{\mathrm{\Xi }}=\left[\begin{array}{cc}& \\ +& \end{array}\right]`$. Obtained, by symmetry, from subcase $`\stackrel{~}{\mathrm{\Xi }}=\left[\begin{array}{cc}+& \\ & \end{array}\right]`$.
Summing up $`\mathrm{\Lambda }`$-values over subcases and subsubcases, we obtain
$$G_{\left[a\right],\xi ,\eta }^{}f_{\stackrel{~}{e}10}^{}=\left(F_1+F_2\right).$$
Case $`f_{\stackrel{~}{e}00}^{}`$. The Garnir relations yield
$$G_{\left[a\right],\xi ,\eta }^{}f_{\stackrel{~}{e}00}^{}=0.$$
We note that $`su=X_pX_{p+1}`$ and evaluate the linear combination
$$\begin{array}{c}\underset{\alpha ,\beta [0,2]}{}X_p^{\left(2\alpha \right)}X_{p+1}^{\left(2\beta \right)}G_{\left[a\right],\xi ,\eta }^{}f_{\stackrel{~}{e}\alpha \beta }^{}\hfill \\ \begin{array}{cccc}=& 1\hfill & 1\hfill & (2(X_pX_{p+1})A+(X_pX_{p+1})(X_pX_{p+1}1)B\hfill \\ & & & +2(X_pX_{p+1})(X_pX_{p+1}+1)D)\hfill \\ +& X_p\hfill & 1\hfill & (2A2(X_pX_{p+1})B2(X_pX_{p+1}+1)D\hfill \\ & & & +2(X_pX_{p+1}+1)H)\hfill \\ +& X_p\left(X_p+1\right)\hfill & 1\hfill & \left(B2H\right)\hfill \\ +& 1\hfill & X_{p+1}\hfill & (2A+2(X_pX_{p+1}1)B\hfill \\ & & & (X_pX_{p+1}1)(C_1+C_2)+4(X_pX_{p+1})D)\hfill \\ +& X_p\hfill & X_{p+1}\hfill & (2B+(C_1+C_2)2D+2H\hfill \\ & & & +(X_pX_{p+1})(F_1+F_2))\hfill \\ +& X_p\left(X_p+1\right)\hfill & X_{p+1}\hfill & \left(\left(F_1+F_2\right)\right)\hfill \\ +& 1\hfill & X_{p+1}\left(X_{p+1}+1\right)\hfill & \left(B\left(C_1+C_2\right)+2D\right)\hfill \\ +& X_p\hfill & X_{p+1}\left(X_{p+1}+1\right)\hfill & \left(\left(F_1+F_2\right)\right)\hfill \\ +& X_p\left(X_p+1\right)\hfill & X_{p+1}\left(X_{p+1}+1\right)\hfill & \left(0\right)\hfill \\ =& 0.\hfill & & \end{array}\hfill \end{array}$$
###### Lemma A.2 (full version of (2.10))
Suppose $`g=p<k`$ and $`s,t2`$. There exist elements $`x^{},y^{}\eta `$, $`x^{}y^{}`$, $`z\overline{\xi }`$, which we choose and fix. For $`\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e}0)`$, we let $`x_\gamma :=a_{j,\overline{\gamma }(j,1)}`$, where $`j[g+2,k+1]`$ is minimal with $`1\times j\mathrm{\Xi }_\gamma `$, and $`y_\gamma :=a_{j,\overline{\gamma }(j,2)}`$, where $`j[g+2,k+1]`$ is minimal with $`2\times j\mathrm{\Xi }_\gamma `$. I.e. we pick the entries $`x_\gamma ,y_\gamma `$ that ‘cross the column’ $`g+1`$ under the operation of $`\gamma `$. The set of maps
$$\begin{array}{ccc}\hfill [1,\lambda _1]\backslash \left\{g+1\right\}& \text{}& [0,2]\hfill \\ \hfill j& \text{}& \stackrel{~}{e}_j\hfill \end{array}$$
that send $`g`$ and $`k+1`$ to $`\stackrel{~}{e}_g=\stackrel{~}{e}_{k+1}=2`$, and that map $`j[1,\lambda _1]\backslash [g,k+1]`$ to $`e_j=0`$, is denoted by $`\stackrel{~}{E}`$. For $`\stackrel{~}{e}\stackrel{~}{E}`$ and $`\beta [0,2]`$, we denote by $`\stackrel{~}{e}\beta `$ the prolongation of $`\stackrel{~}{e}`$ to $`[g,k+1]`$ by $`(\stackrel{~}{e}\beta )_{g+1}=\beta `$. For $`\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e}0)`$, we write
$$U_\gamma :=a^\gamma \epsilon _\gamma (\overline{\xi }\backslash z,\eta \backslash \{x^{},y^{}\})(x_\gamma ,x^{})(y_\gamma ,y^{})$$
and let
$$\begin{array}{ccc}A_{\left[a\right],\xi ,\eta ,\stackrel{~}{e}}\hfill & :=& \underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}0\right)}{}U_\gamma \underset{w^{}\overline{\eta }}{}(w^{},z)\hfill \\ B_{\left[a\right],\xi ,\eta ,\stackrel{~}{e}}\hfill & :=& \underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}0\right)}{}U_\gamma \left(1\underset{w\overline{\xi }\backslash z}{}(w,z)\right)\hfill \\ C_{1,\left[a\right],\xi ,\eta ,\stackrel{~}{e}}\hfill & :=& \underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}0\right)}{}U_\gamma (z,x_\gamma )\hfill \\ C_{2,\left[a\right],\xi ,\eta ,\stackrel{~}{e}}\hfill & :=& \underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}0\right)}{}U_\gamma (z,y_\gamma ).\hfill \end{array}$$
We obtain
$$\begin{array}{c}\hfill G_{\left[a\right],\xi ,\eta }^{}f^{}=(X_g+2)\underset{\stackrel{~}{e}\stackrel{~}{E}}{}\left(\underset{j[g+2,k]}{}X_j^{\left(2\stackrel{~}{e}_j\right)}\right)(2A_{\left[a\right],\xi ,\eta ,\stackrel{~}{e}}+(X_g+1)B_{\left[a\right],\xi ,\eta ,\stackrel{~}{e}}\\ \hfill X_{g+1}(C_{1,\left[a\right],\xi ,\eta ,\stackrel{~}{e}}+C_{2,\left[a\right],\xi ,\eta ,\stackrel{~}{e}})).\end{array}$$
Having fixed a map $`\stackrel{~}{e}\stackrel{~}{E}`$, we need to evaluate the expression
$$\underset{\beta [0,2]}{}X_{g+1}^{\left(2\beta \right)}G_{\left[a\right],\xi ,\eta }^{}f_{\stackrel{~}{e}\beta }^{}.$$
Given $`\beta [0,2]`$, given a subset $`\stackrel{~}{\mathrm{\Xi }}[1,2]\times \left\{g+1\right\}`$ of cardinality $`\mathrm{\#}\stackrel{~}{\mathrm{\Xi }}=\beta `$ and given a subset $`\eta ^{}[1,2]`$ such that
$$()\eta ^{}\times \left(g+1\right)\stackrel{~}{\mathrm{\Xi }},$$
we let
$$\dot{\mathrm{\Gamma }}(\stackrel{~}{e},\stackrel{~}{\mathrm{\Xi }},\eta ^{}):=\left\{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}\beta \right)\right|\mathrm{\Xi }_\gamma \left([1,2]\times \left\{g+1\right\}\right)=\stackrel{~}{\mathrm{\Xi }},\gamma ^1\left(\eta \right)=\eta ^{}\times \left\{g+1\right\}\}$$
and form the partial sum
$$\mathrm{\Lambda }(\left[a\right],\stackrel{~}{e},\stackrel{~}{\mathrm{\Xi }},\eta ^{}):=\frac{1}{s!t!}\underset{\gamma \dot{\mathrm{\Gamma }}(\stackrel{~}{e},\stackrel{~}{\mathrm{\Xi }},\eta ^{})}{}a^\gamma \epsilon _\gamma \left(\xi \eta \right)$$
so that we can recover
$$G_{\left[a\right],\xi ,\eta }^{}f_{\stackrel{~}{e}\beta }^{}=\underset{\stackrel{~}{\mathrm{\Xi }}[1,2]\times \left\{g+1\right\}}{}\underset{\eta ^{}[1,2]\text{ subject to }()}{}\mathrm{\Lambda }(\left[a\right],\stackrel{~}{e},\stackrel{~}{\mathrm{\Xi }},\eta ^{}).$$
Calculation of $`\mathrm{\Lambda }`$-values. To distinguish subcases and subsubcases, we adapt the according notation of (A.1).
Case $`f_{\stackrel{~}{e}2}^{}`$.
Subcase $`\stackrel{~}{\mathrm{\Xi }}=\left[\begin{array}{c}+\\ +\end{array}\right]`$.
Subsubcase $`(\stackrel{~}{\mathrm{\Xi }},\eta ^{})=\left[\begin{array}{c}\\ \end{array}\right]`$. We calculate
$$\begin{array}{cc}& \mathrm{\Lambda }(\left[a\right],\stackrel{~}{e},\text{},\{1,2\})\hfill \\ =& \frac{t\left(t1\right)}{s!t!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}0\right)}{}a^\gamma \epsilon _\gamma (x_\gamma ,x^{})(y_\gamma ,y^{})\left(\xi \eta \right)\hfill \\ \stackrel{\text{(}\text{1.4}\text{)}}{=}& \frac{t\left(t1\right)\left(s+2\right)!}{s!t!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}0\right)}{}a^\gamma \epsilon _\gamma (x_\gamma ,x^{})(y_\gamma ,y^{})(\overline{\xi }\backslash z,\eta \backslash \{x^{},y^{}\})\overline{\xi }\hfill \\ =& \left(s+2\right)\left(s+1\right)B_{\left[a\right],\xi ,\eta ,\stackrel{~}{e}}.\hfill \end{array}$$
Subsubcase $`(\stackrel{~}{\mathrm{\Xi }},\eta ^{})=\left[\begin{array}{c}\\ +\end{array}\right]`$. We calculate
$$\begin{array}{cc}& \mathrm{\Lambda }(\left[a\right],\stackrel{~}{e},\text{},\left\{1\right\})\hfill \\ =& \frac{t}{s!t!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}0\right)}{}\underset{y_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x^{})(y_\gamma ,y_1^{})\left(\xi \eta \right)\hfill \\ \stackrel{\text{(}\text{1.4}\text{)}}{=}& \frac{t\left(s+1\right)!}{s!t!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}0\right)}{}\underset{y_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x^{})(y_\gamma ,y_1^{})(\overline{\xi }\backslash z,\eta \backslash \{x^{},y^{}\})(y_1^{},y^{})\left(\overline{\xi }y_1^{}\right)\hfill \\ =& \frac{\left(s+1\right)}{\left(t1\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}0\right)}{}\underset{y_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x^{})(y_\gamma ,y_1^{})(\overline{\xi }\backslash z,\eta \backslash \{x^{},y^{}\})(y_1^{},y^{})(y_1^{},y_\gamma )\left(\overline{\xi }y_1^{}\right)\hfill \\ =& \frac{\left(s+1\right)}{\left(t1\right)!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}0\right)}{}\underset{y_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x^{})(\overline{\xi }\backslash z,\eta \backslash \{x^{},y^{}\})(y_\gamma ,y^{})\left(\overline{\xi }y_1^{}\right)\hfill \\ =& \left(s+1\right)\left(uB_{\left[a\right],\xi ,\eta ,\stackrel{~}{e}}A_{\left[a\right],\xi ,\eta ,\stackrel{~}{e}}\right).\hfill \end{array}$$
Subsubcase $`(\stackrel{~}{\mathrm{\Xi }},\eta ^{})=\left[\begin{array}{c}+\\ \end{array}\right]`$. By symmetry, we obtain from subsubcase $`\left[\begin{array}{c}\\ +\end{array}\right]`$
$$\mathrm{\Lambda }(\left[a\right],\stackrel{~}{e},\text{},\left\{2\right\})=\left(s+1\right)\left(uB_{\left[a\right],\xi ,\eta ,\stackrel{~}{e}}A_{\left[a\right],\xi ,\eta ,\stackrel{~}{e}}\right).$$
Subsubcase $`(\stackrel{~}{\mathrm{\Xi }},\eta ^{})=\left[\begin{array}{c}+\\ +\end{array}\right]`$. We calculate
$$\begin{array}{cc}& \mathrm{\Lambda }(\left[a\right],\stackrel{~}{e},\text{},\left\{\right\})\hfill \\ =& \frac{1}{s!t!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}0\right)}{}\underset{x_1^{}y_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x_1^{})(y_\gamma ,y_1^{})\left(\xi \eta \right)\hfill \\ \stackrel{\text{(}\text{1.4}\text{)}}{=}& \frac{s!}{s!t!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}0\right)}{}\underset{x_1^{}y_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x_1^{})(y_\gamma ,y_1^{})(\overline{\xi }\backslash z,\eta \backslash \{x^{},y^{}\})(x_1^{},x^{})(y_1^{},y^{})\left(\overline{\xi }x_1^{}y_1^{}\right)\hfill \\ =& \frac{1}{t!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}0\right)}{}\underset{x_1^{}y_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma \hfill \\ & (x_\gamma ,x_1^{})(y_\gamma ,y_1^{})(\overline{\xi }\backslash z,\eta \backslash \{x^{},y^{}\})(x_1^{},x^{})(y_1^{},y^{})(x_1^{},x_\gamma )(y_1^{},y_\gamma )\left(\overline{\xi }x_1^{}y_1^{}\right)\hfill \\ =& \frac{1}{t!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}0\right)}{}\underset{x_1^{}y_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (\overline{\xi }\backslash z,\eta \backslash \{x^{},y^{}\})(x_\gamma ,x^{})(y_\gamma ,y^{})\left(\overline{\xi }x_1^{}y_1^{}\right)\hfill \\ =& \left(u1\right)uB_{\left[a\right],\xi ,\eta ,\stackrel{~}{e}}2\left(u1\right)A_{\left[a\right],\xi ,\eta ,\stackrel{~}{e}}.\hfill \end{array}$$
Case $`f_{\stackrel{~}{e}1}^{}`$.
Subcase $`\stackrel{~}{\mathrm{\Xi }}=\left[\begin{array}{c}+\\ \end{array}\right]`$.
Subsubcase $`(\stackrel{~}{\mathrm{\Xi }},\eta ^{})=\left[\begin{array}{c}\\ \end{array}\right]`$. We calculate
$$\begin{array}{cc}& \mathrm{\Lambda }(\left[a\right],\stackrel{~}{e},\text{},\left\{1\right\})\hfill \\ =& \frac{t}{s!t!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}0\right)}{}a^\gamma \epsilon _\gamma (x_\gamma ,x^{})\left(\xi \eta \right)\hfill \\ \stackrel{\text{(}\text{1.4}\text{)}}{=}& \frac{t\left(s+1\right)!}{s!t!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}0\right)}{}a^\gamma \epsilon _\gamma (x_\gamma ,x^{})(\overline{\xi }\backslash z,\eta \backslash \{x^{},y^{}\})(y_\gamma ,y^{})\left(\overline{\xi }y_\gamma \right)\hfill \\ =& \left(s+1\right)\left(B_{\left[a\right],\xi ,\eta ,\stackrel{~}{e}}C_{2,\left[a\right],\xi ,\eta ,\stackrel{~}{e}}\right).\hfill \end{array}$$
Subsubcase $`(\stackrel{~}{\mathrm{\Xi }},\eta ^{})=\left[\begin{array}{c}+\\ \end{array}\right]`$. We calculate
$$\begin{array}{cc}& \mathrm{\Lambda }(\left[a\right],\stackrel{~}{e},\text{},\left\{\right\})\hfill \\ =& \frac{1}{s!t!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}0\right)}{}\underset{x_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x_1^{})\left(\xi \eta \right)\hfill \\ \stackrel{\text{(}\text{1.4}\text{)}}{=}& \frac{s!}{s!t!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}0\right)}{}\underset{x_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x_1^{})(\overline{\xi }\backslash z,\eta \backslash \{x^{},y^{}\})(x_1^{},x^{})(y_\gamma ,y^{})\left(\overline{\xi }x_1^{}y_\gamma \right)\hfill \\ =& \frac{1}{t!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}0\right)}{}\underset{x_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (x_\gamma ,x_1^{})(\overline{\xi }\backslash z,\eta \backslash \{x^{},y^{}\})(x_1^{},x^{})(y_\gamma ,y^{})(x_\gamma ,x_1^{})\left(\overline{\xi }x_1^{}y_\gamma \right)\hfill \\ =& \frac{1}{t!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}0\right)}{}\underset{x_1^{}\overline{\eta }}{}a^\gamma \epsilon _\gamma (\overline{\xi }\backslash z,\eta \backslash \{x^{},y^{}\})(x_\gamma ,x^{})(y_\gamma ,y^{})\left(\overline{\xi }x_1^{}y_\gamma \right)\hfill \\ =& \left(uB_{\left[a\right],\xi ,\eta ,\stackrel{~}{e}}uC_{2,\left[a\right],\xi ,\eta ,\stackrel{~}{e}}A_{\left[a\right],\xi ,\eta ,\stackrel{~}{e}}\right).\hfill \end{array}$$
Subcase $`\stackrel{~}{\mathrm{\Xi }}=\left[\begin{array}{c}\\ +\end{array}\right]`$. Obtained, by symmetry, from subcase $`\stackrel{~}{\mathrm{\Xi }}=\left[\begin{array}{c}+\\ \end{array}\right]`$.
Case $`f_{\stackrel{~}{e}0}^{}`$.
Subcase $`\stackrel{~}{\mathrm{\Xi }}=\left[\begin{array}{c}\\ \end{array}\right]`$.
Subsubcase $`(\stackrel{~}{\mathrm{\Xi }},\eta ^{})=\left[\begin{array}{c}\\ \end{array}\right]`$. We calculate
$$\begin{array}{cc}& \mathrm{\Lambda }(\left[a\right],\stackrel{~}{e},\text{},\left\{\right\})\hfill \\ =& \frac{1}{s!t!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}0\right)}{}a^\gamma \epsilon _\gamma \left(\xi \eta \right)\hfill \\ \stackrel{\text{(}\text{1.4}\text{)}}{=}& \frac{s!}{s!t!}\underset{\gamma \dot{\mathrm{\Gamma }}\left(\stackrel{~}{e}0\right)}{}a^\gamma \epsilon _\gamma (\overline{\xi }\backslash z,\eta \backslash \{x^{},y^{}\})(x_\gamma ,x^{})(y_\gamma ,y^{})\left(\overline{\xi }x_\gamma y_\gamma \right)\hfill \\ =& B_{\left[a\right],\xi ,\eta ,\stackrel{~}{e}}C_{1,\left[a\right],\xi ,\eta ,\stackrel{~}{e}}C_{2,\left[a\right],\xi ,\eta ,\stackrel{~}{e}}.\hfill \end{array}$$
We note that $`su=X_gX_{g+1}`$ and evaluate the linear combination
$$\begin{array}{c}\underset{\beta [0,2]}{}X_{g+1}^{\left(2\beta \right)}G_{\left[a\right],\xi ,\eta }^{}f_{\stackrel{~}{e}\beta }^{}\hfill \\ \begin{array}{ccc}=& 1\hfill & (2(X_gX_{g+1}+2)A_{\left[a\right],\xi ,\eta ,\stackrel{~}{e}}\hfill \\ & & +(X_gX_{g+1}+2)(X_gX_{g+1}+1)B_{\left[a\right],\xi ,\eta ,\stackrel{~}{e}})\hfill \\ +& X_{g+1}\hfill & (2A_{\left[a\right],\xi ,\eta ,\stackrel{~}{e}}+2(X_gX_{g+1}+1)B_{\left[a\right],\xi ,\eta ,\stackrel{~}{e}}\hfill \\ & & (X_gX_{g+1}+1)(C_{1,\left[a\right],\xi ,\eta ,\stackrel{~}{e}}+C_{2,\left[a\right],\xi ,\eta ,\stackrel{~}{e}}))\hfill \\ +& X_{g+1}\left(X_{g+1}+1\right)\hfill & \left(B_{\left[a\right],\xi ,\eta ,\stackrel{~}{e}}\left(C_{1,\left[a\right],\xi ,\eta ,\stackrel{~}{e}}+C_{2,\left[a\right],\xi ,\eta ,\stackrel{~}{e}}\right)\right)\hfill \end{array}\hfill \\ =2\left(X_g+2\right)A_{\left[a\right],\xi ,\eta ,\stackrel{~}{e}}+\left(X_g+2\right)\left(X_g+1\right)B_{\left[a\right],\xi ,\eta ,\stackrel{~}{e}}X_{g+1}\left(X_g+2\right)\left(C_{1,\left[a\right],\xi ,\eta ,\stackrel{~}{e}}+C_{2,\left[a\right],\xi ,\eta ,\stackrel{~}{e}}\right).\hfill \end{array}$$
## Appendix B References
Carter, R.W., Lusztig, G.
* On the Modular Representations of the General Linear and Symmetric Groups, Math. Z. 136, p. 139-242, 1974.
Carter, R.W., Payne, M.T.J.
* On homomorphisms between Weyl modules and Specht modules, Math. Proc. Camb. Phil. Soc. 87, p. 419-425, 1980.
Garnir, H.
* Théorie de la répresentation linéaire des groupes symétriques, Thèse, Mém. Soc. Roy. Sc. Liège, (4), 10, 1950.
James, G.D.
* The Representation Theory of the Symmetric Groups, SLN 682, 1978.
Künzer, M.
* Ties for the $`\text{Z}𝒮_n`$, thesis, http://www.mathematik.uni-bielefeld.de/$``$kuenzer, Bielefeld, 1999.
Specht, W.
* Die irreduziblen Darstellungen der Symmetrischen Gruppe, Math. Z. 39, p. 696-711, 1935.
Invariant address:
Matthias Künzer
Fakultät für Mathematik
Universität Bielefeld
Postfach 100131
D-33501 Bielefeld
kuenzer@mathematik.uni-bielefeld.de
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# Binomial Residues
## 1. Introduction
By a binomial residue we mean a rational function in $`2n`$ variables $`x_1,\mathrm{},x_n`$, $`y_1,\mathrm{},y_n`$, which is defined by a residue integral of the form
(1.1)
$$R_\mathrm{\Gamma }(x,y):=_\mathrm{\Gamma }\frac{t^\gamma }{(x_1+t^{a_1}y_1)^{\beta _1}\mathrm{}(x_n+t^{a_n}y_n)^{\beta _n}}\frac{dt_1}{t_1}\mathrm{}\frac{dt_d}{t_d}.$$
Here $`a_1,a_2,\mathrm{},a_n`$ are non-zero lattice vectors which span $`^d`$, $`\gamma `$ is any vector in $`^d`$, $`\beta _1,\mathrm{},\beta _n`$ are positive integers, and $`\mathrm{\Gamma }`$ ranges over a certain collection, specified in (3.7) below, of compact $`d`$-cycles in the torus $`(^{})^d`$. In this paper we study analytic, combinatorial, and geometric properties of binomial residues. On the analytic side, we view binomial residues as hypergeometric integrals \[18, page 223\] and, consequently, as rational solutions of a certain $`A`$-hypergeometric system of differential equations, in the sense of Gel’fand, Kapranov and Zelevinsky . The $`A`$-hypergeometric system annihilating (1.1) is the left ideal in the $`2n`$-dimensional Weyl algebra generated by the operators
$`_x^u_y^v_x^v_y^u\text{whenever}u,v^n\text{and}{\displaystyle \underset{i=1}{\overset{n}{}}}u_ia_i={\displaystyle \underset{i=1}{\overset{n}{}}}v_ia_i,`$
(1.2) $`x_i_{x_i}+y_i_{y_i}+\beta _i\text{for}i=1,2,\mathrm{},n,\text{and}`$
$`a_{j1}y_1_{y_1}+a_{j2}y_2_{y_2}+\mathrm{}+a_{jn}y_n_{y_n}+\gamma _j\text{for}j=1,2,\mathrm{},d.`$
Here $`_x^u=_{x_1}^{u_1}\mathrm{}_{x_n}^{u_n}`$ for $`u^n`$. In the notation of , this is the system $`H_A(\beta ,\gamma )`$ associated with the $`(n+d)\times 2n`$-matrix
(1.3)
$$A:=\left(\begin{array}{cc}I_n& I_n\\ \mathrm{0\hspace{0.17em}0}\mathrm{}\mathrm{\hspace{0.17em}0}& a_1a_2\mathrm{}a_n\end{array}\right),$$
where $`I_n`$ denotes the $`n\times n`$ identity matrix. The matrix $`A`$ is called the Lawrence lifting of $`a_1,a_2,\mathrm{},a_n`$. Such matrices play an important role in combinatorics \[3, §9.3\] and Gröbner bases \[17, §7, page 55\].
We next introduce a combinatorial invariant associated with a configuration of vectors. For the Lawrence lifting $`A`$, this invariant agrees with that of the submatrix $`M:=(a_1,\mathrm{},a_n)`$. The matroid complex of $`M`$ is the simplicial complex $`\mathrm{\Delta }(M)`$ consisting of all subsets $`I\{1,\mathrm{},n\}`$ such that the corresponding vectors $`a_i,iI`$ are linearly independent. Let $`\chi (M)`$ denote the Euler characteristic of the matroid complex $`\mathrm{\Delta }(M)`$, i.e. the sum of $`(1)^{|I|}`$ for $`I\mathrm{\Delta }(M)`$. The integer $`\chi (A)=\chi (M)`$ equals the Möbius invariant of the dual matroid \[2, Proposition 7.4.7\] and, via Zaslavsky’s Theorem , it counts the regions of the hyperplane arrangement (2.5). Lemma 2.10 implies
(1.4)
$$|\chi (A)|\left(\genfrac{}{}{0pt}{}{n1}{d}\right),$$
with equality if all $`d`$-tuples $`\{a_{i_1},\mathrm{},a_{i_d}\}`$ are linearly independent.
We note that $`\chi (A)=0`$ if and only if $`A`$ has a coloop, i.e., some linear functional on $`^d`$ vanishes on all but one of the points $`a_1,\mathrm{},a_n`$. If this is the case, then every $`A`$-hypergeometric function is a monomial times a solution of a smaller system (1) gotten by contracting the coloops. Thus, we will assume without loss of generality that $`\chi (A)0.`$
A rational function $`f`$ in $`x_1,\mathrm{},x_n`$, $`y_1,\mathrm{},y_n`$ is called unstable if it is annihilated by some iterated derivative $`_x^u_y^v`$. Otherwise we say that $`f`$ is stable. Thus $`f`$ is unstable if it is a linear combination of rational functions that depend polynomially on at least one of the variables. We denote by $`R(\beta ,\gamma )`$ the vector space of rational solutions of $`H_A(\beta ,\gamma )`$, by $`U(\beta ,\gamma )`$ the subspace of unstable rational solutions, and we set
(1.5)
$$𝒮(\beta ,\gamma ):=R(\beta ,\gamma )/U(\beta ,\gamma ).$$
Our main result gives an integral representation for stable rational $`A`$-hypergeometric functions, when $`A`$ is the Lawrence configuration (1.3).
###### Theorem 1.1.
Let $`\beta _{>0}^n`$ and $`\gamma ^d`$. The space $`𝒮(\beta ,\gamma )`$ of stable rational $`A`$-hypergeometric functions of degree $`(\beta ,\gamma )`$ has dimension $`|\chi (A)|`$ and is spanned by binomial residues $`R_\mathrm{\Gamma }(x,y)`$.
We illustrate this theorem with three examples. First consider $`d=1,n=3,a_1=a_2=a_3=1^1`$, $`\beta _1=\beta _2=\beta _3=1`$, and $`\gamma =3`$. The Euler characteristic is $`\chi (A)=2`$. The binomial residues are the integrals
(1.6)
$$_\mathrm{\Gamma }\frac{t^3}{(x_1+ty_1)(x_2+ty_2)(x_3+ty_3)}\frac{dt}{t}.$$
By integrating around the three poles $`t=x_i/y_i`$, we obtain
$$R_1=\frac{x_1^2}{(x_1y_2x_2y_1)(x_1y_3x_3y_1)y_1}$$
$$R_2=\frac{x_2^2}{(x_3y_2x_2y_3)(x_1y_2x_2y_1)y_2}$$
$$R_3=\frac{x_3^2}{(x_2y_3x_3y_2)(x_1y_3x_3y_1)y_3}$$
These residues form a solution basis for the hypergeometric system
$`H_A(\beta ,\gamma )=`$ $`\{_{x_1}_{y_2}_{x_2}_{y_1},_{x_1}_{y_3}_{x_3}_{y_1},_{x_2}_{y_3}_{x_3}_{y_2},`$
$`x_1_{x_1}+y_1_{y_1}+1,x_2_{x_2}+y_2_{y_2}+1,x_3_{x_3}+y_3_{y_3}+1,`$
$`y_1_{y_1}+y_2_{y_2}+y_3_{y_3}+3\}.`$
This is the Aomoto-Gel’fand system for a $`2\times 3`$-matrix, which is holonomic of rank $`3`$; see \[18, §1.5\]. The space $`𝒮(\beta ,\gamma )`$ of rational solutions modulo unstable rational solutions has dimension $`\mathrm{\hspace{0.17em}2}=|\chi (A)|`$, since
$$R_1+R_2+R_3=\frac{1}{y_1y_2y_3}$$
contains no $`x_i`$ and is hence unstable. This identity expresses the fact that the sum of all local residues of a rational $`1`$-form over $`^1`$ is zero.
Our second example is the Lawrence lifting of the twisted cubic curve:
(1.7)
$$d=2,n=4,A=\left(\begin{array}{cccccccc}1& 0& 0& 0& 1& 0& 0& 0\\ 0& 1& 0& 0& 0& 1& 0& 0\\ 0& 0& 1& 0& 0& 0& 1& 0\\ 0& 0& 0& 1& 0& 0& 0& 1\\ 0& 0& 0& 0& 1& 1& 1& 1\\ 0& 0& 0& 0& 0& 1& 2& 3\end{array}\right)$$
Fix $`\beta =(1,1,1,1)`$ and $`\gamma =(1,1)`$. The space of $`A`$-hypergeometric functions is $`10`$-dimensional, and the subspace of rational solutions is $`3`$-dimensional. A basis for $`R(\beta ,\gamma )`$ consists of the three binomial residues
$$R_{23}=\frac{x_2y_3^2y_2}{(x_3^2y_2y_4x_2x_4y_3^2)(x_1x_3y_2^2x_2^2y_1y_3)}$$
$$R_{24}=\frac{y_4y_2(x_2^2y_3y_1+x_1x_3y_2^2)}{(x_3^2y_2y_4x_2x_4y_3^2)(x_2^3y_1^2y_4x_1^2x_4y_2^3)}$$
$$R_{34}=\frac{x_3x_4y_3^3y_4}{(y_2y_4x_3^2y_3^2x_2x_4)(y_3^3x_4^2x_1x_3^3y_1y_4^2)}$$
Other residues can be computed by the Orlik-Solomon relations (cf. §5):
$$R_{14}=R_{24}R_{34},R_{13}=R_{23}+R_{34},R_{12}=R_{23}+R_{24}.$$
For our third example take $`\{a_1,\mathrm{},a_n\}`$ to be the positive roots in the root system of type $`A_d`$. This means $`n=\left(\genfrac{}{}{0pt}{}{d}{2}\right)`$ and (1.1) looks like
$$_\mathrm{\Gamma }\frac{t_1^{\gamma _1}\mathrm{}t_{d1}^{\gamma _{d1}}}{_{1i<jd}(x_{ij}+t_it_j^1y_{ij})^{\beta _{1j}}}\frac{dt_1}{t_1}\mathrm{}\frac{dt_{d1}}{t_{d1}},$$
where $`t_d=1.`$ This is the Selberg type integral studied by Kaneko and many others; see \[18, Example 5.4.7\]. The holonomic rank of the associated $`A`$-hypergeometric system equals $`d^{d2}`$, the number of labeled trees on $`d`$ vertices. The following explicit formula for the number of stable rational hypergeometric functions of Selberg type is given in :
$$|\chi (A_d)|=(d2)\underset{k=0}{\overset{[(d3)/2]}{}}\left(\genfrac{}{}{0pt}{}{d3}{2k}\right)(d1)^{d32k}\underset{i=1}{\overset{k}{}}(2i1).$$
This paper is organized as follows. In §2 we examine hypergeometric Laurent series solutions, and we derive the upper bound in Theorem 1.1. In §3 we establish the connection to toric geometry, by expressing binomial residues as toric residues in the sense of Cox ; see also . Formulas and algorithms for computing binomial residues are presented in §4. In §5, we complete the proof of Theorem 1.1, and we prove Conjecture 5.7 from our previous paper in the Lawrence case.
## 2. Laurent series expansions and Gale duality
In this section we establish the upper bound in Theorem 1.1 for arbitrary rational $`A`$-hypergeometric functions. The Lawrence hypothesis is not needed for this. The main idea is to look at series expansions, which leads to counting cells in a hyperplane arrangement. We fix an arbitrary integer $`r\times s`$-matrix $`A`$ of rank $`r`$ and an integer vector $`\alpha ^r`$.
###### Definition 2.1.
. The $`A`$-hypergeometric system is the left ideal $`H_A(\alpha )`$ in the Weyl algebra $`x_1,\mathrm{},x_s,_1,\mathrm{},_s`$ generated by the toric operators $`^u^v`$, for $`u,v^s`$ such that $`Au=Av`$, and the Euler operators $`_{j=1}^sa_{ij}x_j_j\alpha _i`$ for $`i=1,\mathrm{},r`$. A function $`f(x_1,\mathrm{},x_s)`$, holomorphic on an open set $`U^s`$, is said to be $`A`$-hypergeometric of degree $`\alpha `$ if it is annihilated by the left ideal $`H_A(\alpha )`$.
A rational $`A`$-hypergeometric function admits Laurent series expansions convergent in a suitable open set. In the terminology of these are logarithm-free hypergeometric series with integral exponents. We review their construction and refer to \[18, §3.4\] for proofs and details.
Given a vector $`v^s`$, we define its negative support by
$$\mathrm{nsupp}(v):=\{i\{1,\mathrm{},s\}:v_i_{<0}\}$$
A vector $`v^s`$ is said to have minimal negative support if there is no integer vector $`u`$ in the kernel of $`A`$ such that $`\mathrm{nsupp}(u+v)`$ is properly contained in $`\mathrm{nsupp}(v)`$. The following set of integer vectors,
$$N_v:=\{u\mathrm{ker}_{}(A):\mathrm{nsupp}(u+v)=\mathrm{nsupp}(v)\},$$
is used to define the formal Laurent series
(2.1)
$$\varphi _v(x):=\underset{uN_v}{}\frac{[v]_u_{}}{[v+u]_{u_+}}x^{v+u}$$
$$\text{where}[v]_u_{}=\underset{i:u_i<0}{}\underset{j=1}{\overset{u_i}{}}(v_ij+1)\text{and}[v+u]_{u_+}=\underset{i:u_i>0}{}\underset{j=1}{\overset{u_i}{}}(v_i+j).$$
The Weyl algebra acts on formal Laurent series by multiplication and differentiation. The following is Proposition 3.4.13 in :
###### Proposition 2.2.
Let $`\alpha =Av`$. The series $`\varphi _v(x)`$ is annihilated by $`H_A(\alpha )`$ if and only if the vector $`v^s`$ has minimal negative support.
In order to ensure that the $`A`$-hypergeometric series $`\varphi _v(x)`$ have a common domain of convergence, we fix a generic weight vector $`w^s`$. A vector $`v^s`$ is called an exponent for $`H_A(\alpha )`$ with respect to $`w`$ if $`v`$ has minimal negative support and
(2.2)
$$Av=\alpha \text{and}w,v=\mathrm{min}\{w,u:uv+N_v\}.$$
The following is a restatement of \[18, Theorem 3.4.14, Corollary 3.4.15\]:
###### Theorem 2.3.
The set $`\{\varphi _v:v^s\text{and }v\text{ is an exponent}\}`$ is a basis for the space of hypergeometric functions of degree $`\alpha `$ admitting a Laurent expansion convergent in a certain open subset $`U_w`$ of $`^s`$.
For a more precise description of hypergeometric Laurent series, we next introduce the oriented hyperplane arrangement defined by the Gale dual (or matroid dual) to $`A`$. Set $`m:=sr`$ and let $`B`$ be an integral $`s\times m`$ matrix whose columns are a $``$-basis of $`\mathrm{ker}_{}(A)`$. The matrix $`B`$ has rank $`m`$ and $`AB=0`$. Note that $`B`$ is well-defined modulo right multiplication by elements of $`GL(m,)`$. We identify $`B`$ with its set of row vectors, and we call this configuration the Gale dual of $`A`$:
$$B=\{b_1,\mathrm{},b_s\}^m.$$
Our assumption $`\chi (A)0`$ translates into the condition $`b_j0`$ for all $`j=1,\mathrm{},s`$. As remarked in the Introduction, the study of $`A`$-hypergeometric functions, for arbitrary $`A`$, easily reduces to this case.
Fix an exponent $`v^s`$. We identify the lattice $`^m`$ with the sublattice $`image_{}(B)+v=\mathrm{ker}_{}(A)+v`$ of $`^s`$ via the affine isomorphism $`\lambda B\lambda +v`$. Under this identification, the affine hyperplane
(2.3)
$$\{\lambda ^m:b_j,\lambda =v_j\}$$
corresponds to the coordinate hyperplane $`x_j=0`$ in $`\mathrm{ker}_{}(A)+v^s`$. Let $``$ denote the arrangement in $`^m`$ consisting of the hyperplanes (2.3) for $`j=1,\mathrm{},s`$. We define the negative support of a vector $`\lambda `$ in $`^m`$ as the negative support of its image under the above isomorphism:
$$\mathrm{nsupp}(\lambda ):=\{j\{1,2,\mathrm{},s\}:b_j,\lambda <v_j\}.$$
The set of points with the same negative support will be called a cell of the hyperplane arrangement $``$. Note that our definition of cell differs slightly from the familiar subdivision into relatively open polyhedra by the hyperplanes in $``$. Our cells are unions of these: they are also polyhedra but they are usually not relatively open.
Consider the following attributes of a cell $`\mathrm{\Sigma }`$ in $``$. We say that:
* $`\mathrm{\Sigma }`$ is bounded if $`\mathrm{\Sigma }`$ is a bounded subset of $`^m`$.
* $`\mathrm{\Sigma }`$ is minimal if the set $`\mathrm{\Sigma }^s`$ is nonempty and the support of the elements in this set is minimal with respect to inclusion.
* $`\mathrm{\Sigma }`$ is $`w`$-positive, for a given vector $`w`$ on $`^n`$, if there exists a real number $`\rho `$ such that $`w,\lambda \rho `$ for all $`\lambda \mathrm{\Sigma }`$.
We can now rewrite the hypergeometric series (2.1) as follows:
(2.4)
$$\varphi _\mathrm{\Sigma }:=\varphi _v=\underset{\lambda \mathrm{\Sigma }^m}{}\frac{[v]_{(B\lambda )_{}}}{[v+B\lambda ]_{(B\lambda )_+}}x^{B\lambda +v}$$
If $`\mathrm{\Sigma }`$ is bounded then $`\varphi _\mathrm{\Sigma }`$ is a Laurent polynomial, and if $`\mathrm{\Sigma }`$ is $`w`$-positive then $`\varphi _\mathrm{\Sigma }`$ lies in the Nilsson ring (cf. \[18, §3.4\]) associated with $`w`$, and hence defines an $`A`$-hypergeometric function on $`U_w`$ when $`\mathrm{\Sigma }`$ is minimal.
The following is an immediate consequence of Theorem 2.3:
###### Proposition 2.4.
The series $`\varphi _\mathrm{\Sigma }`$ where $`\mathrm{\Sigma }`$ runs over all $`w`$-positive minimal cells in $``$ form a basis for the space of $`A`$-hypergeometric functions of degree $`\alpha `$ admitting a Laurent expansion convergent in $`U_w^s`$. Restricting to bounded cells $`\mathrm{\Sigma }`$, we get a basis for the subspace of hypergeometric Laurent polynomials.
Recall that a rational function $`f`$ in is called unstable if there exists $`u^s`$ such that the partial derivative $`^u(f)`$ is identically zero.
###### Lemma 2.5.
$`A`$-hypergeometric Laurent polynomials are unstable.
Proof. By Proposition 2.4, it is enough to show that for any bounded minimal chamber $`\mathrm{\Sigma },`$ the common negative support of all monomials in $`\varphi _\mathrm{\Sigma }`$ does not equal $`\{1,\mathrm{},s\}.`$ In fact, suppose
$$\mathrm{\Sigma }=\{\lambda ^m:b_j,\lambda <v_j\text{for all}j=1,2,\mathrm{},m\}.$$
The negative support of any lattice point in $`^m\backslash \mathrm{\Sigma }`$ is a proper subset of $`\{1,2,\mathrm{},n\}`$, and we conclude that $`\mathrm{\Sigma }`$ is not minimal. ∎
If we differentiate an $`A`$-hypergeometric function of degree $`\alpha `$ with respect to $`x_i`$ then we get an $`A`$-hypergeometric function of degree $`\alpha a_i`$. If we iterate this process long enough, for all variables, then only the stable functions survive. The following definition is intended to make this more precise. The Euler-Jacobi cone is the open cone in $`^s`$:
$$\mathrm{Int}(\mathrm{pos}(A))=\{\nu _1a_1+\nu _2a_2+\mathrm{}+\nu _sa_s:\nu _i_{<0}\text{ for all }i\}.$$
Note that $`(\beta ,\gamma )`$ lies in the Euler-Jacobi cone in Example (1.7).
###### Proposition 2.6.
If $`\alpha \mathrm{Int}(\mathrm{pos}(A))`$, then every $`A`$-hypergeometric series of degree $`\alpha `$ is stable.
Proof. If suffices to show that none of the hypergeometric series $`\varphi _v`$ is unstable. Fix a strictly negative vector $`\nu _{<0}^s`$ with $`A\nu =Av=\alpha `$. Let $`k`$ be a positive integer such that $`k\nu ^s`$. For each integer $`\mathrm{}`$, the vector $`v+\mathrm{}(v\nu )`$ has negative support contained in $`\mathrm{nsupp}(v)`$. Since $`v`$ is minimal, we conclude that $`\mathrm{nsupp}(v+\mathrm{}(v\nu ))=\mathrm{nsupp}(v)`$ for all $`\mathrm{}N`$. Let $`I:=\{i\{1,\mathrm{},s\}:v_i0\}.`$ For all $`iI,`$ we have $`v_i>\nu _i`$, and so all the coordinates in $`I`$ of the vectors $`v+\mathrm{}(v\nu )`$ strictly increase with $`\mathrm{}`$. This shows that $`\varphi _v`$ cannot be decomposed as a finite sum of Laurent series that depend polynomially on one variable. ∎
###### Theorem 2.7.
If $`\alpha \mathrm{Int}(\mathrm{pos}(A))`$, then the dimension of the space of $`A`$-hypergeometric Laurent series of degree $`\alpha `$ with a common domain of convergence is bounded above by the Euler characteristic $`|\chi (A)|`$.
Proof. Consider the central hyperplane arrangement gotten from $``$ by translating all $`s`$ hyperplanes so as to pass through the origin. This central arrangement consists of the $`s`$ hyperplanes
(2.5)
$$\{\lambda ^m:b_j,\lambda =\mathrm{\hspace{0.17em}0}\}\text{for}j=1,2,\mathrm{},s.$$
Since $`\alpha `$ is in the Euler-Jacobi cone, the minimal cells $`\mathrm{\Sigma }`$ of $``$ are all unbounded and correspond to certain maximal cones of the central arrangement (2.5). Fix a generic linear functional $`w`$ on $`^m`$. A basis for the relevant space of $`A`$-hypergeometric Laurent series is indexed by the $`w`$-bounded, minimal cells of $``$. Their number is bounded above by the number of $`w`$-bounded maximal cones in the central arrangement.
A classical result in combinatorics due to Zaslavsky states that the number of $`w`$-bounded maximal cones is the absolute value of the Möbius invariant $`\mu (B)`$ of the matroid associated with $`B`$. Our assertion now follows from the following identity from \[2, Proposition 7.4.7 (i)\]:
(2.6)
$$|\mu (B)|=|\chi (A)|.$$
In words, the Möbius invariant of a matroid equals (up to sign) the Euler characteristic of the dual matroid. ∎
###### Corollary 2.8.
For any $`\alpha ^d`$, the complex vector space of rational $`A`$-hypergeometric functions of degree $`\alpha `$ modulo the subspace of unstable functions has dimension at most $`|\chi (A)|`$.
Proof. We represent the rational $`A`$-hypergeometric functions by Laurent series expansions which have a common domain of convergence. Hence it suffices to prove the asserted dimension bound for the space of convergent $`A`$-hypergeometric Laurent series modulo unstable ones.
Choose $`u^s`$ so that $`\alpha Au`$ lies in the Euler-Jacobi cone. The operator $`^u`$ induces a monomorphism from $`𝒮(\alpha )`$ into $`𝒮(\alpha Au)`$. By Proposition 2.6, $`𝒮(\alpha Au)R(\alpha Au)`$, hence the dimension bound follows from Theorem 2.7 applied to $`\alpha Au`$. ∎
Passing from $`\{a_1,\mathrm{},a_n\}`$ to its Lawrence lifting (1.3) corresponds under Gale duality to the operation of replacing $`\{b_1,\mathrm{},b_n\}`$ by its symmetrization $`\{b_1,\mathrm{},b_n,b_1,\mathrm{},b_n\}`$; see \[3, Proposition 9.3.2\]. This process does not change the geometry of the hyperplane arrangement (2.5) and hence it does not change the Möbius invariant $`\mu (B)`$. In view of (2.6), we conclude that the Euler characteristic of $`\{a_1,\mathrm{},a_n\}`$ equals the Euler characteristic of its Lawrence lifting as stated in the Introduction. Corollary 2.8 implies the upper bound in Theorem 1.1.
###### Corollary 2.9.
The space $`𝒮(\beta ,\gamma )`$ has dimension at most $`|\chi (A)|`$
We conclude this section with one more result from matroid theory which we need to complete the proof of Theorem 1.1. A maximally independent subset of $`B`$ is a basis of $`B`$. Note that $`\{b_j:jJ\}`$ is a basis of $`B`$ if and only if $`\{a_j:jJ\}`$ is a basis of $`A`$. A minimally-dependent subset of $`B`$ is a circuit of $`B`$. If $`C=\{b_{i_1},\mathrm{},b_{i_t}\}`$ is a circuit and $`i_1<\mathrm{}<i_t`$ then the set $`C\backslash \{b_{i_t}\}`$ is a broken circuit. A basis of $`B`$ is called an nbc-basis if it contains no broken circuits.
###### Lemma 2.10.
The number of nbc-bases of $`B`$ equals $`|\chi (A)|`$.
Proof. This result follows from (2.6) and Proposition 7.4.5 in . ∎
## 3. Binomial residues and toric geometry
This section is concerned with global residues of meromorphic forms whose polar divisor is a union of hypersurfaces defined by binomials. The analogous case when the polar divisor is defined by linear forms has been extensively studied, for instance, by Varchenko and Brion-Vergne . Our situation can be regarded as a multiplicative analogue to that theory. The binomial hypersurfaces are embedded in a suitable projective toric variety, which places binomial residues into the framework of toric residues . This will allow us in §5 to find bases of $`A`$-hypergeometric stable rational functions for Lawrence liftings in terms of binomial residues, and to give a geometric meaning to the linear dependencies among binomial residues. We refer to for the definition and basic properties of Grothendieck residues.
Let $`X`$ be a complete $`d`$-dimensional toric variety and $`S`$ its homogeneous coordinate ring in the sense of Cox . Homogeneous polynomials in $`S`$ may be thought of as sections of coherent sheaves over $`X`$ and, consequently, their zero-loci are well defined divisors in $`X`$. Let $`T(^{})^d`$ denote the dense torus in $`X`$. Suppose $`G_0,G_1,\mathrm{},G_d`$ are homogeneous polynomials in $`S`$ whose divisors $`D_i`$ satisfy
(3.1)
$$D_0D_1\mathrm{}D_d=\mathrm{}.$$
Any homogeneous polynomial $`H`$ of critical degree determines a meromorphic $`d`$-form on $`X`$ with polar divisor contained in $`D_0\mathrm{}D_d`$,
$$\mathrm{\Phi }(H)=\frac{H\mathrm{\Omega }_X}{G_0G_1\mathrm{}G_d},$$
where $`\mathrm{\Omega }_X`$ is a choice of an Euler form on $`X`$ . The $`d`$-form $`\mathrm{\Phi }(H)`$ defines a Čech cohomology class $`[\mathrm{\Phi }(H)]H^d(X,\widehat{\mathrm{\Omega }}_X^d)`$ relative to the open cover $`\{X\backslash D_i\}_{i=0,\mathrm{},d}`$ of $`X`$. Here $`\widehat{\mathrm{\Omega }}_X^d`$ denotes the sheaf of Zariski $`d`$-forms on $`X`$. The class $`[\mathrm{\Phi }(H)]`$ is alternating with respect to permutations of $`G_0,\mathrm{},G_d`$. If $`H`$ lies in the ideal $`G_0,\mathrm{},G_d`$ of $`S`$ then $`\mathrm{\Phi }(H)`$ is a Čech coboundary. Thus, $`[\mathrm{\Phi }(H)]`$ depends only on the image of the polynomial $`H`$ in the quotient ring $`S/G_0,\mathrm{},G_d`$.
The toric residue $`\mathrm{Res}_G^X(\mathrm{\Phi }(H))`$ is given by the formula
$$\mathrm{Res}_G^X(\mathrm{\Phi }(H))=\mathrm{Tr}_X([\mathrm{\Phi }(H)]),$$
where $`\mathrm{Tr}_X:H^d(X,\widehat{\mathrm{\Omega }}_X^d)`$ is the trace map.
The following proposition can be deduced from Stokes Theorem (cf. , \[19, §7.2\]). It follows directly from the definition of toric residue.
###### Proposition 3.1.
If the polar locus of the $`d`$-form $`\mathrm{\Phi }(H)`$ is contained in the union of only $`d`$ divisors, say $`D_1\mathrm{}D_d`$, then $`\mathrm{Res}_G^X(\mathrm{\Phi }(H))=0`$.
The relationship between toric residues and the usual notion of multidimensional residues is given by the following result.
###### Theorem 3.2.
Let $`G_0,\mathrm{},G_dS`$ satisfy (3.1) and suppose
(3.2)
$$V^0:=D_1\mathrm{}D_dT$$
Then
(3.3)
$$\mathrm{Res}_G^X(\mathrm{\Phi })=\underset{\xi V^0}{}\mathrm{Res}_\xi (\mathrm{\Phi }|_T)$$
where $`\mathrm{Res}_\xi (\mathrm{\Phi }|_T)`$ denotes the (local) Grothendieck residue at $`\xi `$ of the meromorphic form $`\mathrm{\Phi }`$ restricted to the torus and relative to the divisors $`D_1T,\mathrm{},D_dT`$.
Proof. We note, first of all, that (3.1) implies that $`V^0`$ is a finite set and hence the sum in (3.3) makes sense. Moreover, as shown in \[19, §II.7.2\], the local residues in the right-hand side of (3.3) depend only on the divisors $`D_iT`$ and not on the choice of local defining equations.
If $`X`$ is simplicial, then (3.3) is the content of Theorem 0.4 in . For general $`X`$ we argue as in the proof of Theorem 4 in . ∎
We consider now the binomial case which is relevant in this paper. Let $`a_1,\mathrm{},a_n^d`$ as in the Introduction. Let $`\mathrm{\Delta }_i`$ denote the segment $`[0,a_i]^d`$ and $`\mathrm{\Delta }=\mathrm{\Delta }_1+\mathrm{}+\mathrm{\Delta }_n`$ their Minkowski sum. This is a zonotope, that is, a polytope all of whose faces are centrally symmetric \[3, §2.2\]. Let $`\eta _1,\mathrm{},\eta _{2p}`$ denote the inner normals of the facets of the zonotope $`\mathrm{\Delta }`$, where $`\eta _j=\eta _{p+j}`$. We can write
$$\mathrm{\Delta }=\{m^d:m,\eta _j\underset{i:\eta _j,a_i<0}{}\eta _j,a_i;j=1,\mathrm{},2p\}$$
We consider the associated projective toric variety $`X_\mathrm{\Delta }`$. The homogeneous coordinate ring of $`X_\mathrm{\Delta }`$ is the polynomial ring $`S=[z_1,\mathrm{},z_{2p}]`$. The monomials $`t_j:=_{i=1}^{2p}\left(z_i^{\eta _{ij}}\right)`$, for $`j=1,2,\mathrm{},d`$, have degree zero and define coordinates in the torus $`TX_\mathrm{\Delta }`$.
To each binomial $`f_i:=x_i+y_it^{a_i}`$ in the denominator of the kernel of (1.1) we associate the homogeneous polynomial
$$F_i(z):=x_i\underset{\eta _j,a_i<0}{}z_j^{\eta _j,a_i}+y_i\underset{\eta _j,a_i>0}{}z_j^{\eta _j,a_i}.$$
The divisor $`Y_i:=\{F_i(z)=0\}X_\mathrm{\Delta }`$ is the closure of the divisor $`\{f_i(t)=0\}T`$. Moreover, for $`\beta _{>0}^n`$ and $`\gamma ^d`$, the $`d`$-form on $`T`$,
(3.4)
$$\varphi (\beta ,\gamma )=\frac{t^\gamma }{f_1^{\beta _1}\mathrm{}f_n^{\beta _n}}\frac{dt_1}{t_1}\mathrm{}\frac{dt_d}{t_d},$$
extends to the following meromorphic $`d`$-form on the toric variety $`X_\mathrm{\Delta }`$:
(3.5)
$$\mathrm{\Phi }(\beta ,\gamma )=\frac{z^{h(\beta ,\gamma )}}{F_1^{\beta _1}\mathrm{}F_n^{\beta _n}}\mathrm{\Omega }_\mathrm{\Delta },$$
$$\text{where}h_j(\beta ,\gamma )=\eta _j,\gamma \underset{\eta _j,a_i<0}{}\eta _j,\beta _ia_i1,j=1,\mathrm{},2p.$$
The polar divisor of $`\mathrm{\Phi }(\beta ,\gamma )`$ is the union of the divisors $`Y_1,\mathrm{},Y_n`$ and coordinate divisors $`\{z_{\mathrm{}}=0\}`$ for indices $`\mathrm{}`$ with $`h_{\mathrm{}}(\beta ,\gamma )<0`$. For degrees in the Euler-Jacobi cone such indices $`\mathrm{}`$ do not exist. Indeed,
(3.6)
$$\mathrm{Int}(\mathrm{pos}(A))=\{(\beta ,\gamma )^{n+d}:\beta _i>0;h_j(\beta ,\gamma )+1>0\}$$
Thus, if $`(\beta ,\gamma )`$ lies in the Euler-Jacobi cone, the polar divisor of $`\mathrm{\Phi }(\beta ,\gamma )`$ equals $`Y_1\mathrm{}Y_n`$.
We are now prepared to give a precise definition of binomial residues. Fix an index set $`I=\{1i_1<\mathrm{}<i_dn\}`$ such that the corresponding vectors $`a_i`$, $`iI`$, are linearly independent. For $`k=1,\mathrm{},d`$, set $`G_k^I=F_{i_k}`$ and $`D_k=\{G_k^I=0\}`$. For generic values of the coefficients $`x_i,y_i`$, $`iI`$, the divisors $`D_1,\mathrm{},D_d`$ satisfy (3.2).
###### Definition 3.3.
For $`\beta _{>0}^n`$ and $`\gamma ^d`$, let
$$G_0^I=\left(\underset{\mathrm{}:h_{\mathrm{}}(\beta ,\gamma )<0}{}z_{\mathrm{}}\right)\left(\underset{jI}{}F_j\right).$$
Define the following quantity which depends on $`x_1,\mathrm{},x_n,y_1,\mathrm{},y_n`$:
$$R_I(\beta ,\gamma ):=\mathrm{Res}_{G^I}^X(\mathrm{\Phi }(\beta ,\gamma )).$$
Each local residue in the right-hand side of (3.3) may be written as an integral over a $`d`$-cycle “around” the point $`\xi V^0`$. Since for generic values of the coefficients, the map $`f_I=(f_{i_1},\mathrm{},f_{i_d}):T^d`$ is proper, it follows from \[19, §II.8\] that the total sum of residues (3.3) may be written as a single integral,
(3.7)
$$R_I(\beta ,\gamma )=\left(\frac{1}{2\pi i}\right)^d_{\mathrm{\Gamma }(I,x,y)}\frac{t^\gamma }{f_1^{\beta _1}\mathrm{}f_n^{\beta _n}}\frac{dt_1}{t_1}\mathrm{}\frac{dt_d}{t_d},$$
where $`\mathrm{\Gamma }(I,x,y)`$ is the compact real $`d`$-cycle $`\mathrm{\Gamma }(I,x,y)T`$ defined by $`\{|f_{i_1}|=\epsilon _1,\mathrm{},|f_{i_d}|=\epsilon _d\}`$ for small positive $`\epsilon _1,\mathrm{},\epsilon _d.`$ Moreover, the cycle $`\mathrm{\Gamma }(I,x,y)`$ can be locally replaced by a cohomologous cycle $`\mathrm{\Gamma }(I)`$ independent of $`(x_1,\mathrm{},x_n,y_1,\mathrm{},y_n)`$. See \[18, §5.4\] for further details.
We close this section with the observation that the “basic binomial residue” $`R_I(\beta ,\gamma )`$ is indeed a rational $`A`$-hypergeometric function.
###### Lemma 3.4.
The toric residue $`R_I(\beta ,\gamma )`$ is a rational function of $`(x,y)`$ and is annihilated by the hypergeometric system (1).
Proof. For any choice of polynomials $`G_0,\mathrm{},G_d`$, the trace map $`\mathrm{Tr}_X`$ in the definition of the toric residue has its image in the subfield of $``$ generated by the coefficients of the $`G_i`$. This implies that $`R_I(\beta ,\gamma )`$ is an element in the rational function field $`(x_1,\mathrm{},x_n,y_1,\mathrm{},y_n)`$.
The kernel of the integral (3.4) is annihilated by the toric operators $`_x^u_y^v_x^v_y^u`$ in (1). Hence so is the integral itself, by diffentiating under the integral sign. Specifically, it follows from \[6, Lemma 6\] that
(3.8)
$$_{x_i}R_I(\beta ,\gamma )=\beta _iR_I(\beta +e_i,\gamma ),\text{and}$$
(3.9)
$$_{y_i}R_I(\beta ,\gamma )=\beta _iR_I(\beta +e_i,\gamma +a_i),$$
where $`e_1,\mathrm{},e_d`$ is the standard basis of $`^d`$. The verification of the homogeneity equations is immediate from the expression (3.4) for the form $`\varphi (\beta ,\gamma )`$. Hence $`R_I(\beta ,\gamma )`$ is a rational solution of $`H_A(\beta ,\gamma )`$. ∎
## 4. Computing binomial residues
In this section we present methods for computing the binomial residue $`R_I(\beta ,\gamma )`$. Here $`I=\{i_1,\mathrm{},i_d\}`$ is a fixed column basis of the matrix $`M=(a_1,\mathrm{},a_n)`$. Let $`M_I`$ denote the non-singular $`d\times d`$ matrix with columns $`a_i`$, $`iI`$. Write $`M_I^1=(\mu _{ij})GL(d,)`$. We set $`V_I=\{\xi T:f_i(\xi )=0`$ for all $`iI\}`$. The points in $`V_I`$ are in bijection with the characters $`\theta \mathrm{Hom}(^d,^{})`$ satisfying $`\theta (a_i)=1`$, for all $`iI`$. The point $`\xi ^\theta =(\xi _1^\theta ,\mathrm{},\xi _d^\theta )V_I`$ indexed by $`\theta `$ has coordinates
$$\xi _j^\theta =\theta (e_j)\underset{iI}{}\left(\frac{x_i}{y_i}\right)^{\mu _{ij}}$$
There are $`det(M_I)`$-many simple roots $`\xi ^\theta `$ provided all $`x_i,y_i`$ are nonzero.
Let $`g`$ be a function meromorphic on the torus $`T=(^{})^d`$ and regular at a simple root $`\xi V_I`$. Then the local Grothendieck residue of the meromorphic $`d`$-form $`\frac{g}{f_{i_1}\mathrm{}f_{i_d}}\frac{dt_1}{t_1}\mathrm{}\frac{dt_d}{t_d}`$ at the point $`\xi `$ equals
(4.1)
$$R_{I,\xi }[g]=\frac{g(\xi )}{J_I(\xi )}$$
where $`J_I`$ denotes the toric Jacobian of the binomials $`f_i=x_i+y_it^{a_i}`$:
$$J_I(t)=det\left(t_j\frac{f_i}{t_j}\right)_{iI}^{j=1,\mathrm{},d}=detM_I(\underset{iI}{}y_i)t^{a_I}.$$
Here $`a_I=a_{i_1}+\mathrm{}+a_{i_d}`$. We deduce the following identity
(4.2)
$$J_I(\xi )=(1)^ddetM_I(\underset{iI}{}x_i)\text{for all}\xi V_I.$$
We obtain the following procedure for summing (4.1) over all $`\xi V_I`$.
###### Algorithm 4.1.
(Computing global residues using Gröbner bases)
Input: A $`d\times d`$-integer matrix $`M_I`$ of rank $`d`$, a Laurent polynomial $`g(t)`$.
Output: The global residue
$$R_I[g]:=\underset{\xi V_I}{}R_{I,\xi }[g]$$
* Fix the field $`K=(x_1,\mathrm{},x_n,y_1,\mathrm{},y_n)`$ and write the Laurent polynomial ring over $`K`$ as a quotient of a polynomial ring:
$$K[t_1,\mathrm{},t_d,t_1^1,\mathrm{},t_d^1]=K[t_0,t_1,\mathrm{},t_d]/t_0t_1\mathrm{}t_d1.$$
* Compute any Gröbner basis $`G`$ for its ideal $`f_{i_1},\mathrm{},f_{i_d}`$.
* Let $`B`$ be the set of standard monomials for $`G`$ in $`K[t_0,\mathrm{},t_d]`$
* Compute the trace of $`g`$ modulo $`B`$ as follows:
$$\underset{\xi V_I}{}g(\xi )=\underset{t^bB}{}\mathrm{coeff}_{t^b}\left(\mathrm{normalform}_G(t^bg(t))\right)$$
* Output the result of step (4) divided by the monomial in (4.2).
The output produced by the above algorithm is a rational function in $`x_i,y_i`$ and the coefficients of $`g`$. In the case when $`g`$ is a Laurent monomial, one can give a completely explicit formula for that output.
###### Lemma 4.2.
Let $`\gamma ^d`$. If $`\nu =M_I^1\gamma `$ lies in the lattice $`^d`$ then
(4.3)
$$R_I[t^\gamma ](x,y)=\frac{(1)^{|\nu |+d}}{det(M_I)}\underset{iI}{}x_i^{\nu _i1}y_i^{\nu _i}.$$
Otherwise the global residue $`R_I[t^\gamma ]`$ is zero.
Proof. It follows from (4.1) and (4.2) that
$$R_{I,\xi ^\theta }[t^\gamma ](x,y)=\theta (\gamma )\frac{(1)^d}{det(M_I)}\underset{iI}{}x_i^{\nu _i1}y_i^{\nu _i}$$
where $`\nu _i:=_{j=1}^d\mu _{ij}\gamma _j`$, $`iI`$. Thus, the global residue is given by
$$R_I[t^\gamma ](x,y)=\left(\underset{\theta }{}\theta (\gamma )\right)\frac{(1)^d}{det(M_I)}\underset{iI}{}x_i^{\nu _i1}y_i^{\nu _i}$$
and consequently it vanishes unless $`\gamma M_I^d`$. In this case we have (4.3) for $`\gamma =_{iI}\nu _ia_i`$, and $`|\nu |:=_{iI}\nu _i`$. ∎
We now compute the binomial residue $`R_I(\beta ,\gamma )`$ for $`I=\{i_1,\mathrm{},i_d\}`$ as above. In view of (3.8) and (3.9), it suffices to consider the case $`\beta =\mathrm{𝟏}:=(1,\mathrm{},1)`$. Set $`J:=\{1,\mathrm{},n\}\backslash I`$ and let $`M_J`$ denote the matrix whose columns are the vectors $`a_j,jJ`$. Since the coefficients are generic, none of the polynomials $`f_j`$, $`jJ`$ vanishes on any point of $`V_I`$ and hence
(4.4)
$$R_I(\mathrm{𝟏},\gamma )=R_I[t^\gamma /f_J(t)](x,y)\text{where}f_J(t)=\underset{jJ}{}f_j(t).$$
This gives rise to the following symbolic algorithm for binomial residues.
###### Algorithm 4.3.
(Computing binomial residues)
Input: Vectors $`a_1,\mathrm{},a_n`$ and $`\gamma `$ as above, and a basis $`I=\{i_1,\mathrm{},i_d\}`$.
Output: The rational function $`R_I(\mathrm{𝟏},\gamma )`$ of $`x_1,\mathrm{},x_n,y_1,\mathrm{},y_n`$.
* Run steps (1), (2) and (3) of Algorithm 4.1.
* Using linear algebra over the field $`K`$, compute the unique polynomial $`g(t)=_{t^bB}c_bt^b`$ such that all $`c_b`$ lie in $`K`$ and $`g(t)f_J(t)t^\gamma `$ reduces to zero modulo the Gröbner basis $`G`$.
* Run steps (4) and (5) of Algorithm 4.1.
The output of this algorithm is an element of the field $`K`$. It is nonzero and has the following expansion as a Laurent series in $`x_i,y_i`$.
###### Proposition 4.4.
Suppose $`\gamma M^n`$. Then $`R_I(\mathrm{𝟏},\gamma )0`$ and
(4.5)
$$R_I(\mathrm{𝟏},\gamma )=\frac{1}{det(M_I)}(1)^{d+|\nu |+|\mu |}\underset{iI}{}\frac{x_i^{\nu _i1}}{y_i^{\nu _i}}\underset{jJ}{}\frac{y_j^{\mu _j}}{x_j^{\mu _j+1}},$$
where the sum is over $`\nu ^I`$ and $`\mu ^J`$ such that $`M_I\nu M_J\mu =\gamma `$. Moreover, for every $`\beta _{>0}^n`$, the residue $`R_I(\beta ,\gamma )`$ is a stable rational hypergeometric function.
Proof. We expand
(4.6)
$$t^\gamma f_J(t)^1=\underset{\mu ^J}{}\underset{jJ}{}\left(y_j^{\mu _j}x_j^{\mu _j1}\right)t^{\gamma +M_J\mu }.$$
Applying (4.3) to each term of (4.6) yields the Laurent expansion (4.5).
Suppose now that $`\gamma =M_I\nu _0M_J\mu _0`$, $`\nu _0^I`$, $`\mu _0^J`$. There exists a vector $`m_{>0}^J`$ such that $`m_ja_jM_I^d`$. Hence for $`k`$,
$$\gamma +M_J(\mu _0+km)M_I^d$$
and $`\mu _0+km`$ is non-negative for $`k0`$. Hence, the series (4.5) contains infinitely many non-zero terms. This shows that $`R_I(\mathrm{𝟏},\gamma )0`$.
Suppose now that $`\beta _{>0}^n`$ is arbitrary. In view of (3.8), it suffices to show that the derivative $`_x^{\beta \mathrm{𝟏}}`$ of the series (4.5) contains infinitely many powers of each of the variables $`x_{\mathrm{}}`$, $`\mathrm{}=1,\mathrm{},n`$. The previous argument shows that this is indeed the case for $`x_j`$, $`jJ`$ and also for a variable $`x_{i_0}`$, $`i_0I`$, unless every vector $`a_j`$, $`jJ`$, is in the $``$-span of $`\{a_i,iI,ii_0\}`$. But this would mean that the points $`a_k`$, $`ki_0`$ would define a coloop in $`A`$ which is impossible by assumption. ∎
Our final task in this section is to identify the irreducible factors in the denominators of these binomial residues. Let $`C\{1,\mathrm{},n\}`$ be a circuit, i.e., the set $`\{a_i,iC\}`$ obeys a unique (up to sign) linear relation $`_{iC}m_ia_i=0`$ over $``$ such that $`\mathrm{gcd}(m_i,iC)=1`$. Then
$$\mathrm{Res}(C;x,y)=\underset{m_i>0}{}x_i^{m_i}\underset{m_j<0}{}y_j^{m_j}(1)^{|C|}\underset{m_i>0}{}y_i^{m_i}\underset{m_j<0}{}x_j^{m_j}$$
is the resultant of the binomials $`f_i,iC`$. In fact, the singular locus of $`H_A(\beta ,\gamma )`$ is described by the product of all the variables and all the resultants $`\mathrm{Res}(C;x,y)`$ as $`C`$ ranges over the circuits (cf. ,). Let $`I`$ be a basis as above. Note that for each $`jI`$, there exists a unique subset $`I^{}(j)I`$, such that $`I(j):=I^{}(j)\{j\}`$ is a circuit.
###### Theorem 4.5.
The binomial residue, defined by $`I,\beta ,\gamma `$ as above, equals
(4.7)
$$R_I(\beta ,\gamma )=\frac{P(x,y)}{x^ay^b_{jI}\mathrm{Res}(I(j);x,y)^{c_j}}\text{with all }c_j>0$$
where $`P(x,y)`$ is a polynomial relatively prime from the denominator.
Proof. We may assume that $`\beta =\mathrm{𝟏}`$. It follows from a variant of Theorem 1.4 in that $`R_I(\mathrm{𝟏},\gamma )`$ is a rational function whose denominator divides a monomial times
$$\underset{jI}{}\mathrm{Res}(f_{i_1},\mathrm{},f_{i_d},f_j).$$
Since $`\{a_k|kI(j)\}`$ is the unique essential subset of $`\{a_i|iI\{j\}\}`$, with “essential” as defined in , we have that
$$\mathrm{Res}(f_{i_1},\mathrm{},f_{i_d},f_j)=\mathrm{Res}(I(j);x,y).$$
We know by Proposition 4.4 that $`P`$ is non zero. Moreover, if any of the factors $`\mathrm{Res}(I(j);x,y)`$ were missing from the denominator of $`R_I(\mathrm{𝟏},\gamma )`$, then the Laurent series (4.5) would contain only finitely many powers of $`x_j`$. The formula in Proposition 4.4 implies that is impossible. ∎
For unimodular bases, Theorem 4.5 can be refined as follows:
###### Proposition 4.6.
Suppose that $`\{a_i|iI\}`$ is a $``$-basis of $`^d`$. Then
(4.8)
$$R_I(\mathrm{𝟏},\gamma )=\frac{x^ay^b}{_{jI}\mathrm{Res}(I(j);x,y)}$$
where $`x^a`$ and $`y^b`$ are monomials specified in the proof.
Proof. Choose $`\nu ,n_j^I`$, $`jJ`$, so that $`\gamma =M_I\nu `$, $`a_j=M_In_j`$, Then
$$\gamma +\underset{jJ}{}\mu _ja_j=M_I(\nu +\underset{jJ}{}\mu _jm_j)\text{ for all }\mu ^J,$$
and consequently, the Laurent series (4.5) reduces, up to sign, to
$`R_I(\mathrm{𝟏},\gamma )`$ $`=`$ $`{\displaystyle \frac{x_I^{\nu \mathrm{𝟏}}}{y_I^\nu x_J}}{\displaystyle \underset{\mu ^J}{}}{\displaystyle \underset{jJ}{}}{\displaystyle \underset{iI}{}}x_i^{n_{ij}\mu _j}y_i^{n_{ij}\mu _j}y_j^{\mu _j}x_j^{\mu _j}`$
$`=`$ $`{\displaystyle \frac{x_I^{\nu \mathrm{𝟏}}}{y_I^\nu }}{\displaystyle \underset{n_{ij}>0}{}}y_i^{n_{ij}}{\displaystyle \underset{n_{ij}<0}{}}x_i^{n_{ij}}{\displaystyle \underset{jJ}{}}\mathrm{Res}(I(j);x,y)^1\text{}`$
## 5. The lower bound and the linear relations
In this section we establish the lower bound in Theorem 1.1 by exhibiting $`|\chi (A)|`$ many linearly independent binomial residues $`R_I(\beta ,\gamma )`$ for fixed $`\beta ,\gamma `$ and fixed Lawrence matrix
$$A:=\left(\begin{array}{cc}I_n& I_n\\ \mathrm{0\hspace{0.17em}0}\mathrm{}\mathrm{\hspace{0.17em}0}& a_1a_2\mathrm{}a_n\end{array}\right).$$
We will show that all linear relations among the $`R_I(\beta ,\gamma )`$ arise from Proposition 3.1 and correspond to Orlik-Solomon relations \[16, §3.1\].
The Gale dual to the Lawrence matrix $`A`$ has the form
(5.1)
$$B=\{b_1,\mathrm{},b_n,b_1,\mathrm{},b_n\},$$
where $`B_0=\{b_1,\mathrm{},b_n\}^{nd}`$ is a Gale dual of $`\{a_1,\mathrm{},a_n\}`$. According to Corollary 2.8 and Lemma 2.10, the dimension of the space of stable rational $`A`$-hypergeometric functions of degree $`(\beta ,\gamma )`$ is at most the number of nbc-bases in $`B`$, which agrees with the number of nbc-bases in $`B_0`$. The following converse will imply Theorem 1.1.
###### Theorem 5.1.
Let $`\beta _{>0}^n`$ and $`\gamma ^d`$. Then the set of binomial residues $`R_I(\beta ,\gamma )`$, where $`\{1,\mathrm{},n\}\backslash I`$ runs over all nbc-bases of $`B_0`$, is linearly independent modulo the space of unstable rational functions.
It is convenient to use the following characterization for being an nbc-basis of the dual matroid. The proof of Lemma 5.2 is straightforward.
###### Lemma 5.2.
The set $`\{1,\mathrm{},n\}\backslash I`$ is an nbc-basis of $`B_0=\{b_1,\mathrm{},b_n\}`$ if and only if, for each $`i_0I`$, there exists $`j_0\{1,\mathrm{},n\}\backslash I`$ such that $`j_0>i_0`$ and $`I\backslash \{i_0\}\{j_0\}`$ is a basis of $`\{a_1,\mathrm{},a_n\}^d`$.
Proof of Theorem 5.1 Consider the space $`𝒮(\beta ,\gamma )`$ of stable rational hypergeometric functions defined in the Introduction. The derivative $`_{x_i}`$ induces a monomorphism from $`𝒮(\beta ,\gamma )`$ into $`𝒮(\beta +e_i,\gamma )`$, while $`_{y_i}`$ induces an monomorphism into $`𝒮(\beta +e_i,\gamma +a_i)`$. Binomial residues are mapped to binomial residues, with the set of irreducible factors in their denominators preserved. We may thus assume $`\beta =\mathrm{𝟏}`$. All linear spaces in this proof are understood modulo unstable rational functions.
By Theorem 4.5, for any basis $`I`$ of $`\{a_1,\mathrm{},a_n\}`$, the denominator of $`R_I(\mathrm{𝟏},\gamma )`$ equals a monomial multiplied by
(5.2)
$$\underset{jI}{}\mathrm{Res}(I(j);x,y)$$
Let $`_0`$ denote the set of indices $`I`$ complementary to nbc-bases of $`B_0`$. Let $`R__0`$ denote the linear span of binomial residues $`R_I(\mathrm{𝟏},\gamma )`$, $`I_0`$. Clearly, $`nI`$ for any $`I_0`$. Our goal is to show $`dim_{}(R__0)=\mathrm{\#}_0`$.
Let $`K`$ be a circuit of $`\{a_1,\mathrm{},a_n\}`$ which contains the index $`n`$. Define $`R__0(K)`$ to be the span of all binomial residues $`R_I(\mathrm{𝟏},\gamma )`$ with $`I_0`$ and $`I(n)=K`$, i.e., $`K`$ is the unique circuit in $`I\{n\}`$. We may decompose
(5.3)
$$R__0=\underset{K}{}R__0(K)$$
The sum in (5.3) is direct because no element in $`_{K^{}K}R__0(K^{})`$ contains $`\mathrm{Res}(K;x,y)`$ in its denominator, while all elements in $`R__0(K)`$ do.
Thus, it suffices to fix $`K=K_0`$ and show that the binomial residues $`R_I(\mathrm{𝟏},\gamma )`$ with $`I_0`$ and $`I(n)=K_0`$ are linearly independent. Let
$$_1=\{I_0:I(n)=K_0\}.$$
Let $`n_1`$ denote the largest index which does not belong to $`K_0`$, then note that $`n_1I`$ for any $`I_1`$. Indeed, if $`n_1I`$, $`I_1`$, then we would not be able to replace $`a_{n_1}`$ by $`a_j`$ with $`j>n_1`$ and still have a basis; this would contradict Lemma 5.2. This means that we can repeat the previous argument with $`_1`$ in place of $`_0`$ and $`n_1`$ in place of $`n`$ and obtain a decomposition of $`R__1`$ as a direct sum of subspaces $`R__1(K)`$ spanned by binomial residues $`R_I(\mathrm{𝟏},\gamma )`$ with $`I_1`$ and $`I(n_1)=K`$. Continuing in this manner, all subspaces $`R__p(K)`$ will eventually be one-dimensional. Then, the desired result follows from Proposition 4.4. ∎
We next describe all linear relations among the binomial residues $`R_I(\beta ,\gamma )`$ as $`I`$ varies. In the identity below, it is essential to keep track of signs. Namely, if $`I^{}`$ is taken to be ordered then we must multiply $`R_I^{}\mathrm{}(\beta ,\gamma )`$ by the sign of the permutation which orders $`I^{}\{\mathrm{}\}`$.
###### Theorem 5.3.
Let $`I^{}`$ be a $`(d1)`$-subset of $`\{1,\mathrm{},n\}`$ and $`\mathrm{ind}I^{}`$ the set of indices $`\mathrm{}`$ such that $`\{a_{\mathrm{}}\}\{a_i:iI^{}\}`$ is a basis of $`^d`$. Then
$$\underset{\mathrm{}\mathrm{ind}I^{}}{}R_I^{}\mathrm{}(\beta ,\gamma )\mathrm{\hspace{0.17em}\hspace{0.17em}0}\text{modulo unstable rational functions},$$
and these span all the $``$-linear relations relations among the $`R_I(\beta ,\gamma )`$.
Proof. By Proposition 4.4, all $`R_I(\beta ,\gamma )`$ residues are stable. We have established that the spaces $`𝒮(\beta ,\gamma )`$ have the same dimension $`|\chi (A)|`$ for all $`\beta `$ and $`\gamma `$. It follows that the maps $`_{x_i}:𝒮(\beta ,\gamma )𝒮(\beta +e_i,\gamma )`$ and $`_{y_i}:𝒮(\beta ,\gamma )𝒮(\beta +e_i,\gamma +a_i)`$ are isomorphisms. Iterating, we can assume that $`(\beta ,\gamma )`$ lies in the Euler-Jacobi cone $`\mathrm{Int}(\mathrm{pos}(A))`$. By Proposition 2.6, there are no unstable rational $`A`$-hypergeometric functions, so we are claiming that $`_{\mathrm{}\mathrm{ind}I^{}}R_I^{}\mathrm{}(\beta ,\gamma )`$ is zero.
We may assume that $`\{a_i:iI^{}\}`$ is linearly independent. On the $`B`$-side, the complement of $`I^{}`$ has $`nd+1`$ elements and therefore defines a dependent set $`\{b_i,iI^{}\}`$. We can consider as in §2, the central hyperplane arrangement $`𝒜`$ defined by $``$. Consider the socle of the Orlik-Solomon algebra of that hyperplane arrangement \[16, §3.1\]. The linear relation in Theorem 5.3 is the translation to the $`A`$-side of the relation in the socle degree of the Orlik-Solomon algebra defined by $`\{b_i,iI^{}\}`$. In view of \[16, Theorem 3.4\] and Theorem 1.1, it suffices to show that the asserted relations are valid. It will then follow by dimension reasons that they span all $``$-linear relations.
We now prove the identity $`_{\mathrm{}\mathrm{ind}I^{}}R_I^{}\mathrm{}(\beta ,\gamma )=\mathrm{\hspace{0.17em}0}`$ using the formulation in terms of toric residues given in §2. By (3.6), all $`h_j(\beta ,\gamma )`$ are non negative, and so the polar divisor of the form $`\mathrm{\Phi }(\beta ,\gamma )`$ in (3.5) is contained in the union of the divisors $`Y_i=\{F_i=0\},i=1,\mathrm{},n.`$
For $`k=1,\mathrm{},d1`$, set $`G_k^I^{}=F_{i_k}`$. Set also $`G_d^I^{}=_{jI^{}}F_j`$ and let $`G_0^I^{}=z_1\mathrm{}z_{2p}.`$ Then, $`G_0^I^{},\mathrm{},G_d^I^{}`$ define divisors with empty intersection in $`X=X_\mathrm{\Delta }`$ for generic values of the coefficients and moreover
$$\mathrm{\Phi }(\beta ,\gamma )=\frac{z^{h(\beta ,\gamma )}\mathrm{\Omega }_\mathrm{\Delta }}{G_1\mathrm{}G_d}.$$
Proposition 3.1 implies that the corresponding toric residue vanishes:
$$\mathrm{Res}_{G^I^{}}^X(\mathrm{\Phi }(\beta ,\gamma ))=0.$$
On the other hand, consider also the following $`nd+1`$ families of divisors: for any $`\mathrm{}I^{}`$, set $`G_k^{I^{},\mathrm{}}=G_k^I^{}`$ for any $`k=1,\mathrm{},d1,`$ $`G_d^{I^{},\mathrm{}}=F_{\mathrm{}}`$ and $`G_0^{I^{},\mathrm{}}=_{jI^{}\{\mathrm{}\}}F_j.`$ Again, these divisors have empty intersection on $`X`$ for generic values of the coefficients and the poles of $`\mathrm{\Phi }(\beta ,\gamma )`$ are contained in their union, and so we can consider the toric residues $`\mathrm{Res}_{G^{I^{},\mathrm{}}}^X(\mathrm{\Phi }(\beta ,\gamma )).`$ These toric residues are non-zero precisely when $`\mathrm{}\mathrm{ind}I^{}`$. We conclude that the following relations hold:
$`{\displaystyle \underset{\mathrm{}\mathrm{ind}I^{}}{}}\mathrm{Res}_{G^{I^{},\mathrm{}}}^X(\mathrm{\Phi }(\beta ,\gamma ))={\displaystyle \underset{\mathrm{}j}{}}\mathrm{Res}_{G^{I^{},\mathrm{}}}^X(\mathrm{\Phi }(\beta ,\gamma ))=\mathrm{Res}_{G^I^{}}^X(\mathrm{\Phi }(\beta ,\gamma ))=\mathrm{\hspace{0.17em}0}.`$
The second equality follows from a variation on \[19, §II.7\]. Translating back to binomial residues completes the proof of Theorem 5.3. ∎
In , we studied the problem of classifying vector configurations $`A`$ for which there exist rational a $`A`$-hypergeometric function which is not a Laurent polynomial. We conjectured \[8, Conjecture 1.3\] that such a configuration has to have a facial subset which is an essential Cayley configuration. It is easy to see that Lawrence liftings are Cayley configurations of segments; they are essential if and only if $`n=d+1`$. We also conjectured \[8, Conjecture 5.7\] that a rational $`A`$-hypergeometric function has an iterated derivative which is a linear combination of toric residues associated with facial subsets of $`A`$.
###### Theorem 5.4.
Conjecture 5.7 in holds for Lawrence configurations.
Proof. Let $`A`$ be a Lawrence configuration. The assertion of \[8, Conjecture 5.7\] is obvious for unstable rational hypergeometric functions. On the other hand, given a stable rational hypergeometric function, a suitable derivative will have degree in the Euler-Jacobi cone and hence, by Theorem 1.1, will be a linear combination of toric residues. ∎
Acknowledgements: Alicia Dickenstein was partially supported by UBACYT TX94 and CONICET, Argentina, and the Wenner-Gren Foundation, Sweden. Bernd Sturmfels was partially supported by NSF Grant DMS-9970254.
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# Steady-states and kinetics of ordering in bus-route models: connection with the Nagel-Schreckenberg model
## I Introduction
Systems of interacting particles driven far from equilibrium are of current interest in statistical physics . Microscopic models of such systems often capture some aspects of vehicular traffic. In such ”particle-hopping” models of vehicular traffic the particles represent vehicles and the nature of the interactions among these particles is determined by the manner in which the vehicles influence the motion of each other . The dynamics of these models are often formulated in terms of ”update rules” using the language of cellular automata (CA) . For example, the Nagel-Schreckenberg (NaSch) model is the most popular minimal CA model of vehicular traffic on highways while, to our knowledge, the first CA model of city traffic was developed by Biham, Middleton and Levin . The results obtained for these models, using the techniques of statistical mechanics, are not only of fundamental intetest for understanding truly nonequilibrium phenomena but may also find practical use in traffic science and engineering. . Among such results are the time-headway and distance-headway distributions. The time-headway(TH) is defined as the time interval between the departures (or arrivals) of two successive vehicles recorded by a detector placed at a fixed position on the route while the distance between the successive vehicles can be defined as the corresponding distance-headway (DH). The distributions of TH and DH not only contain detailed informations on the nature of the spatio-temporal organization of the vehicles but are also of practical interest to traffic engineers because larger headways provide greater margins of safety whereas higher capacities of the highway require smaller headways.
In a $`1998`$ paper, O’Loan et al have developed a one- dimensional lattice model of bus-route where the buses are represented by particles which move from one site to the next; each site of this model represents a bus stop along the route. The motion of the buses in this bus route model with random sequential updating (BRMRSU) is strongly influenced by the passengers waiting at the bus stops. The BRMRSU model may be viewed as a generalization of a simple particle-hopping model, namely, the totally asymmetric simple exclusion process (TASEP) by coupling the dynamics of the particles to another new variable which represents the presence (or absence) of passengers waiting at the bus stops. The bus route model in does not deal with overcrowded buses; it implicitly assumes that either the buses have infinite capacity or that the passenger arrival rate is slow enough to avoid overcrowding.
The BRMRSU exhibits a Bose-Einstein-condensation-like phenomenon which has been observed earlier in the TASEP and in the NaSch model when quenched random hopping rates are associated with the particles . However, unlike the stable Bose-Einstein-condensed states observed at sufficiently low densities in the TASEP (and in the NaSch model) with random hopping rates, those in the BRMRSU are metastable. The main characteristic of the spatially-inhomogeneous Bose-Einstein-condensed state is the existence of a macroscopically long gap in front of a cluster of vehicles led by the slowest one. In finite systems, for small $`\lambda `$ ($`\lambda `$ is the rate of passenger arrival at a bus stop), the bus clusters (or, equivalently, the gaps between clusters) in the BRMRSU exhibit interesting coarsening phenomena as the system evolves from a random initial state. O’Loan et al. find that, after sufficiently long time, the typical size of the large gaps in the system grows with time $`t`$ according to a power growth law $`t^{1/2}`$.
In this paper we use a BRM with parallel updating (BRMPU) which is obtained from the BRMRSU by replacing the random sequential updating rule with parallel updating, with the aim of relating it with the NaSch model where updating is done in parallel. We also propose here two extensions of the NaSch model (from now onwards referred to as the models Y and Z) by replacing the constant hopping rates with two different time-/space-dependent hopping rates which we shall specify explicitly in section II. We have computed the TH distributions in the steady-states of the BRMPU as well as in the models Y and Z through computer simulations. Comparison of these distributions are shown in section III. Such comparisons elucidate the connection between the BRMPU and the NaSch model. Then, approximating the BRMPU as a generalization of the NaSch model with a time-dependent hopping rate for the buses, we calculate in section IV the TH distribution in the BRMPU from the corresponding analytical expression in the NaSch model. We compare the TH distributions thus derived from analytical considerations with the corresponding results of computer simulations of the BRMPU. These comparisions do not merely point out the regimes of validity of our analytical results but also indicate the differences arising from the different natures of the low-density steady-states in the BRMPU and the NaSch model. Finally, we investigate in section V, interesting kinetic phenomena at low densities of the BRMPU by computing the appropriate correlation functions (to be defined in section V). We extract the universal laws governing the growth of the clusters of buses in finite samples of BRMPU at low densities where the system approaches a Bose-Einstein-like ”condensed” state evolving from random initial states.
## II The models and methods
Let us first summarize how the totally asymmetric exclusion process (TASEP) , the NaSch model and the bus route models are defined.
### A TASEP and the NaSch model
In the ”particle-hopping” models of traffic the position, speed, acceleration as well as time are treated as discrete variables. In this approach, a lane is represented by a one-dimensional lattice. Each lattice site represents a ”cell” which can be either empty or occupied by at most one ”vehicle” at a given instant of time. At each discrete time step $`tt+1`$, the state of the system is updated following a well defined prescription. In the TASEP a randomly chosen particle can move forward, by one lattice spacing, with probability $`q`$ if the lattice site immediately in front of it is empty. In the NaSch model, the speed $`v`$ of each vehicle can take one of the $`v_{max}+1`$ allowed integer values $`v=0,1,\mathrm{},v_{max}`$. If the random-sequential updating scheme of the TASEP is replaced by parallel updating then it becomes identical to the NaSch model with $`v_{max}=1`$ and random braking probability $`p=1q`$. Our interest in the NaSch model is to unravel its connections to the BRM. For this purpose, we only need the NaSch model only with $`v_{max}=1`$. Thus in what follows, by the NaSch model, we shall mean NaSch model with $`v_{max}=1`$, unless explicitly stated otherwise.
### B BRM with parallel and random-sequential updatings
In the BRM each of the lattice sites represents a bus stop and these stops are labeled by an index $`i`$ ($`i=1,2,\mathrm{},L`$) . In each step of updating, each bus attempts to hop from one stop to the next. Note that in the TASEP and the NaSch model one can label the lattice sites by the index $`i`$ ($`i=1,2,\mathrm{},L`$) and describe the state of each of the sites by associating a variable $`\sigma _i`$ with it; $`\sigma _i=1`$ if the site $`i`$ is occupied and $`\sigma _i=0`$ if the site $`i`$ is empty. In contrast, in the BRM, two binary variables $`\sigma _i`$ and $`\varphi _i`$ are assigned to each site $`i`$: (i) If the site $`i`$ is occupied by a bus then $`\sigma _i=1`$; otherwise $`\sigma _i=0`$. (ii) If site $`i`$ has passengers waiting for a bus then $`\varphi _i=1`$; otherwise $`\varphi _i=0`$. A site cannot have both $`\sigma _i=1`$ and $`\varphi _i=1`$ simultaneously since a site cannot have simultaneously a bus and waiting passengers. The state of the system is updated according to the following random sequential update rules: a site $`i`$ is picked up at random. Then, (i) if $`\sigma _i=0`$ and $`\varphi _i=0`$ (i.e, site $`i`$ contains neither a bus nor waiting passengers), then $`\varphi 1`$ with probability $`\lambda `$, where $`\lambda `$ is the probability per unit time of the arrival ($`i.e.`$ the arrival rate) of the first passenger at the empty bus stop. (Arrival of the subsequent passengers does not affect the time evolution.) (ii) If $`\sigma _i=1`$ (i.e., there is a bus at the site $`i`$) and $`\sigma _{i+1}=0`$, then the hopping rate $`\mu `$ of the bus from site $`i`$ to $`i+1`$ is defined as follows: (a) if $`\varphi _{i+1}=0`$, then $`\mu =\alpha `$ but (b) if $`\varphi _{i+1}=1`$, then $`\mu =\beta `$, where $`\alpha `$ is the hopping rate of a bus onto a stop which has no waiting passengers and $`\beta `$ is the hopping rate onto a stop with waiting passenger(s). Generally, $`\beta <\alpha `$, which reflects the fact that a bus has to slow down when it has to pick up passengers. In the BRMRSU one can set $`\alpha =1`$ without loss of generality. However, for reasons which will become clear soon, we shall keep $`\beta <\alpha <1`$. When a bus hops onto a stop $`i`$ with waiting passengers, $`\varphi _i`$ is reset to zero as the bus takes all the passengers. Note that the density of buses $`c=N/L`$ is a conserved quantity whereas that of the passengers is not.
In the BRMPU the random sequential update rules of the BRMRSU are replaced by parallel updating but all the other aspects of the updating remain unchanged. BRMPU is related to the NaSch model in two extreme limits of $`\lambda `$: In the unphysical limit of $`\lambda =0`$ (which means the passengers never arrive at a busstop), the BRMPU reduces to the NaSch model with $`v_{max}=1`$ and $`q=1p=\alpha `$. In the opposite limit of maximum value of $`\lambda `$ in BRMPU, $`\lambda =\mathrm{}`$ (very fast rate of passenger arrival at a busstop), the BRMPU is equivalent to the NaSch model with $`v_{max}=1`$ and $`q=1p=\beta `$. Note that since the time between two updating steps is the unit of time, all values of $`\lambda 1`$ are synonymous with $`\lambda =\mathrm{}`$. This is because any value of $`\lambda 1`$ will bring at least one passenger to an empty bus stop between two updating times. Interesting results in the model occur only for values of $`\lambda 1`$.
Note that if we take $`\alpha =1`$, then the limit $`\lambda =0`$ would correspond to the limit $`q=1`$ (i.e., $`p=0`$) of the NaSch model which is a deterministic CA and does not exhibit jammed states . Since we are interested in exploring the connection between the BRMPU and the NaSch model with arbitrary $`q`$ throughout this paper we consider $`\alpha <1`$.
### C Extended NaSch models
It has been realized over the last few years that different modifications of the braking rule in the NaSch model can lead to different types of phenomena which are interesting from the perspective of statistical physics. For example, such modifications can lead to self-organized criticality as well as metastability and phase segregation . Klauck and Schadschneider considered a model where the particle is allowed to hop forward by one site or by two sites with two different hopping rates. It has also been established that assigning quenched random hopping rates can lead to the formation of clusters of vehicles .
In a similar vein, we now extend the NaSch model by replacing its constant (time-independent) hopping rate $`q`$ by two other alternatives which are intended to mimic the situations in the BRMPU.
In one of these two alternatives (from now onwards referred to as Model Y) the hopping rate of a vehicle at a given site $`x`$ is given by
$$q_x=\beta +(\alpha \beta )e^{\mathrm{\Lambda }T_{x+1}}$$
(1)
where $`\mathrm{\Lambda }>0`$ is a constant and $`T_{x+1}`$ is the time interval that has elapsed since the leading vehicle (LV) left the site $`x+1`$. $`T_{x+1}`$ is therefore the time interval between the departure of LV and the arrival of the following vehicle (FV) at the site $`x+1`$.
In the other extended NaSch model (from now onwards referred to as Model Z) the hopping rate of the $`n`$-th vehicle depends on its instantaneous DH $`\mathrm{\Delta }x_n`$ :
$$q_n=\beta +(\alpha \beta )e^{\mathrm{\Lambda }\mathrm{\Delta }x_n/\beta }$$
(2)
where $`\beta `$ ($`<1`$) is a constant. At first sight it may seem more appropriate to have $`v`$, rather than $`\beta `$, in the exponential in equation (2). However, in section III and fig.3(b), we show that the extended NaSch model with the hopping rates of the form (2) is a good approximation of the BRMPU in a wide range of circumstances.
Both the models Y and Z reduce to the NaSch model with constant hopping rates $`\alpha `$ in the limit $`\mathrm{\Lambda }=0`$, and reduce to the NaSch model with a constant hopping rate $`\beta `$ in the limit $`\mathrm{\Lambda }=\mathrm{}`$.
Models Y and Z are devised to capture the essential features of a bus-route model where the time- /space- dependent hopping rates of the vehicles depend on the presence or absence of waiting passengers.
### D Methods of Simulation
For the numerical calculations of the various quantities through computer simulations, we let the system evolve from a random initial state following the appropriate updating rules mentioned above. We compute the quantities relevant for the investigation of the kinetics of the system during the time-evolution of the system towards its steady-state. After the system reaches steady-state, we compute its steady-state properties, e.g, the TH distribution, by letting it evolve for the next $`5\times 10^4`$ time steps to obtain the required data. We then repeat the calculation with a different random initial state and, finally, average the data over $`100`$ different random initial states of the system.
The largest systems we have simulated have a total length $`L=10^5`$; each sample of these was allowed to evolve upto a maximum of $`10^6`$ time steps which is not long enough to reach the corresponding steady-state but were used for the study of the kinetics. For the computation of the average steady-state properties we have used smaller systems (typically $`L=10^4`$) which require shorter time to reach steady-state. In all our simulations we have used a periodic boundary condition.
## III Results of the extended NaSch models Y and Z
In fig.1 we plot the TH distributions in the models Y and Z for $`\mathrm{\Lambda }=0.01`$, $`\alpha =0.9`$, $`\beta =0.5`$ at two different densities, namely, $`c=0.1`$ and $`c=0.5`$. Note that in the special case $`\mathrm{\Lambda }=0`$, both the models Y and Z as well as the BRMPU reduce to the NaSch model with $`q=\alpha `$. The data in fig.1 establish that, when $`\mathrm{\Lambda }`$ is sufficiently small (e.g., $`\mathrm{\Lambda }=0.01`$), the results of the models Y and Z agree well with those of the BRMPU at all densities for identical values of the set of parameters.
In order to emphasize the effects of time-/space-dependence of the hopping rates on the TH distribution we plot in fig. 2 the exact TH distributions in the NaSch model with $`v_{max}=1`$ for $`q=0.9`$ and $`q=0.5`$ at the same two densities as used in fig.1. The observation that the TH distribution in the NaSch model for $`q=0.9,c=0.5`$ is narrow can be explained by fact that small noise ($`p=0.1`$) gives rise to only a small width of the $`\delta `$-function-like TH distribution, centered at $`2`$, that one would observe at $`c=0.5`$ in the deterministic limit $`q=1.0`$ of the NaSch model. Comparing fig.2 with the fig.1 we find that, except for $`q=0.9,c=0.1`$, the TH distributions in the NaSch model have much longer tail than those in the BRMPU as well as in the model Y and model Z for the parameters $`\alpha =0.9,\beta =0.5`$.
Thus, the BRMPU is well approximated by both the models Y and Z at $`\mathrm{\Lambda }`$ as small as $`0.01`$. However, we find a larger difference between the TH distributions in the models Y and Z at larger values of $`\mathrm{\Lambda }`$ (see fig.(3a)). For $`\lambda 1`$ the BRMPU reduces to the NaSch model with $`q=\beta `$ and the corresponding TH distribution is in excellent agreement with that in the model Z but differs significantly from that in the model Y (see fig.(3b)).
The results in figures 1,2 and 3 show that in a certain density regime, the BRMPU is well approximated by the NaSch model with time-/space-dependent hopping rates. We expect the passenger arrival rate $`\lambda `$ of the BRMPU and the hopping rate $`\mathrm{\Lambda }`$ in models Y and Z to be related. We assume that the two are equal, even though we continue to use the two different symbols in order to allow the possibility of a difference between them in future simulations.
## IV Analytical results of the BRMPU
For analytical calculation of the TH distributions we label the site (i.e., the bus stop) where the detector is located by $`j=0`$, the stop immediately in front of it by $`j=1`$, and so on. The detector clock resets to $`\tau =0`$ everytime a bus leaves the detector site. We begin our analytical calculations by writing $`𝒫_{th}(\tau )`$, the probability of a TH $`\tau `$ between the LV and the FV of a pair, as
$$𝒫_{th}(\tau )=\underset{t_1=1}{\overset{\tau 1}{}}P(t_1)Q(\tau t_1|t_1)$$
(3)
where $`P(t_1)`$ is the probability that there is a time interval $`t_1`$ between the departure of the LV and the arrival of the FV at the detector site and $`Q(\tau t_1|t_1)`$ is the conditional probability given that the FV arrives at the detector site $`t_1`$ time steps after the departure of the LV, it halts for $`\tau t_1`$ time steps at that site.
Encouraged by the success of the models Y and Z in capturing the TH distributions over moderate and high density regimes, we now approximate the BRMPU as an extended NaSch model with a time-dependent hopping rate which is closely related to (but slightly different from) those in the models Y and Z. Thus for our analytical calculation of $`𝒫_{th}(\tau )`$ in the BRMPU, we approximately treat it as an extended NaSch model (with $`v_{max}=1`$) where the hopping rate $`q`$, instead of being a constant, is a time-dependent quantity given by the expression
$$q=\beta +(\alpha \beta )e^{\mathrm{\Lambda }t_1}.$$
(4)
This form can be compared to those given in equations (1) and (2).
The exact analytical expression for $`P(t_1)`$ in the NaSch model (with $`v_{max}=1`$) has been derived earlier using a 2-cluster approximation which goes beyond the simple mean field approximation. Following the same arguments we now get
$$P_{cl}(t_1)=𝒞(1|\underset{¯}{0})q\left[𝒞(0|\underset{¯}{0})q+p\right]^{t_11}$$
(5)
where $`q`$ is given by (4) and $`𝒞`$ gives the 2-cluster steady-state configurational probability for the argument configuration; the underline under an argument of $`𝒞`$ implies the associated condition. The expressions for the various $`𝒞`$s are given by
$$𝒞(\underset{¯}{1}|0)=𝒞(0|\underset{¯}{1})=\frac{y}{c}$$
(6)
$$𝒞(\underset{¯}{0}|1)=𝒞(1|\underset{¯}{0})=\frac{y}{d}$$
(7)
$$𝒞(\underset{¯}{1}|1)=𝒞(1|\underset{¯}{1})=1\frac{y}{c}$$
(8)
$$𝒞(\underset{¯}{0}|0)=𝒞(0|\underset{¯}{0})=1\frac{y}{d}$$
(9)
where
$$y=\frac{1}{2q}\left(1\sqrt{14qcd}\right),$$
(10)
$`q=1p`$ and $`d=1c`$.
On the other hand, in the simple mean field approximation, the 2-cluster probabilities reduce to $`𝒞(1|\underset{¯}{0})c`$ and $`𝒞(0|\underset{¯}{0})1c`$ and, hence,
$$P_{mf}(t_1)=cq[1cq]^{t_11}$$
(11)
We shall calculate the TH distribution, $`𝒫_{th}(\tau )`$, given in (3), using the expression (5) (together with (7), (9) and (10)) and then compare with the corresponding TH distribution obtained by using (11), instead of (5), to emphasize the importance of correlations.
In order to obtain $`𝒫_{th}(\tau )`$, let us next calculate $`Q(\tau t_1t_1)`$. Again, following the arguments used earlier in the calculation of the TH distribution in the NaSch model, we get
$`Q(\tau t_1t_1)`$ $`=`$ $`(1\overline{g}^{t_1})p^{\tau t_11}q`$ (13)
$`+\overline{g}^{t_1}gq{\displaystyle \frac{[(\overline{g})^{\tau t_11}(p)^{\tau t_11}]}{\overline{g}p}}`$
where $`g`$ is the probability that a vehicle moves in the next time step (i.e., in the $`(t+1)^{th}`$ time step) and $`\overline{g}=1g`$. In the 2-cluster approximation
$$g_{cl}=q𝒞(\underset{¯}{1}|0)$$
(14)
which, in the simple mean field approximation reduces to
$$g_{mf}=q(1c)$$
(15)
Substituting (5) and (13) into (3) and using (14) for $`g`$ and (4) for $`q`$ (together with (6)-(9) for the configurational probabilities and (10)), we get $`𝒫_{th}(\tau )`$ in the 2-cluster approximation by carrying out the summation over $`t_1`$ in (3) numerically. We shall refer to this result as the 2-cluster estimate of the TH distribution. Similarly, substituting (11) and (13) into (3) and using (15) for $`g`$ and (4) for $`q`$ (together with (6)-(9) and (10)) we get the simple mean field estimate of $`𝒫_{th}(\tau )`$ by again summing over $`t_1`$ numerically. Note that in both cases, $`𝒫_{th}(\tau )`$, in addition to its $`\tau `$ dependence, depends on parameters $`c,\alpha ,\beta ,\mathrm{\Lambda }.`$
In fig.4 we show these analytic results for three different values of $`\mathrm{\Lambda }`$ and two different values of densities $`c`$. At sufficiently low density of buses, there is hardly any difference between the 2-cluster estimate and the simple mean field estimate of the TH distribution (see fig.4(a)). However, with the increase of the density of the buses, the difference between these two estimates increases (see fig.4(b)).
As noted earlier, the TH distribution in the BRMPU changes continuously with the variation of $`\lambda `$; the results for $`\lambda =0`$ and $`\lambda =1`$ are identical to those in the NaSch model with $`q=\alpha `$ and $`q=\beta `$, respectively.
In fig. 5 we compare the 2-cluster analytic estimate of the TH distribution in the BRMPU (approximated as the extended NaSch model) for three different values of $`\mathrm{\Lambda }`$, namely, $`\mathrm{\Lambda }=0.01,0.10,0.50`$, with the corresponding numerical data we have obtained from direct computer simulations of the BRMPU model. Figure 5 (a) shows the comparison for the higher density value $`c=0.5`$. We note very good agreement between the 2-cluster estimates and the computer simulation data. Similar comparison at a lower density $`c=0.1`$ is shown in fig. 5 (b). The poor agreement between the 2-cluster estimate and simulation data for the TH distribution in the BRMPU at low densities is a consequence of the fact that at low densities, the vehicles in a finite system form a cluster in the steady-state where there is at most one or two empty sites in between each pair of vehicles. It is well known that the 2-cluster approximation scheme is not good enough for such states where correlation extends over distances which are much longer than what can be captured by a 2-cluster approximation. On the other hand, at higher densities there is no clustering of the buses in the steady-state of the BRMPU and the physics of the system is very similar to that in the steady-states of the NaSch model. This is because most of the buses stop on account of another bus at the next site, rather than due to waiting passengers.
Nevertheless, since the states with clusters of buses are metastable in infinitely long samples of BRMPU, the 2-cluster approximation is expected to yield good estimates for the stable steady-states of the BRMPU even at low densities. However, since it is extremely difficult (requires very long simulation time) to achieve these stable steady-states in any computer simulation at small $`\lambda `$ we have not been able to demonstrate this explicitly.
## V Kinetics in the BRMPU
In traffic models like BRMPU, the kinetics are governed by two coupled fields, local passenger density $`\varphi _i(t)`$ and local bus density $`\sigma _i(t)`$. In this paper a binary approximation (zero or nonzero) to the $`\varphi _i(t)`$ field is made. $`\sigma _i(t)`$ is globally conserved and $`\varphi _i(t)`$ is a nonconserved field. It is not clear whether the kinetics seen in our simulations of BRMPU and models Y and Z are derivable from a free energy functional. We are currently inquiring into such a possibility. If this turns out to be the case, then the model of kinetics appropriate to our simulations is model C in the Halperin- Hohenberg classification scheme of critical dynamics.
Let us define the correlation function
$$𝒞(r,t)=\left[\frac{1}{L}\underset{i=1}{\overset{L}{}}\sigma _i(t)\sigma _{i+r}(t)c^2\right]$$
(16)
where $`t=0`$ corresponds to the initial state. The symbol $`[.]`$ indicates average over random initial conditions. By definition, $`𝒞(r,t)`$ vanishes in the absence of any correlation in the occupation of the sites by the buses. Also at any time $`t`$, $`𝒞(r=0,t)=c(1c).`$ This correlation function has been calculated earlier for the NaSch model analytically for $`v_{max}=1`$ and numerically for higher values of $`v_{max}`$. However, the nature of this correlation function in the BRMPU is expected to differ qualitatively, particularly at low densities, from that in the NaSch model because of the formation of clusters in the steady-state of the BRMPU.
We compute $`𝒞(r,t)`$ during the time-evolution of the system from random initial states. Our simulations for the time evolution are done only for one set of parameters: $`c=0.1,\lambda =0.01,\alpha =0.9`$ and $`\beta =0.5.`$ In fig.6 we plot the normalized correlation function
$$G(r,t)=𝒞(r,t)/𝒞(r=0,t)=𝒞(r,t)/(c(1c))$$
(17)
as a function of $`r`$ for values of $`t`$ upto $`5\times 10^6`$.
The value $`r=R`$ corresponding to the first zero-crossing of $`G(r,t)`$ is taken as a measure of the typical size of the clusters of buses at time $`t`$ . The fact that $`R`$ increases with $`t`$ indicates the coarsening of these clusters. It is worth mentioning here that in systems with conserved order parameters (the so-called model B, of which the binary alloy is a physical realization) the coarsening follows the Lifshitz-Slyozov law $`R(t)t^{1/3}`$. In the case of the BRMPU, $`R(t)`$ may appear to follow the same Lifshitz-Slyozov law if the data upto $`t10^6`$ are shown on a log-log plot (see fig.7). However, the upward turn of the data beyond $`t10^6`$ in fig.7 indicates more subtle features of this growth. In fact, fitting the raw data $`R(t)`$ to the curve
$$R(t)=R_0+At^{1/2}$$
(18)
we have estimated the parameters $`R_0`$ and $`A`$. We found that $`R_055`$ and $`A0.2`$. If the growth of $`R(t)`$, indeed, follows the law (18) then the $`t^{1/2}`$ growth law is expected to become clearly visible directly in fig.7 for times long enough to satisfy the condition $`At^{1/2}>>R_0`$, i.e., for $`t>>10^5`$. This argument, together with our estimate $`R_055`$, explains why the true growth law (18) can be anticipated in fig.7 only beyond $`t10^6`$. In fact, plotting $`R(t)`$ against $`t^{1/2}`$ and comparing with $`55+0.2t^{1/2}`$ in fig.8 we do, indeed, see clear evidence of the $`t^{1/2}`$ growth for $`t>>10^5`$.
Finally, in fig.9 we plot the normalized correlation function $`G`$ against the scaled variable $`r/R(t)`$. Since the data for $`t`$ as widely separated as $`t=5\times 10^3`$ to $`t=5\times 10^6`$ superpose, the validity of dynamic scaling is convincing for kinetics of the BRMPU model. The $`t^{1/2}`$ growth in our BRMPU model that is akin to model C implies that it is the nonconserved passenger density field $`\varphi `$ that is driving the kinetics, for the parameter set that we have simulated.
## VI Conclusions
The models and simulations presented in this paper were inspired by the work on BRMRSU in ref. . The model in uses random sequential updating (RSU), whereas our complimentary study is based on parallel updating (PU). As expected, the properties of the steady states including the dependence of average velocity and current on the particle density are similar (not explicitly displayed in figures). We have in addition computed and analytically obtained (in the $`2`$-cluster approximation) the time headway (TH) distributions for a wide range of parameters. At moderate and high densities, the simulation and analytic results agree well; but the comparison fails at low densities. We have also studied kinetics of the BRMPU which is also complimentary to that done for the BRMRSU in ref. . We find the bus clusters to grow in size as $`t^{1/2}`$. This shows that the growth exponent is robust with respect to the updating schemes. We have also computed the equal time pair correlation function of the local bus density, and find that it obeys a dynamical scaling ansatz: $`G(r,t)=G(r/R(t))`$, where $`R(t)`$ is the bus cluster size.
We have also found connection between the models of BRMPU type and the NaSch model. The NaSch model is the minimal CA model of vehicular traffic on idealized single-lane highways. This model has been extended in various ways to incorporate some aspects of real traffic which are not captured by the minimal model. All the bus route models which we have considered in this paper, namely, the BRMPU, model Y and model Z, may be regarded as extensions of the NaSch model with $`v_{max}=1`$. In each of these models the dynamics of the vehicles are coupled to another non-conserved field, namely, the passenger density field $`\varphi `$, resulting in space- and time-dependent hopping rates of the vehicles. However, there is a difference in the length and time scales in the bus route models and in the NaSch model. In the NaSch model the motion of the vehicles from one cell to the next is described on a time scale which is roughly the reaction time of each driver. In contrast, the lattice constant in the BRM is of the order of the distance between successive stops on the bus route. Thus, each time step corresponds to a much longer real time than that in the NaSch model. In other words, the interactions of the buses with the traffic, on its way from one stop to the next, are included in the BRM models only through the phenomenological rate constant $`\alpha `$. It would be interesting to extend the bus route models further by including the interaction between a bus and other vehicles as it moves from one stop to the next. Inter-vehicle interactions are a natural ingredient of the NaSch model, and this prescription can be incorporated in the bus route models.
Acknowledgements: We thank A. Schadschneider and D. Stauffer for valuable comments and suggestions, and the Department of Physics at the University of Toronto for providing us free CPU time on Helios2 without which this work could not have been completed. This work was also supported by NSERC of Canada.
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# COSMOLOGICAL STRUCURE FORMATION WITH TOPOLOGICAL DEFECTS
## 1 Introduction
Even if the big bang is an “irreproducible experiment”, we want to learn from it as much as possible about the physics at high energies. We have reasons to hope that it may have left traces from energies much higher than those reached in any astrophysical event or terrestrial experiment. Therefore, even if it is irreproducible and hence not as controllable as we might want, we simply cannot afford to ignore the information it may have left.
The initial fluctuations in the cosmic matter density and geometry may represent one such trace. In fact, presently there are two relatively worked out ideas for cosmic initial fluctuations, both relying on the physics at very high energies. In the first model, cosmic initial perturbations are due to quantum fluctuations which ’freeze in’ as classical fluctuations when they become super-horizon during an inflationary era.
The second possibility is that topological defects which may have formed during a phase transition in the early universe have induced structure formation. This second possibility is the topic of this talk.
Here, a pedagogical remark may be in order: Often these two alternatives have been represented as ’inflation versus defects’. This is of course not quite correct, as topological defects have nothing to say about the flatness, the horizon and the monopole or moduli (or whatever unwanted relicts) problems which inflation also solves. It is, however, easy to construct inflationary models where the amplitude of initial fluctuations is much too small to be relevant for structure formation. Therefore, in a model, where cosmic structure is due to topological defects, one needs either inflation prior to defect formation or another mechanism to solve the flatness, horizon and relict problems.
The reminder of this talk is organized as follows: In Section 2 I give a short overview on the formation of topological defects during cosmological phase transitions. In Section 3 I discuss the problem of structure formation with topological defects. I will first describe some generic insights and then discuss results for specific models. Conclusions are presented in Section 4.
## 2 Topological defects
During adiabatic expansion the universe cools down from a very hot initial state. It is natural to expect that the cosmic plasma undergoes several symmetry breaking phase transitions. In the process of such a transition an initial symmetry group $`G`$ is broken down to a subgroup $`H`$. Depending on the topology of the vacuum manifold $``$, which generically is topologically equivalent to the homogeneous space $`G/H`$, topological defects may form.
This is described by an order parameter or Higgs field, $`\varphi `$, with a temperature dependent effective potential. The field values which minimize the potential form the vacuum manifold $``$. After the phase transition the field will assume different values in $``$ in different positions of physical space, which are uncorrelated if, e.g. the spatial separation is larger than the present particle horizon, $`l_Ht`$. If the topology of the vacuum manifold is non-trivial, the Kibble mechanism $`^\mathrm{?}`$ generically leads to the formation of topological defects: the field $`\varphi `$ may vary in space in such a way that there are points, where $`\varphi `$ has to leave the vacuum manifold by continuity reasons and assume values with higher potential energy. Such points have to form a connected sub-manifold of spacetime.
For example if $``$ is not connected, $`\pi _0()\{0\}`$, in different positions $`\varphi `$ can assume values which belong to disconnected parts of $``$ and therefore is has to leave the vacuum manifold somewhere in between . The sub-manifold of higher energy is in this case three dimensional in spacetime and is called domain wall. (Domain walls from high energy phase transitions are disastrous for cosmology.) Similarly, a non simply connected vacuum manifold, $`\pi _1()\{0\}`$, leads to the formation of two dimensional defects, cosmic strings. Domain walls and cosmic strings are either infinite or closed. If $``$ contains non shrinkable two spheres, $`\pi _2()\{0\}`$, one dimensional defects, monopoles form. Finally, if $`\pi _2()\{0\}`$, zero dimensional textures appear, which are events of higher energy. By Derrick’s theorem one can argue that a scalar field configuration with non-trivial $`\pi _3`$ winding number (i.e. a texture knot) contracts and eventually unwinds producing a space-time ’point’ of higher energy. A summary of this is given in Table 1; more details can be found in Refs. $`^{\mathrm{?},\mathrm{?}}`$.
Topological defects are also very well known in solid state physics. For example the vortex lines in type II super conductors are nothing else than cosmic strings. Also in liquid crystals $`^\mathrm{?}`$ (see Fig. 1) or super fluid Helium $`^\mathrm{?}`$ a variety of topological defects form during symmetry breaking phase transitions.
The defects are called local, if a gauge symmetry is broken and global if they emerge from global symmetry breaking. In the case of local defects, gradients in the scalar field are ’compensated’ by the gauge field and the energy density of the defect is confined to the defect manifold with very small transverse dimension of the order of the symmetry breaking scale. Soon after formation, local defects therefore seize to interact over distances larger than the inverse symmetry breaking scale.
The energy density of global defects is dominated by gradient energy and hence of the order of $`\rho _{\mathrm{defect}}T_c^2/t^2`$ where $`T_c`$ is the symmetry breaking temperature and $`t`$ is the horizon scale, the typical scale over which the scalar field varies. As the energy density of the cosmic fluid also decays like $`1/t^2`$, global defects always scale<sup>a</sup><sup>a</sup>aUp to logarithmic corrections to the scaling law which are especially important in the case of global cosmic strings. and lead to fluctuations with a typical amplitude of
$$\rho _{\mathrm{defect}}/\rho 4\pi GT_c^2=ϵ.$$
(1)
In the case of local defects only cosmic strings scale and obey (1). Local monopoles soon come to dominate the cosmic energy density and are therefore ruled out from observations. Local texture die out. To be relevant for structure formation, the defects have to induce scaling fluctuations with an amplitude $`ϵ10^5`$ with implies
$$T_c10^{16}\mathrm{GeV},$$
a grand unified energy scale. Topological defects which form at lower temperature are of no relevance for structure formation<sup>b</sup><sup>b</sup>bWith a possible exception of ’soft domain walls’, see e.g. $`^\mathrm{?}`$ or the contribution of M. Bucher to these proceedings..
## 3 Structure formation with topological defects
We discuss especially the differences of structure formation with topological defects from inflationary initial perturbations. I first highlight some very generic features, then we discuss results for specific models.
### 3.1 Generics
The large scale fluctuations in the cosmic microwave background (CMB) are of the same order as the deviation of the cosmic metric from a Friedmann metric. Since these fluctuations are small, linear perturbation theory is justified. For a cosmic fluid consisting of radiation, massless neutrinos, baryons, cold dark matter, possibly hot dark matter and/or a cosmological constant, we obtain linear perturbation equations (in Fourier space). For each wave vector $`𝐤`$ they are of the form
$$DX=𝒮,$$
(2)
where $`X`$ is a long vector describing all the random perturbation variables, $`D`$ is a deterministic linear first order differential operator and $`𝒮`$ is a random source term which consists of linear combinations of the energy momentum tensor of the defect network. More details can be found, e.g. in Ref. $`^\mathrm{?}`$.
For inflationary perturbations $`𝒮=0`$ and the solutions are determined entirely by the random initial conditions, $`X(𝐤,t_{\mathrm{in}})`$. For most inflationary models $`X(𝐤,t_{\mathrm{in}})`$ is a set of Gaussian random variables and hence their statistical properties are entirely determined by the spectra $`𝒫`$ (the Fourier transforms of the two point functions),
$$X_i((t_{\mathrm{in}},𝐤)X_j^{}((t_{\mathrm{in}},𝐤^{})𝒫_{ij}(𝐤)\delta (𝐤𝐤^{}).$$
(3)
Here the Dirac delta is a consequence of statistical homogeneity which we want to assume for the random process leading to the initial perturbations.
Be $`A_i(k,t)`$ the solution with initial condition $`X_j(k,t_{\mathrm{in}})=\delta _{ij}`$. The spectra of the solution with initial ’spectrum’ given by Eq. (3) is then just
$$X_i((t_0,𝐤)X_j^{}((t_0,𝐤^{})=A_i(k,t_0)A_j^{}(k,t_0)𝒫_{ij}(k)\delta (𝐤𝐤^{}).$$
Tehrefore, if $`A_i`$ is oscillating, e.g. as a function of $`kt`$ so will $`|X_i|^2`$. This leads to a very important feature in the CMB anisotropy spectrum, the acoustic peaks: Prior to recombination, due to radiation pressure the photon/baryon plasma undergoes acoustic oscillations on subhorizon scales. At recombination the photons become suddenly free and ’stream’ into our antennas without further interaction. Since the acoustic oscillations of a given wave number $`k`$ are all in phase, the have a fixed amplitude at decoupling. This phenomenon imprints itself in the CMB anisotropy spectrum as a series of peaks. On very small scales, the finite thickness of the recombination shell and free streaming have to be taken into account which leads to an exponential damping of the peaks (Silk damping). As we shall see below, the acoustic peaks are very characteristic of inflationary perturbations.
If the source term $`\mathrm{SS}`$ does not vanish, the situation is different. Equation (2) can be solved by means of a Green’s function (kernel), $`𝒢(t,t^{})`$, in the form
$$X_j(t_0,𝐤)=_{t_{in}}^{t_0}𝑑t𝒢_{jl}(t_0,t,𝐤)\mathrm{SS}_l(t,𝐤).$$
(4)
Power spectra or, more generally, quadratic expectation values of the form $`X_j(t_0,𝐤)X_l^{}(t_0,𝐤)`$ are then given by
$$X_j(t_0,𝐤)X_l^{}(t_0,𝐤)=_{t_{in}}^{t_0}𝑑t_{t_{in}}^{t_0}𝑑t^{}𝒢_{jm}(t_0,t,𝐤)𝒢_{ln}^{}(t_0,t^{},𝐤)\mathrm{SS}_m(t,𝐤)\mathrm{SS}_n^{}(t^{},𝐤).$$
(5)
The only information about the source random variable which we really need in order to compute power spectra are therefore the unequal time two point correlators
$$\mathrm{SS}_m(t,𝐤)\mathrm{SS}_n^{}(t^{},𝐤).$$
(6)
This nearly trivial fact has been exploited by many workers in the field, for the first time probably in Ref. $`^\mathrm{?}`$ where the decoherence of models with seeds has been discovered, and later in Refs. $`^{\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?}}`$ and others.
To determines the correlators (6) one has to calculate the unequal time correlators of the defect energy momentum tensor by means of numerical simulations. To solve the enormous problem of dynamical range, ’scaling’, statistical isotropy and causality have to be used.
Seeds from global topological defects and from cosmic strings are ’scaling’ in the sense that their correlation functions $`C_{\mu \nu \rho \lambda }`$ defined by
$`\mathrm{\Theta }_{\mu \nu }(𝐤,t)`$ $`=`$ $`M^2\theta _{\mu \nu }(𝐤,t),`$ (7)
$`C_{\mu \nu \rho \lambda }(𝐤,t,t^{})`$ $`=`$ $`\theta _{\mu \nu }(𝐤,t)\theta _{\rho \lambda }^{}(𝐤,t^{})`$ (8)
are scale free; i.e. the only dimensional parameters in $`C_{\mu \nu \rho \lambda }`$ are the variables $`t,t^{}`$ and $`𝐤`$ themselves. Here the energy scale $`M`$ corresponds to the symmetry breaking scale. One can set $`M=T_c`$. Up to a certain number of dimensionless functions $`F_n`$ of $`z=k\sqrt{tt^{}}`$ and $`r=t/t^{}`$, the correlation functions are then determined by the requirement of statistical isotropy, symmetries and by their dimension. Causality requires the functions $`F_n`$ to be analytic in $`z^2`$. A more detailed investigation of these arguments and their consequences is presented in Ref. $`^\mathrm{?}`$. There it is shown that statistical isotropy and energy momentum conservation reduce the correlators (8) for global defects to five such functions $`F_1`$ to $`F_5`$. Since cosmic strings loose energy by gravitational radiation, which is crucial to ensure scaling, in this case 14 functions $`F_n`$ are needed to fully describe the correlators. However, numerical simulations show that for cosmic strings the density-density correlator is significantly larger than all the other components of $`C_{\mu \nu \rho \lambda }`$ which again simplifies the problem $`^\mathrm{?}`$.
Since analytic functions generically are constant for small arguments $`z^21`$, $`F_n(0,r)`$ actually determines $`F_n`$ for all values of $`k`$ with $`z=k\sqrt{tt^{}}\stackrel{<}{}0.5`$. Furthermore, the correlation functions decay inside the horizon and we can safely set them to zero for $`z\stackrel{>}{}40`$ where they have decayed by about two orders of magnitude. In Fig. 2 I show one of these functions for global $`O(4)`$-texture (a) and for the large $`N`$ limit of global $`O(N)`$ models$`^\mathrm{?}`$ (b).
For the induced perturbations in the cosmic fluids, the presence of of this source term has several important consequences. First of all, as is clear from Eqn. (4), the randomness of the source term enters at all times (as long as the source term is non-zero). Therefore, fluctuations of a given wave number $`k`$ are in general not in phase, and the distinctive series of acoustic peaks present in inflationary models is blurred into one ’broad hump’. This phenomenon has been termed ’decoherence’ $`^\mathrm{?}`$. A key ingredient for decoherence to happen is the non-linearity of the time evolution of the source term<sup>c</sup><sup>c</sup>cIf the source term would evolve linearly it could just be added to the components of $`X`$ and we would obtain a new, somewhat longer linear system of equations where again randomness can enter only via the initial conditions.. Even though time evolution is deterministic, different Fourier modes mix due to non-linearity, and the randomness in one mode ’sweeps’ into the other modes. In the case of topological defects, $`\mathrm{SS}`$ is given by linear combinations of the defect stress energy tensor, $`\theta _{\mu \nu }`$, quadratic in the defect field, which itself obeys non-linear evolution equations. Only in the large $`N`$ limit, the evolution of the ’defect field’ becomes linear and decoherence is much weaker. The non-linearity of the source evolution also leads to the non Gaussianity of defect models. Even if the initial field configuration would be Gaussian (which it usually is not due to non-linear constraints), the non-linear time evolution renders the source term and therefore also the perturbations highly non Gaussian.
In Table 2 we highlight the similarities and differences of inflationary and defect models of structure formation.
### 3.2 Results
As we have seen, there are several important differences between defect models and inflationary models of structure formation. First of all, defect models generically predict scalar, vector and tensor perturbations with comparable amplitudes at horizon scale, whereas in inflationary models vector perturbations are absent (they simply have decayed from their initial values) and tensor perturbations are often significantly smaller than scalar modes. Furthermore, inflationary perturbations are usually adiabatic. This leads to an important cancelation in the temperature fluctuations due to gravity, given by $`\left(\frac{\mathrm{\Delta }T}{T}\right)_{\mathrm{grav}}=2\mathrm{\Phi }`$, where $`\mathrm{\Phi }`$ denotes denotes the Newtonian potential, and the intrinsic temperature fluctuation on large scales, which is $`\left(\frac{\mathrm{\Delta }T}{T}\right)_{\mathrm{int}}=\frac{1}{4}\delta _{rad}=\frac{1}{3}\delta _{mat}=\frac{5}{3}\mathrm{\Phi }`$ in the adiabatic case. The net result becomes $`\left(\frac{\mathrm{\Delta }T}{T}\right)_{\mathrm{SW}}=\frac{1}{3}\mathrm{\Phi }`$, the ordinary Sachs-Wolfe effect for adiabatic perturbations $`^\mathrm{?}`$.
Both these effects lower the temperature fluctuations of inflationary models on very large scales if compared to those from defect models. This leads to the result that the amplitude of fluctuations on very large scales, the height of the ’Sachs-Wolfe plateau’ is comparable to the amplitude of intermediate scales, the acoustic peak(s). This has first been noted in Ref. $`^\mathrm{?}`$. Furthermore, the isocurvature nature of defect models leads to a shift of the first acoustic peak towards smaller angular scales. For flat cosmologies the peak position is around $`\mathrm{}_{\mathrm{peak}}350450`$, depending on the specific model (to be compared with $`\mathrm{}_{\mathrm{peak}}220`$ for inflationary models).
Thorough numerical simulations from two different groups $`^{\mathrm{?},\mathrm{?}}`$ now show that CMB anisotropies from global $`O(N)`$ models do not agree with present data see Fig. 3). There models also require a very high bias to fit the galaxy power spectrum and exhibit much too low bulk flows on large scales. For example the bulk velocity on $`50h^1`$Mpc for the texture model is $`V_{50}60`$km/s whereas the measured value is more like $`V_{50}300\pm 100`$km/s.
The results for cosmic strings are somewhat more promising due to a variety of effects. Most notably the following:
* The cosmic string energy density seems to be considerably higher in the radiation era than in the matter era, therefore boosting the fluctuations on scales which enter the horizon already in the radiation dominated era of the universe, $`\stackrel{<}{}50h^1`$Mpc, just the scales where global $`O(N)`$ models are missing power.
* $`T_0^0=\rho `$ is much larger than the other components of the string energy momentum tensor. Being of scalar nature it induces only scalar perturbations so that vector and tensor perturbations are suppressed in the case of strings.
* Cosmic strings loose power on scales inside the horizon by inter-commutation and gravitational radiation. These processes are slower than the speed of light with which global defects decay. Therefore, the energy momentum tensor persists to later times, up to larger values of $`kt`$ than for $`O(N)`$ models. This induces larger fluctuations in the dark matter.
The induced fluctuations in the dark matter may even be too large on small scales, a problem which can be solved by introducing hot dark matter $`^\mathrm{?}`$. The persistence of the string energy momentum tensor induces even more decoherence $`^\mathrm{?}`$ than for $`O(N)`$ models.
Therefore, cosmic strings may lead to one broad ’acoustic hump’ but certainly not to a series of peaks. The precise height of the hump depends sensitively on several unknowns, for example on how one models the string energy momentum non conservation $`^\mathrm{?}`$ and on the small scale structure of the string network $`^\mathrm{?}`$, see, e.g., Fig. 4.
Decoherence which leads to a ’smearing’ of acoustic peaks (if they are there) is one of the few features about which all results on cosmic strings agree. The height of the ’acoustic hump’ may be about two to four times the height of the plateau at low $`\mathrm{}`$. The position of the hump is not very well defined and depends on the details of the modelling, but it is typically at $`\mathrm{}\stackrel{>}{}400`$ for a flat universe, which is in disagreement with the new data shown in Fig. 3. The bias factor needed in the dark matter spectra (maybe between 2 and 5) are still quite uncertain. Some recent work on this subject can be found in Refs. $`^{\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?}}`$.
## 4 Conclusions
All the defect models studied in detail are in disagreement with current observations. They exhibit no acoustic peaks (global $`O(N)`$ models) or only one broad hump on too small scales (cosmic strings). Decoherence, which is inherent to the non-linear evolution of the defect source term smears out the distinguished series of acoustic peaks expected in inflationary models. The width of the first peak measured by the the Toco$`^\mathrm{?}`$ and BoomerangNA$`^\mathrm{?}`$ experiments is relatively narrow, which already clearly disfavors a model where decoherence is important. Secondary peaks in the CMB anisotropy spectrum will finally be a unambiguous sign for a (quasi-)linear process of structure formation like, e.g., inflation.
It has been shown, however, that linearly evolving causal scaling seeds might mimic an inflationary CMB and dark matter power spectrum $`^{\mathrm{?},\mathrm{?}}`$. Nevertheless, due to causality they differ from inflation in the CMB polarization spectrum$`^\mathrm{?}`$. Clearly, such seeds are not topological defects and there is so far no convincing physical motivation to introduce them.
## Acknowledgments
It is a pleasure to acknowledge useful discussions with Martin Kunz, Joao Magueijo and Alessandro Melchiorri. This work is supported by the Swiss National Science Foundation.
## References
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# References
CAN A KASNER UNIVERSE WITH A VISCOUS COSMOLOGICAL FLUID BE ANISOTROPIC?
I. Brevik<sup>1</sup><sup>1</sup>1E-mail address: iver.h.brevik@mtf.ntnu.no
Division of Applied Mechanics, Norwegian University of Science and Technology,
N-7491 Trondheim, Norway
S. V. Pettersen<sup>2</sup><sup>2</sup>2E-mail address: svp@fysel.ntnu.no
Department of Physical Electronics, Norwegian University of Science and Technology,
N-7491 Trondheim, Norway
PACS numbers: 98.80.Hw, 98.80.Bp
March 2000
## Abstract
A Bianchi type -I metric of Kasner form is considered, when the space is filled with a viscous fluid. Whereas an ideal (nonviscous) fluid permits the Kasner metric to be anisotropic provided that the fluid satisfies the Zel’dovich equation of state, the viscous fluid does not permit the Kasner metric to be anisotropic at all. In the latter case, we calculate the Kasner (isotropic) metric expressed by the fluid’s density, pressure, and bulk viscosity, at some chosen instant $`t=t_0`$. The equation of state is also calculated. The present paper is related to a recent Comment of Cataldo and del Campo \[Phys. Rev. D, scheduled to April 15, 2000\], on a previous work of the present authors \[Phys. Rev. D 56, 3322 (1997)\].
Consider a cosmic fluid, endowed with a bulk viscosity $`\zeta `$ and a shear viscosity $`\eta `$. In a homogeneous (possibly anisotropic) space, $`\zeta `$ and $`\eta `$ are independent of position, but will in general depend on time. If $`U^\mu =(U^0,U^i)`$ is the fluid’s four-velocity, the energy-momentum tensor is
$$T_{\mu \nu }=\rho U_\mu U_\nu +(p\zeta \theta )h_{\mu \nu }2\eta \sigma _{\mu \nu }.$$
(1)
Here $`h_{\mu \nu }=g_{\mu \nu }+U_\mu U_\nu `$ is the projection tensor, $`\theta =\theta _\mu ^\mu =U_{}^{\mu }{}_{;\mu }{}^{}`$ is the scalar expansion, $`\theta _{\mu \nu }=\frac{1}{2}(U_{\mu ;\alpha }h_\nu ^\alpha +U_{\nu ;\alpha }h_\mu ^\alpha )`$ is the expansion tensor, and $`\sigma _{\mu \nu }=\theta _{\mu \nu }\frac{1}{3}h_{\mu \nu }\theta `$ is the shear tensor.
Assume now that the metric of the background space, in which the cosmic fluid resides, is of the Kasner form:
$$ds^2=dt^2+t^{2p_1}dx^2+t^{2p_2}dy^2+t^{2p_3}dz^2.$$
(2)
The three numbers $`p_1,p_2,p_3`$ are required to be constants. The Kasner metric is a subclass of the Bianchi type-I metrics. The space is anisotropic if at least two of the three $`p_i`$ are different. Defining numbers $`S`$ and $`Q`$ by $`S=_{i=1}^3p_i`$ and $`Q=_{i=1}^3p_i^2`$ we have, for a strict Kasner universe in the classical sense corresponding to a pure vacuum, that $`S=Q=1`$. Once a real, in general viscous fluid is present, however, these simple relationships are lost.
Consider next the Einstein equations. With the cosmological constant $`\mathrm{\Lambda }`$ set equal to zero, and with the notation $`\kappa =8\pi G`$, we obtain from $`R_{\mu \nu }=\kappa (T_{\mu \nu }\frac{1}{2}g_{\mu \nu }T_\alpha ^\alpha )`$ the two equations
$$SQ+\frac{3}{2}\kappa t\zeta S=\frac{1}{2}\kappa t^2(\rho +3p),$$
(3)
$$p_i(1S2\kappa t\eta )+\frac{1}{2}\kappa t(\zeta +\frac{4}{3}\eta )S=\frac{1}{2}\kappa t^2(\rho p).$$
(4)
These are the governing field equations. Adding the three equations (4) we get
$$2S(S1)=3\kappa t\zeta S+3\kappa t^2(\rho p),$$
(5)
whereas from Eqs. (3) and (5) we obtain the simple equation
$$S^2Q=2\kappa t^2\rho .$$
(6)
Faced with these equations, we see that there are several factors determining the state of the cosmic fluid:
(i) The viscosity coefficients $`\zeta `$ and $`\eta `$.
(ii) The constants $`S`$ and $`Q`$, which may be regarded as input parameters in the governing equations.
(iii) In turn, $`S`$ and $`Q`$ are for a viscous fluid determined by the energy density $`\rho =\rho _0`$ and pressure $`p=p_0`$ evaluated at some chosen instant $`t=t_0`$ which, in terms of an appropriately scaled time variable, will be taken at $`t_0=1`$. The time variations of $`\rho =\rho (t)`$ and $`p=p(t)`$ follow from the governing equations themselves. By the same token, the time variations of $`\zeta =\zeta (t)`$ and $`\eta =\eta (t)`$ follow. Once the proportionaly constants $`\{\rho _0,p_0,\zeta _0,\eta _0\}`$ are known, the Kasner parameters $`p_i`$ and, as mentioned, also $`S`$ and $`Q`$, are known.
(iv) The possible equation of state for the fluid. Conventionally, the equation of state is written in the form $`p=(\gamma 1)\rho `$, where $`\gamma `$ is a constant. An important point, which can be resolved only after a detailed analysis of the formalism, is to what extent the equation of state is fixed once $`\{\rho _0,p_0,\zeta _0,\eta _0\}`$ are known.
In a previous paper we analysed, on basis of the governing equations above, the possible equation of state for the fluid. A recent Comment of Cataldo and del Campo claims that one of the options discussed in (actually a singular option), involving a viscous fluid residing in an anisotropic space, is not physically realizable because it runs into conflict with the dominant energy condition. The remarks of Cataldo and del Campo are interesting, since they seem to narrow down, on physical grounds, the range of applicability of the anisotropic Kasner metric for a realistic cosmic fluid. In turn, this may even lead to consequences for the applicability of anisotropic viscous universe models in general. We have found it desirable, therefore, to discuss this point in some detail, from a more broad perspective than in Ref. . This is the purpose of the present paper.
Let us first consider the physical meaning of the dominant energy condition. This condition is usually expressed mathematically by saying that in a local rest orthonormal frame the magnitude of the stress components $`T_{\widehat{i}\widehat{k}}`$ are always less than or equal to the energy density component :
$$|T_{\widehat{i}\widehat{k}}|\rho .$$
(7)
This ought to be contrasted with the weak energy condition, which says that the energy density in an orthonormal rest inertial frame is always non-negative, $`\rho 0`$. The dominant energy condition is the weak energy condition with the additional requirement that the stress components - in practice the diagonal components of $`T_{\widehat{i}\widehat{k}}`$ \- should not exceed the energy density. Moreover, since the velocity of sound in a fluid equals the square root of the adiabatic derivative of $`p`$ with respect to $`\rho `$, the dominant condition is strongly related to the physical property of a sound wave that it cannot propagate faster than light.
In the present case, $`\sigma _{ii}=(p_i\frac{1}{3}S)t^{2p_i1}`$ (no sum over $`i`$) and $`\theta =S/t`$ . From Eqs. (1) and (7) we then derive, in the local rest orthonormal frame,
$$\left|p\frac{\zeta S}{t}\frac{2\eta }{t}(p_i\frac{1}{3}S)\right|\rho .$$
(8)
This is the viscous generalization of the dominant energy condition. For an ideal fluid, Eq. (8) reduces to the well known condition $`|p|\rho `$.
Let us now consider the time dependence of the physical quantities. From Eqs. (3)-(6) it is evident that $`\rho (t)`$ and $`p(t)`$ are proportional to $`t^2`$, whereas $`\zeta (t)`$ and $`\eta (t)`$ are proportional to $`t^1`$. We write
$$\rho (t)=\rho _0t^2,p(t)=p_0t^2,$$
(9)
$$\zeta (t)=\zeta _0t^1,\eta (t)=\zeta _0t^1,$$
(10)
where $`\{\rho _0,p_0,\zeta _0,\eta _0\}`$ are the proportionality constants. The proportionalities when written in this form imply that the instant $`t=0`$ is taken to be a singular point. This is in accordance with the fact that the Kasner metric when written such as in Eq. (2) implies the true singularity (i. e., the divergence of the contracted Riemann tensor components) to occur at $`t=0`$ (cf., for instance, Ref. ). We can now write the field equations (3)-(6) as
$$SQ+\frac{3}{2}\kappa \zeta _0S=\frac{1}{2}\kappa (\rho _0+3p_0),$$
(11)
$$p_i(1S2\kappa \eta _0)+\frac{1}{2}\kappa (\zeta _0+\frac{4}{3}\eta _0)S=\frac{1}{2}\kappa (\rho _0p_0),$$
(12)
$$2S(S1)=3\kappa \zeta _0S+3\kappa (\rho _0p_0),$$
(13)
$$S^2Q=2\kappa \rho _0.$$
(14)
These equations contain time-independent terms only. The dominant energy condition (8) becomes correspondingly
$$\left|p_0\zeta _0S2\eta _0(p_i\frac{1}{3}S)\right|\rho _0.$$
(15)
After having established the governing equations and the dominant energy condition, we proceed to categorize the different options for the equation of state for the fluid. We follow essentially the same line of approach as in Ref. , and consider the ideal fluid first.
Ideal fluid. When $`\zeta =\eta =0`$ we consider, as first option, the case
$$p_0=\rho _0,$$
(16)
i. e., a Zel’dovich fluid. This is a maximally stiff equation of state; the velocity of sound is equal to the velocity of light. From Eqs. (11)-(14) we obtain
$$S=1,Q=12\kappa \rho _0.$$
(17)
Thus two of the Kasner parameters, say $`p_1`$ and $`p_2`$, can be found as solutions of the second-degree equation
$$p_1^2+p_1p_2+p_2^2p_1p_2+\kappa \rho _0=0,$$
(18)
which describes a conic section in the $`p_1p_2`$ plane . If
$$\kappa \rho _0<\frac{1}{3},$$
(19)
the curve is an ellipse. In this case there is a continuous family of values for $`p_1`$ and $`p_2`$, permitting the Kasner space to be anisotropic. The third Kasner parameter follows as $`p_3=1p_1p_2`$. The condition (19) is satisfied under usual physical circumstances. If this condition is not satisfied, there is no real locus (imaginary ellipse).
It is remarkable that the initial Kasner form of the metric is, for an ideal fluid, compatible with the Zel’dovich state equation only. For any state equation different from $`p_0=\rho _0`$ (and this is our second option), we see from Eq. (4) that all the three $`p_i`$’s have to be equal. We then find, putting $`p_1=p_2=p_3a`$,
$$a=\frac{1}{6}\left[1+\sqrt{1+6\kappa (\rho _0p_0)}\right].$$
(20)
That is, knowledge about the proportionality constants $`\rho _0`$ and $`p_0`$ suffices to calculate the three equal Kasner parameters in the isotropic space. (The reason why the case $`p_0=\rho _0`$ in Eq. (20) yields $`a=\frac{1}{3}`$, instead of the family of solutions given in Eq. (18), is that we are dealing with two different options.)
Viscous fluid. This is for us of main interest. Consider first the exceptional case when
$$S=12\kappa \eta _0.$$
(21)
From Eq. (12) it might seem possible, at first sight, that in this case anisotropic space is permitted. However, if $`\rho _0p_0`$, there is a conflict between the remaining terms in the equation since the left hand side (involving positive viscosity coefficients) is positive whereas the right hand side becomes negative.
Now, the careful reader may ask: why do we know that $`\rho _0p_0`$? In particular, does this condition follow from the dominant energy condition alone, as stated in ? Let us return to Eq. (15): we see herefrom that $`\rho _0p_0\zeta _0S2\eta _0(p_i\frac{1}{3}S)`$. If $`p_i\frac{1}{3}S`$ this is a weaker condition than the condition $`\rho _0p_00`$. Inserting this inequality into Eq. (12) we find that $`Sp_i`$, which is equivalent to $`p_2+p_30`$, $`p_1+p_30`$, $`p_1+p_20`$. These inequalities are actually acceptable. The dominant energy condition alone is thus not sufficient. What really makes us accept the restriction $`\rho _0p_00`$ is rather the physical property mentioned above: a sound wave cannot propagate faster than light. We end up with the conclusion that the option (21) has to be abandoned, thus in agreement with , although our reasoning is different from that of .
In the isotropic Kasner space we now obtain, from Eq. (12),
$$a=\frac{1}{6}\left[1+\frac{3}{2}\kappa \zeta _0+\sqrt{(1+\frac{3}{2}\kappa \zeta _0)^2+6\kappa (\rho _0p_0)}\right],$$
(22)
showing that the isotropic Kasner metric (2) is completely determined when $`\rho _0`$ and $`p_0`$, and the bulk viscosity $`\zeta _0`$, are known. It is notable that the shear viscosity $`\eta _0`$ is absent in Eq. (22); this is in accordance with the fact that the shear viscosity is a concept related to an anisotropic state of motion for a fluid. If the fluid is ideal, Eq. (22) is seen to reduce to Eq. (20).
As for the equation of state for the fluid, we note from the governing equations that
$$\kappa \rho _0=3a^2,$$
(23)
$$\kappa p_0=2a3a^2+3a\kappa \zeta _0.$$
(24)
That means that we are allowed to write the equation of state in the conventional form $`p=(\gamma 1)\rho `$, or $`p_0=(\gamma 1)\rho _0`$, where the constant $`\gamma `$ is
$$\gamma =\frac{2+3\kappa \zeta _0}{3a}.$$
(25)
It ought finally to be emphasized that we have been considering the Kasner subclass of Bianchi type-I metrics only. We cannot, on basis of the remarks above, make any conclusion about the existence of viscous cosmic fluids in the whole Bianchi type-I class.
Acknowledgment
It is a pleasure to acknowledge valuable remarks by Professor Øyvind Grøn on the present manuscript.
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